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warning/0006/cond-mat0006403.html	ar5iv	text	# Flow equation approach to the sine–Gordon modelfootnote footnote This work is dedicated to Prof. Wegner on the occasion of his 60th birthday. ## I Introduction ### A Motivation Perturbative scaling theory plays a key role for analyzing the large class of physical systems with a mismatch between the high–energy scale of the model and the experimentally interesting low–energy scale. For example, in field theory one is generally interested in the universal properties at energies much lower than the UV (ultraviolet)–cutoff, or in condensed matter physics in energies and temperatures much smaller than the Fermi energy/temperature usually of order of a few 1000 K. In order to link high–energy and low–energy regimes it is of fundamental importance to perform perturbation theory in a stable order by first analyzing the effect from high–energy scales, and then progressively smaller energies. An elementary example for this procedure is provided by atomic physics, where one e.g. first establishes the fine structure of a spectrum before using these states to evaluate the hyperfine splittings. For systems with continuous energy scales, like in field theory, the above observations have led to the development of perturbative scaling theory. Perturbative scaling ideas have become a key theoretical tool for analyzing physical systems with many degrees of freedom. The principal idea is to study perturbatively the effect of lowering the high–energy cutoff by finding a Hamiltonian with this reduced cutoff and renormalized couplings that describe the same low–energy physics as the original Hamiltonian. In a path integral formulation this is conveniently achieved by successively integrating out the high–energy degrees of freedom. This procedure leads to the well–known renormalization group (RG) equations that describe the flow of the running coupling constants upon lowering the UV–cutoff. For the important class of strong–coupling problems, however, the RG–equations lead to running coupling constants that grow larger and larger at smaller energy scales (and often eventually even diverge). Since the RG–equations themselves are derived perturbatively, this means that the perturbative scaling approach breaks down for strong–coupling problems. Well–known examples for this class of models are the Kondo model in condensed matter physics or QCD in elementary particle physics. In spite of its eventual breakdown, perturbative scaling can still contribute significantly to the understanding of strong–coupling problems. For example, in the Kondo model, the divergence of the running coupling constant occurs at an energy that sets the low–energy Kondo scale of the model, which already allows considerable insight into the problem. Still the approach becomes uncontrolled since the coupling constants grow very large, and it has, so far, not been possible to extend the perturbative scaling approach in such a way that a controlled systematic expansion emerges that links weak– to strong–coupling behavior. One can sum up these observations by noting that the perturbative scaling approach often allows us to identify the relevant low–energy scale of a strong–coupling problem, but frequently not the physical behavior associated with this energy scale or the crossover behavior linking high and low energies. For an excellent review of these issues see Ref. . This paper will exemplify the way in which Wegner’s method of flow equations can overcome these shortcomings and provide an analytic description for a weak– to strong–coupling behavior crossover. In the flow equation approach, a continuous sequence of infinitesimal unitary transformations is applied to a many–particle Hamiltonian such that the Hamiltonian becomes successively more diagonal. Wegner has set up this approach in a differential formulation $$\frac{dH}{dB}=[\eta (B),H(B)].$$ (1) Here $`\eta (B)=\eta (B)^{}`$ is an anti–Hermitian operator. Therefore $`H(B)`$ as obtained by the solution of this differential equation describes a one–parameter family of unitarily equivalent Hamiltonians. $`H(B=0)=H`$ is the initial condition relating us to the original Hamiltonian $`H`$ in which we are interested. We want $`H(B=\mathrm{})`$ to be diagonal. In order to achieve this, Wegner has proposed a suitable choice for the generator $`\eta (B)`$ that we will discuss in more detail in Sect. III.A. Wegner’s construction of $`\eta `$ generates a Hamiltonian flow where the interaction matrix elements that couple degrees of freedoms with a large energy difference are removed first (for smaller $`B`$), and more degenerate matrix elements during later stages of the flow. This separation of energy scales is reminiscent of the perturbative scaling approach and allows a stable sequence of approximations. As opposed to the perturbative scaling approach, however, degrees of freedom are not integrated out in the flow equation approach, instead they are successively diagonalized. A similar framework that contains Wegner’s flow equations as a special case has independently been developed by Głazek and Wilson (similarity renormalization scheme) . So far the flow equation approach has been applied to a variety of models in condensed matter theory like the $`n`$–orbital model , impurity models like the spin–boson model and the Anderson impurity model , electron–phonon systems and spin models etc. (for an overview see also Ref. ). One advantage of this scheme lies in the observation that it is a non–perturbative approach due to the separation of energy scales, but still has access to all energy scales since no degrees of freedom are integrated out. Therefore one can investigate correlation functions on all energy scales . Also the flow equation approach allows the systematic derivation of low–energy effective Hamiltonians not plagued with singular interactions that frequently occur in other approaches . However, these applications did not deal with strong–coupling problems as defined above, which would be a very interesting perspective for this new method. Głazek and Wilson undertook a first step in this direction in Ref. . They investigated a quadratic Hamiltonian that shows strong–coupling behavior due to the formation of a bound state from a continuum, and demonstrated how this model can be solved using infinitesimal unitary transformations. However, since they dealt with a quadratic Hamiltonian, this was not a true many–particle strong–coupling problem as would be of most interest in condensed matter theory or high energy theory. Recently, I described the application of the flow equation method to the one–dimensional quantum sine–Gordon model . The sine–Gordon model is a many–particle problem with an interesting phase structure including a strong–coupling regime. It was shown in Ref. that it is possible to use the flow equation scheme to develop a systematic expansion that links weak– to strong–coupling behavior in a controlled way. Already the leading order of this expansion was is close agreement with exact results. In the present paper I will present the various details of the calculation not included in the original Letter in a self–contained manner. The sine–Gordon model is defined by the Hamiltonian $$H=𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2+u\tau ^2\mathrm{cos}\left[\beta \varphi (x)\right]\right),$$ (2) where $`\varphi (x)`$ is a bosonic field and $`\mathrm{\Pi }(x)`$ its conjugate momentum field with the commutator $`[\mathrm{\Pi }(x),\varphi (x^{})]=i\delta (xx^{})`$. $`u>0`$ is a small dimensionless coupling constant and $`\mathrm{\Lambda }\tau ^1`$ an implicit UV–cutoff. We are interested in the universal properties for energies $`\|E\|\mathrm{\Lambda }`$. The sine–Gordon model exhibits a strong–coupling phase for $`\beta ^28\pi `$ with a mass gap and fermionic low–energy excitations (massive solitons). The perturbative scaling analysis leads to a characteristic strong–coupling divergence of the running coupling $`u`$ in this regime. This makes the sine–Gordon model an interesting test model for our new approach. The main emphasis in this paper will not lie in deriving new results in this well–studied model, but in showing how these results follow within the flow equation method, and how therefore our new method can be useful for strong–coupling problems more generally. Other features that make the sine–Gordon model an attractive test model are its interesting phase structure with a Kosterlitz–Thouless type transition to a phase with massless solitons at $`\beta ^2/8\pi 1+O(u)`$, its integrable structure that allows the comparison with exact results , and its relation to a variety of other models like the spin-$`1/2`$ X-Y-Z chain, the $`1d`$ electron gas with backward scattering, the Thirring model in field theory and the $`2d`$ Coulomb gas (for an overview of these relations see Ref. ). Therefore the results from the flow equation approach can be viewed within a variety of model contexts. The main motivation for being interested in the flow equation approach to this integrable model lies, however, in the observation that our new method does not make use of the integrable structure. In our approach a small parameter is identified and used within a suitably renormalized perturbation expansion. The usual perturbative scaling approach fails because the initially small expansion parameter $`u`$ diverges during the RG–procedure. In the flow equation approach the expansion parameter will turn out to be the product of the running coupling $`u`$ and a factor $`(1+\beta ^2/4\pi )`$. This combination will always remain small during the flow. It is therefore feasible to study for example nonintegrable perturbations and correlation functions within our new approach, which should be of considerable interest in a variety of contexts. Although the calculations presented here appear rather lengthy and technical at first, they are straightforward and much closer to conventional many–body techniques than methods building on the integrable structure. ### B Outline The structure of this paper is as follows. Sect. II deals with some general properties of the sine–Gordon model that are important in the sequel. In Sect. II.A the sine–Gordon model and the regularization used in this paper are introduced. Sect. II.B reviews the perturbative scaling analysis, the phase structure, and the strong–coupling behavior. In Sect. II.C various exact results based on the integrable structure of the sine–Gordon model are summed up, especially properties of the point $`\beta ^2=4\pi `$ where the model becomes equivalent to a noninteracting Thirring model. This equivalence will play an important role in understanding the structure of our flow equation approach later on. After setting the stage in Sect. II, Sect. III deals with the actual application of the flow equation approach to the sine–Gordon model. Some general properties of the flow equation method are reviewed in III.A. Then the appropriate generator $`\eta (B)`$ for the sine–Gordon model diagonalization is worked out in III.B and the commutator $`[\eta (B),H_0]`$ evaluated in III.C. The key computational parts of the flow equation approach are contained in III.D and III.E, where the commutators $`[\eta (B),H_{\mathrm{int}}(B)]`$ and $`[\eta (B),H_{\mathrm{diag}}(B)]`$ are evaluated. From these commutators the flow of $`\beta ^2(B)`$ and of the running coupling constant are deduced in Sect. III.F. For the convenience of the reader, all the results from this technical part are summed up in Sect. IV.A, in particular the Hamiltonian $`H(B)`$ along the flow and the set of flow equations governing the various parameters in $`H(B)`$. We will see in IV.B that in the strong–coupling phase $`H(B)`$ flows to an effective low–energy noninteracting Thirring model. The mass gap of the sine–Gordon model can be easily deduced from this low–energy model, and the results are then compared with perturbative scaling analysis and exact integrable model results. The agreement will turn out to be very good. In Sect. IV.C the final diagonal Hamiltonian $`H(B=\mathrm{})`$ is discussed in more detail, in particular the soliton dispersion relation and properties in the crossover region. Finally in Sect. IV.D the approximations and the expansion parameter of our approach are reviewed. Sect. V sums up the conclusions and an outlook to open questions. The Appendix contains important properties of vertex operators that are used throughout this paper. ## II Sine–Gordon Model ### A Definition The one–dimensional quantum sine–Gordon model is defined by the Hamiltonian $$H=𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2+u\tau ^2\mathrm{cos}\left[\beta \varphi (x)\right]\right).$$ (3) $`\varphi (x)`$ is a bosonic field and $`\mathrm{\Pi }(x)`$ its conjugate momentum field with the fundamental commutator $$[\mathrm{\Pi }(x),\varphi (x^{})]=i\delta (xx^{}).$$ (4) In (3) an UV–momentum cutoff $`\mathrm{\Lambda }\tau ^1`$ is implied. $`u`$ is a dimensionless coupling constant. Without loss of generality we will assume $`u>0`$ and $`\beta >0`$. Expanding the fields in normal modes gives $`\varphi (x)`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{4\pi }}}{\displaystyle \underset{k0}{}}{\displaystyle \frac{\sqrt{\|k\|}}{k}}e^{ikx}\left(\sigma _1(k)+\sigma _2(k)\right)`$ (5) $`\mathrm{\Pi }(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{4\pi }}}{\displaystyle \underset{k0}{}}\sqrt{\|k\|}e^{ikx}\left(\sigma _1(k)\sigma _2(k)\right).`$ (6) Sums over wavevectors $`k,p,q,\mathrm{}`$ are to be understood in the sense $$\underset{k}{}\stackrel{\mathrm{def}}{=}\frac{2\pi }{L}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}$$ (7) with $`k=2\pi n/L`$ throughout this paper. $`L`$ is the system size. The basic commutators for $`k,k^{}>0`$ are $`[\sigma _1(k),\sigma _1(k^{})]=[\sigma _2(k),\sigma _2(k^{})]`$ $`=`$ $`\delta _{kk^{}}L/2\pi `$ (8) $`[\sigma _j(k),\sigma _j(k^{})]=[\sigma _j(k),\sigma _j(k^{})]`$ $`=`$ $`0`$ (9) and for $`jj^{}`$ $$[\sigma _j(k),\sigma _j^{}(k^{})]=[\sigma _j(k),\sigma _j^{}(k^{})]=[\sigma _j(k),\sigma _j^{}(k^{})]=0.$$ (10) All other commutators can be derived via $`\sigma _i^{}(k)=\sigma _i^{}(k)`$. The vacuum $`\|\mathrm{\Omega }`$ is defined by $$\sigma _1(k)\|\mathrm{\Omega }=\sigma _2(k)\|\mathrm{\Omega }=0$$ (11) for all $`k>0`$. The notion of the dual field $`\mathrm{\Theta }(x)`$ will also be useful. $`\mathrm{\Theta }(x)`$ is defined by $$_x\mathrm{\Theta }(x)=\mathrm{\Pi }(x),$$ (12) leading to the commutator $$[\mathrm{\Theta }(x),\varphi (x^{})]=i\theta (xx^{}).$$ (13) In terms of normal modes one finds $$\mathrm{\Theta }(x)=\frac{i}{\sqrt{4\pi }}\underset{k0}{}\frac{\sqrt{\|k\|}}{k}e^{ikx}\left(\sigma _1(k)\sigma _2(k)\right).$$ (14) The concept of vertex operators will play an important role in the sequel. Vertex operators $`V_j(\alpha ;x)`$ are defined as normal–ordered exponentials $$V_j(\alpha ;x)=:\mathrm{exp}(\pm \alpha \underset{p0}{}\frac{\sqrt{\|p\|}}{p}e^{\frac{a}{2}\|p\|ipx}\sigma _j(p)):$$ (15) with $`+`$ (upper sign) corresponding to $`j=1`$ and $``$ (lower sign) to $`j=2`$. This sign convention will be used throughout this paper. Normal ordering $`:\mathrm{}:`$ amounts to commuting all the operators that annihilate the vacuum according to (11) to the right. One can rewrite (15) in terms of the field and its dual (12) $$V_j(\alpha ;x)=:\mathrm{exp}(\pm i\alpha \sqrt{\pi }dϵc(ϵ)[\varphi (x+ϵ)\pm \mathrm{\Theta }(x+ϵ)]):$$ (16) with the Lorentzian $$c(ϵ)=\frac{a/2\pi }{ϵ^2+a^2/4}.$$ (17) $`c(ϵ)`$ is normalized $$_{\mathrm{}}^{\mathrm{}}𝑑ϵc(ϵ)=1$$ (18) and $`c(ϵ)\stackrel{a0}{}\delta (ϵ)`$. Further properties of vertex operators, in particular their operator product expansion (OPE), are reviewed in the Appendix. One can rewrite the interaction term of the sine–Gordon model (3) in terms of vertex operators $$H=𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2+\frac{u}{2\pi a^2}\left(\frac{2\pi a}{L}\right)^{\alpha ^2}\left(V_1(\alpha ;x)V_2(\alpha ;x)+V_2(\alpha ;x)V_1(\alpha ;x)\right)\right).$$ (19) The prefactor $`(2\pi a/L)^{\alpha ^2}`$ follows from (275). Here and in the sequel $`\alpha `$ and $`\beta `$ are used interchangeably with the identification $$\alpha \stackrel{\mathrm{def}}{=}\frac{\beta }{\sqrt{4\pi }}.$$ (20) Eq. (19) is the form of the sine–Gordon Hamiltonian that we will investigate with the flow equation approach: No implicit momentum cutoff is implied in (19): The UV–regularization of the Hamiltonian (19) is achieved by the cutoff parameter $`a>0`$ in the vertex operators. The regularizations in (19) and (3) are related by $`a^1\mathrm{\Lambda }\tau ^1`$. For a direct comparison between the original Hamiltonian (3) and the form (19) used here one can also identically express (19) as $$H=𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2+\frac{u}{\pi a^2}\mathrm{cos}\left[\beta 𝑑ϵc(ϵ)\varphi (x+ϵ)\right]\right).$$ (21) The regularization with $`a`$ therefore amounts to smearing out the interaction term. In the limit $`a0`$ one recovers the $`\mathrm{cos}(\beta \varphi (x))`$–interaction term. The universal properties of the sine–Gordon model for energies $`\|E\|a^1`$ are not affected by this choice of regularization. We find, however, notational simplifications and more compact expressions in the course of our calculation when we start with (19) (or equivalently (21)). In order to clarify the main conceptual ideas of the flow equation approach, it will therefore be convenient for us to use the regularization (19) with the UV–cutoff $`a^1`$ built in via the definition of the vertex operators. ### B Perturbative scaling analysis The flow equation approach can be viewed as an extension of perturbative scaling. Therefore it is useful to briefly review the results of the perturbative scaling analysis as applied to the sine–Gordon model. A comprehensive review can be found in Ref. . In 2-loop order there are two renormalization group equations that describe the flow of $`u`$ and $`\beta `$ upon integrating out the degrees of freedom with $`\mathrm{\Lambda }d\mathrm{\Lambda }<\|k\|<\mathrm{\Lambda }`$ (see Ref. ) $`{\displaystyle \frac{d\beta ^2}{d\mathrm{ln}\mathrm{\Lambda }}}`$ $`=`$ $`{\displaystyle \frac{u^2}{4\pi }}`$ (22) $`{\displaystyle \frac{du}{d\mathrm{ln}\mathrm{\Lambda }}}`$ $`=`$ $`\left({\displaystyle \frac{\beta ^2}{4\pi }}2\right)u.`$ (23) The initial conditions are $`u(\tau ^1)=u`$ and $`\beta (\tau ^1)=\beta `$. These scaling equations give rise to the Kosterlitz–Thouless type phase diagram shown in Fig. 1. The two separatrices $`S_\pm `$ originating from $`\beta ^2=8\pi `$ with $`\beta ^2=8\pi (1\pm u)`$ for small $`u`$ divide the parameter space in three sectors: 1. the weak–coupling sector I; 2. the crossover sector II; 3. the asymptotic freedom sector III. Both in II and III the perturbative scaling equations (23) lead to strong–coupling behavior with the running coupling constant $`u`$ growing larger and larger during the flow. Therefore the perturbative RG approach eventually becomes invalid in these sectors. This indicates the opening of mass gap in the spectrum in II and III. Across $`S_+`$ the system undergoes a Kosterlitz–Thouless type phase transition between this massive phase and the massless phase in I. In spite of the strong–coupling divergence in II and III, the perturbative scaling equations allow us to analyze the size of the mass gap by identifying the mass $`M`$ with the scaling invariant of (23). One e.g. finds the following expressions for small $`u>0`$ $`M\mathrm{\Lambda }\left({\displaystyle \frac{u}{2\beta ^2/4\pi }}\right)^{1/(2\beta ^2/4\pi )}`$ $`\text{for }{\displaystyle \frac{1\beta ^2/8\pi }{u}}1`$ (24) $`M\mathrm{\Lambda }\mathrm{exp}\left({\displaystyle \frac{1}{2\beta ^2/4\pi }}\right)`$ $`\text{for }(1{\displaystyle \frac{\beta ^2}{8\pi }}u)^11\text{(along }S_{})`$ (25) $`M\mathrm{\Lambda }\mathrm{exp}\left({\displaystyle \frac{\pi }{2\sqrt{u^2(1\beta ^2/8\pi )^2}}}\right)`$ $`\text{for }(1{\displaystyle \frac{\beta ^2}{8\pi }}+u)^11\text{(along }S_+)`$ (26) In this manner one can obtain information about the renormalized low–energy scale even in the strong–coupling phase. But the perturbative scaling approach does by itself not lead to an understanding of the physical behavior associated with this low–energy scale. This situation is typical for other strong–coupling problems as well, the Kondo model being the paradigm in condensed matter theory . In combination with mappings to other exactly solvable models these shortcomings can sometimes be partially overcome, see Sect. II.C below. However, in general there is considerable interest in theoretical methods that can solve strong–coupling problems in a controlled way. Therefore the flow equation approach might be an interesting tool also for other strong–coupling problems by removing some of the above shortcomings. The phase diagram Fig. 1 remains essentially unchanged in higher loop orders . For latter comparison with the flow equation solution it is interesting to also write down the 3-loop result for the mass gap on $`S_{}`$ (that is for $`\beta ^2=8\pi (1u)`$) in the limit $`u0`$ $$M\mathrm{\Lambda }u^{1/2}\mathrm{exp}\left(\frac{1}{2u}\right).$$ (27) Notice the $`u^{1/2}`$–prefactor as compared to the 2-loop result (25). Higher loop orders beyond 3-loop should only affect the proportionality factor in (27. ### C Integrable structure and relation to other models The sine–Gordon model is one of the best studied integrable models, which makes it a very suitable test model for our new approach. Its spectrum was obtained exactly from an inverse scattering solution and its $`S`$–matrix was calculated by Zamolodchikov . For a recent review see Ref. . In the strong–coupling phase the exact solution confirms the scaling behavior (24) for the mass $`M`$ of the solitons and antisolitons. The exact $`S`$–matrix also shows that for small rapidity differences these solitons and antisolitons behave as fermions. This important observation of a change in statistics for the low–energy excitations will be reproduced in our flow equation framework. In addition, the exact solution shows that new features appear for $`\beta ^2<4\pi `$: Soliton–anitsoliton bound states (breathers) emerge in the spectrum with excitation energies smaller than $`2M`$. There is one breather for $`8\pi /3\beta ^2<4\pi `$, two breathers for $`2\pi \beta ^2<8\pi /3`$ etc. . The sine–Gordon model is related to other integrable models like the spin-1/2 X-Y-Z chain and the Thirring model. Since in particular the relation to the Thirring model will be fruitful in the sequel, it is useful to sum up some of its main properties here: The massive $`1d`$ Thirring model is defined by the Lagrangian density $$_{\mathrm{Th}}=\underset{\mu =0}{\overset{1}{}}\left(\overline{\psi }i\gamma _\mu ^\mu \psi \frac{1}{2}gj_\mu j^\mu \right)M\overline{\psi }\psi $$ (28) in terms of the two–component spinor $`\psi (x)`$ $$\psi (x)=\left(\begin{array}{c}\hfill \psi _1(x)\\ \hfill \psi _2(x)\end{array}\right)$$ (29) with components obeying fermionic anticommutation relations $$\{\psi _j^{}(x),\psi _j^{}^{}(y)\}=\delta _{jj^{}}\delta (xy),\{\psi _j^{}(x),\psi _j^{}^{}(y)\}=\{\psi _j^{}(x),\psi _j^{}^{}(y)\}=0.$$ (30) The current is defined by $`j^\mu \stackrel{\mathrm{def}}{=}\overline{\psi }\gamma ^\mu \psi `$ with $`\overline{\psi }\stackrel{\mathrm{def}}{=}\psi ^{}\gamma _0`$ and the $`\gamma `$–matrices are explicitly given by $$\gamma _0=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right)\text{and}\gamma _1=\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right).$$ (31) The exact Bethe ansatz solution of the Thirring model was obtained by Bergknoff and Thacker . Coleman has shown the equivalence of the sine–Gordon model with the Thirring model (28) order by order in perturbation theory with the following mapping between the coupling constants $$\frac{\beta ^2}{4\pi }=\frac{1}{1+g/\pi }.$$ (32) One notices that $`\beta ^2=4\pi `$ is a special point of the sine–Gordon model since it corresponds to a noninteracting massive Thirring model ($`g=0`$). For $`\beta ^2=4\pi `$ the elementary excitations of the sine–Gordon model are therefore fermionic with the dispersion relation $`\pm E_k`$ with $$E_k=\sqrt{k^2+M^2}.$$ (33) The explicit relation between the bosonic field of the sine–Gordon model and the fermionic field of the Thirring model was found by Mandelstam . For $`\beta ^2=4\pi `$ one has explicitly $$\psi _j(x)=\frac{1}{\sqrt{L}}V_j(1;x)$$ (34) with $`V_j(\alpha ;x)`$ from (15): Notice that the $`V_j(\pm 1;x)`$ obey anticommutation relations (300) in the limit $`a0`$. These Thirring fermionsNotice, however, that $`V_1(\pm 1;x)`$ commutes with $`V_2(\pm 1;x)`$ instead of anticommuting. “Proper” fermions can easily be defined with an additional Jordan–Wigner phase factor, but nothing new can be learned from this Jordan–Wigner construction. correspond to the quantized soliton solutions of the sine–Gordon model . The perturbative scaling approach does not “know” about the special point $`\beta ^2=4\pi `$ where the sine–Gordon model is trivially diagonalizable by using the equivalence to the quadratic Thirring model: The strong–coupling scaling trajectories in Fig. 1 go right through the line $`\beta ^2=4\pi `$. The subsequent strong–coupling divergence of the running coupling constant is then due to the fact that one has generated the nonvanishing energy scale $`M`$. A similar scenario occurs in the Kondo model: There the diagonal Hamiltonian corresponds to the Toulouse point and the nonvanishing energy scale is the Kondo temperature $`T_K`$. A standard approach to avoid the strong–coupling divergence is to scale the model to the exactly solvable line $`\beta ^2=4\pi `$: One stops the scaling once $`\beta ^2(\mathrm{\Lambda }_{\mathrm{eff}})=4\pi `$ and obtains the mass from the value of the running coupling constant. With this approach it is also plausible that the low–energy single–particle/hole excitations in the strong–coupling phase are fermionic with a mass set by the scaling invariant. Though very useful, this is an uncontrolled approximation since the running coupling constant is already large when $`\beta ^2(\mathrm{\Lambda }_{\mathrm{eff}})=4\pi `$. It is therefore difficult/impossible to learn something about the crossover from weak–coupling to strong–coupling or about the effect of irrelevant operators at the strong–coupling fixed point. These shortcomings make it desirable to develop our new method that will allow a systematic expansion describing the full crossover flow. The sine–Gordon model is also related to a variety of other models like the $`2d`$ Coulomb gas with temperature $`T=\beta ^2`$ and fugacity $`zu`$, or a $`1d`$ electron gas with backward scattering. For an overview of these and other relations see Ref. . As a final remark we also want to mention that the mapping to the $`1d`$ electron gas gives a natural interpretation to the separatrices $`S_\pm `$ in Fig. 1 since they correspond to an electron gas with SU(2) spin–symmetric interactions . The sine–Gordon model with $`\beta ^2=8\pi (1\pm u)`$, $`\|u\|1`$ therefore carries a hidden SU(2)–symmetry. ## III Flow equation approach ### A General concepts The idea to apply a sequence of infinitesimal unitary transformations to a Hamiltonian in order to make it more diagonal has been independently put forward by Wegner and Głazek and Wilson . Wegner’s original work focussed on diagonalizing many–particle Hamiltonians, whereas the focus in the work of Głazek and Wilson was to construct effective low–energy Hamiltonians for strong–coupling field theories: Such effective Hamiltonians can then be analyzed by standard techniques in order to find the bound state spectrum, which in turn could be interpreted as e.g. hadrons or mesons. Though the outlook of these approaches is somehow different, the concepts are very similar. In this paper we will follow Wegner’s methodology. The main idea of Wegner’s flow equations is to generate a one–parameter family of unitarily equivalent Hamiltonians $`H(B)`$ labelled by a flow parameter $`B`$.The flow parameter has been denoted by $`\mathrm{}`$ in most other works on flow equations. In order to avoid confusion with the common notation where $`\mathrm{}`$ is the logarithm of the change in length scale in RG–equations, $`B`$ instead of $`\mathrm{}`$ is used in this work. This is achieved by solving a differential equation $$\frac{dH(B)}{dB}=[\eta (B),H(B)]$$ (35) with some anti-Hermitian generator $`\eta (B)=\eta ^{}(B)`$ where $`H(B=0)=H`$ is the initial Hamiltonian. One wants to choose $`\eta (B)`$ such that $`H(B)`$ becomes more diagonal as $`B\mathrm{}`$: Splitting up $`H(B)`$ in its diagonal and interaction parts $$H(B)=H_0(B)+H_{\mathrm{int}}(B),$$ (36) this amounts to requiring $`H_{\mathrm{int}}(B)`$ becomes (in some sense) smaller for $`B\mathrm{}`$. In order to achieve this, Wegner proposed the following generator $$\eta (B)\stackrel{\mathrm{def}}{=}[H_0(B),H_{\mathrm{int}}(B)].$$ (37) With this choice of $`\eta (B)`$ one can show $$\frac{d}{dB}\mathrm{Tr}H_{\mathrm{int}}^2(B)0$$ (38) and in this sense the operator $`H_{\mathrm{int}}(B)`$ becomes smaller along the flow. Notice that $`B`$ has the dimension of (Energy)<sup>-2</sup> with this choice. However, for a many–particle Hamiltonian Eq. (38) is usually not well–defined since the trace is typically infinite. Also higher and higher order interactions are successively generated by the system of equations (35) and (37), which have to be truncated in some way making rigorous statements difficult. Still one finds that (37) is generally a suitable choice for achieving our goal to make the initial Hamiltonian diagonal if $`H_{\mathrm{int}}(B=0)`$ can be viewed as a small perturbation term: Truncating the system of higher order interactions produced by (35) and (37) in some order of the coupling constant then amounts to generating a perturbation expansion in a renormalized coupling constant. From this point of view the flow equation approach is similar to perturbative RG. Matrix elements of $`H_{\mathrm{int}}(B=0)`$ that couple states with large energy differences are eliminated in the initial stages of the flow (for small $`B`$), and matrix elements coupling more degenerate states are eliminated later. This is reminiscent of the energy scale separation underlying the renormalization group approach, which is the suitable perturbation expansion for systems with largely varying energy scales. Explicit applications of these ideas have been discussed for various model Hamiltonians like the $`n`$–orbital model , dissipative quantum systems , systems with electron–phonon coupling and various other models in condensed matter physics . In the present paper it will be shown how this method can be used for the sine–Gordon model as a genuine strong–coupling many–body Hamiltonian. ### B Generator $`\eta `$ The aim of this work is to diagonalize the sine–Gordon model (19) using the method of infinitesimal unitary transformations outlined above. We split up the sine–Gordon Hamiltonian $`H(B)`$ into a free part $`H_0`$ and the interaction part $`H_{\mathrm{int}}(B)`$ $$H(B)=H_0+H_{\mathrm{int}}(B)$$ (39) with $`H_0`$ $`=`$ $`{\displaystyle 𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2\right)}`$ (40) $`=`$ $`{\displaystyle \underset{p>0}{}}p\left(\sigma _1(p)\sigma _1(p)+\sigma _2(p)\sigma _2(p)\right)`$ (41) $`H_{\mathrm{int}}(B)`$ $`=`$ $`{\displaystyle }dxdyu(B;y)(V_1(\alpha ;x)V_2(\alpha ;xy)+\mathrm{h}.\mathrm{c}.).`$ (42) In order to avoid confusion the initial parameters $`u`$ and $`\beta `$ in (19) will from now on be denoted as $`u_0`$ and $`\beta _0`$. The initial condition for (42) then reads $`\alpha `$ $`=`$ $`{\displaystyle \frac{\beta _0}{\sqrt{4\pi }}}`$ (43) $`u(B=0;y)`$ $`=`$ $`{\displaystyle \frac{u_0}{2\pi a^2}}\delta (y)\left({\displaystyle \frac{2\pi a}{L}}\right)^{\alpha ^2}.`$ (44) Notice that we have already allowed for a general nonlocal interaction $`u(B;y)`$ in (42) since the initially local interaction (44) will become nonlocal along the flow (see below). Next we have to evaluate (37). The following commutator is useful $$[\sigma _j(p),V_j^{}(\alpha ;x)]=\delta _{jj^{}}\alpha \frac{\sqrt{\|p\|}}{p}\mathrm{exp}\left(\frac{a}{2}\|p\|+ipx\right)V_j^{}(\alpha ;x)$$ (45) leading to $`[{\displaystyle \underset{p>0}{}}p\sigma _1(p)\sigma _1(p),V_1(\alpha ;x)]`$ $`=`$ $`i_xV_1(\alpha ;x)`$ (46) $`[{\displaystyle \underset{p>0}{}}p\sigma _2(p)\sigma _2(p),V_2(\alpha ;x)]`$ $`=`$ $`i_xV_2(\alpha ;x).`$ (47) Thus we find the following generator $`\eta (B)`$ $`=`$ $`[H_0,H_{\mathrm{int}}(B)]`$ (48) $`=`$ $`2i{\displaystyle }dxdy{\displaystyle \frac{u(B;y)}{y}}(V_1(\alpha ;x)V_2(\alpha ;xy)+\mathrm{h}.\mathrm{c}.).`$ (49) ### C Commutator $`[\eta ,H_0]`$ To study the flow generated by $`\eta `$ we first look at the commutator $`[\eta ,H_0]`$. Using (47) one easily shows $$[\eta ,H_0]=4dxdy\frac{^2u(B;y)}{y^2}(V_1(\alpha ;x)V_2(\alpha ;xy)+\mathrm{h}.\mathrm{c}.).$$ (50) Comparison of the coefficients on the left–hand side of (35) with (50) gives $$\frac{u(B;y)}{B}=4\frac{^2u(B;y)}{y^2},$$ (51) where possible contributions from $`[\eta ,H_{\mathrm{int}}]`$ are still missing. Eq. (51) has the character of a diffusion equation: The initially local interaction becomes increasingly non–local along the flow. In terms of Fourier coefficients $$u(B;y)=\underset{p}{}u(B;p)e^{ipy}$$ (52) one finds the solution $$u(B;p)=\frac{u_0}{4\pi ^2a^2}e^{4p^2B}\left(\frac{2\pi a}{L}\right)^{\alpha ^2}.$$ (53) One sees explicitly that matrix elements $`u(B;p)`$ coupling states with large energy differences $`\|p\|`$ are eliminated in the early stages of the flow (for small $`B`$), whereas matrix elements coupling more degenerate states are decoupled later during the flow. This is a generic feature of Wegner’s generator (37). Later we will see that (53) is modified due to higher–order contributions. Therefore we introduce a more general parametrization $$u(B;p)=\frac{\stackrel{~}{u}(B)}{4\pi ^2a^2}\left(\frac{2\pi a}{L}\right)^{\alpha ^2(B)}v(B;p)$$ (54) with a running coupling $`\stackrel{~}{u}(B)`$, initially $`\stackrel{~}{u}(B=0)=u_0`$. The differential equation for the coefficients $`v(B;p)`$ now reads $$\frac{v(B;p)}{B}=4p^2v(B;p)$$ (55) with the initial condition $`v(B=0;p)=1`$. ### D Commutator $`[\eta ,H_{\mathrm{int}}]`$ #### 1 General properties The evaluation of $`[\eta ,H_{\mathrm{int}}]`$ is the key calculation in the flow equation approach. We first look at some general properties of such commutators. Let $`A_1,A_2,B_1,B_2`$ be arbitrary operators with $`[A_j,B_j^{}]=0`$. We define $`O`$ as the operator with its ground state expectation value subtracted $$O\stackrel{\mathrm{def}}{=}OO,$$ (56) with the notation $`O\stackrel{\mathrm{def}}{=}\mathrm{\Omega }\|O\|\mathrm{\Omega }`$. One easily shows $`[A_1B_1,A_2B_2]`$ $`=`$ $`A_1A_2B_1B_2A_2A_1B_2B_1`$ (60) $`+B_1B_2A_1A_2B_2B_1A_2A_1`$ $`+A_1A_2B_1B_2A_2A_1B_2B_1`$ $`+R`$ with $$R=A_1A_2B_1B_2A_2A_1B_2B_1.$$ (61) In general $`R`$ leads to the generation of higher order interaction terms during the flow. $`R`$ vanishes if the operators fulfill the following exchange relations $`A_1A_2+e^{i\varphi }A_2A_1`$ $`=`$ $`c`$ (62) $`B_1B_2+e^{i\varphi }B_2B_1`$ $`=`$ $`c`$ (63) with fixed $`\varphi `$ and $`c`$. E.g. for $`\varphi =0`$ these are fermionic anticommutation relations, or for $`\varphi =\pi `$ bosonic commutation relations. Then no higher–order interactions are generated and it is possible to close the flow equations without approximations. For general $`\beta _0`$ in the interaction term we will, however, have to develop a suitable approximation for $`R`$ in the next section. #### 2 $`[\eta ,H_{\mathrm{int}}]`$ in the sine–Gordon model There are two structurally different commutators of vertex operators generated by $`[\eta ,H_{\mathrm{int}}]`$ in the sine–Gordon model: $`[V_1^{}V_2^{},V_1^{}V_2^{}]`$ and $`[V_1^{}V_2^{},V_1^{}V_2^{}]`$ (or equivalently $`[V_1^{}V_2^{},V_1^{}V_2^{}]`$). The first term $$[V_1(\alpha ;x_1)V_2(\alpha ;x_1y_1),V_1(\alpha ;x_2)V_2(\alpha ;x_2y_2)]$$ (64) and its hermitian conjugate will turn out to be the leading contributions and are discussed first. Eqs. (290) and (291) give $`V_1(\alpha ;x_1)V_1(\alpha ;x_2)=s_1^{\alpha ^2}`$ , $`V_1(\alpha ;x_2)V_1(\alpha ;x_1)=\overline{s}_1^{\alpha ^2}`$ (65) $`V_2(\alpha ;x_1y_1)V_2(\alpha ;x_2y_2)=s_2^{\alpha ^2}`$ , $`V_2(\alpha ;x_2y_2)V_2(\alpha ;x_1y_1)=\overline{s}_2^{\alpha ^2}`$ (66) with $`s_1`$ $`=`$ $`{\displaystyle \frac{2\pi }{L}}\left(i(x_2x_1)+a\right)`$ (67) $`s_2`$ $`=`$ $`{\displaystyle \frac{2\pi }{L}}\left(i(x_1y_1x_2+y_2)+a\right).`$ (68) Using (60) we then find $`[V_1(\alpha ;x_1)V_2(\alpha ;x_1y_1),V_1(\alpha ;x_2)V_2(\alpha ;x_2y_2)]`$ (69) $`=`$ $`s_1^{\alpha ^2}s_2^{\alpha ^2}\overline{s}_1^{\alpha ^2}\overline{s}_2^{\alpha ^2}`$ (73) $`s_2^{\alpha ^2}V_1(\alpha ;x_1)V_1(\alpha ;x_2)\overline{s}_2^{\alpha ^2}V_1(\alpha ;x_2)V_1(\alpha ;x_1)`$ $`s_1^{\alpha ^2}V_2(\alpha ;x_1y_1)V_2(\alpha ;x_2y_2)\overline{s}_1^{\alpha ^2}V_2(\alpha ;x_2y_2)V_2(\alpha ;x_1y_1)`$ $`+R`$ with $`R`$ $`=`$ $`V_1(\alpha ;x_1)V_1(\alpha ;x_2)V_2(\alpha ;x_1y_1)V_2(\alpha ;x_2y_2)`$ (75) $`V_1(\alpha ;x_2)V_1(\alpha ;x_1)V_2(\alpha ;x_2y_2)V_2(\alpha ;x_1y_1).`$ The key approximation in our method is to use an operator product expansion (OPE) in higher–order interaction terms like $`R`$, and then to neglect contributions with larger scaling dimensions (more irrelevant terms in the RG–sense). From (296) we e.g. conclude $$V_1(\alpha ;x_1)V_1(\alpha ;x_2)=s_1^{\alpha ^2}(i\alpha (x_2x_1)\underset{p0}{}\sqrt{\|p\|}e^{\frac{a}{2}\|p\|ipx_1}\sigma _1(p)+\mathrm{}),$$ (76) where we have neglected the higher order terms in (296). Notice that the c–number contribution has already been removed by subtracting the ground state expectation value. Putting everything together gives $$R=\alpha ^2(x_2x_1)(x_1y_1x_2+y_2)(s_1^{\alpha ^2}s_2^{\alpha ^2}\overline{s}_1^{\alpha ^2}\overline{s}_2^{\alpha ^2})\underset{p,q0}{}\sqrt{\|pq\|}e^{\frac{a}{2}(\|p\|+\|q\|)ipx_1iq(x_1y_1)}\sigma _1(p)\sigma _2(q).$$ (77) The first and second term in (69) are c–numbers and describe a shift in the ground state energy. This is of no particular interest and we will not look into it. The various other terms generated in $`[\eta ,H_{\mathrm{int}}]`$ are discussed in the next subsections. #### 3 $`R`$–term The $`R`$–term in (69) leads to the following contribution from $`[\eta ,H_{\mathrm{int}}]`$ $`2i{\displaystyle 𝑑x_1𝑑x_2𝑑y_1𝑑y_2\frac{u(B;y_1)}{y_1}u(B;y_2)}`$ (79) $`\times 2\alpha ^2(x_2x_1)(x_1y_1x_2+y_2)(s_1^{\alpha ^2}s_2^{\alpha ^2}\overline{s}_1^{\alpha ^2}\overline{s}_2^{\alpha ^2}){\displaystyle \underset{p,q0}{}}\sqrt{\|pq\|}e^{\frac{a}{2}(\|p\|+\|q\|)ipx_1iq(x_1y_1)}\sigma _1(p)\sigma _2(q)`$ $`=`$ $`8\pi i\alpha ^2\left({\displaystyle \frac{L}{2\pi }}\right)^{2\alpha ^2}{\displaystyle \underset{k0}{}}\|k\|t_k\sigma _1(k)\sigma _2(k)`$ (80) with coefficients $`t_k`$ $$t_k=dz_1dz_2dz_3e^{a\|k\|ikz_1}\frac{u(B;z_1)}{z_1}u(B;z_1+z_2)z_3(z_2z_3)((iz_3+a)^{\alpha ^2}(i(z_2z_3)+a)^{\alpha ^2}\mathrm{h}.\mathrm{c}.).$$ (81) Except for an (unimportant) initial transient where $`Ba^2`$, the $`z_1`$–integral leads to the following expression $$t_k=2\pi ia^{42\alpha ^2}\underset{p}{}pu(B;p)u(B;kp)𝑑xe^{i(k+p)ax}I(x)$$ (82) with $$I(x)=𝑑yy(xy)\left((1+iy)^{\alpha ^2}(1iy+ix)^{\alpha ^2}(1iy)^{\alpha ^2}(1+iyix)^{\alpha ^2}\right).$$ (83) Writing $`1+i\left(y\pm \frac{x}{2}\right)=r_\pm e^{i\varphi _\pm }`$ with $`r_\pm `$ $`=`$ $`\sqrt{1+\left(y\pm {\displaystyle \frac{x}{2}}\right)^2}`$ (84) $`\varphi _\pm `$ $`=`$ $`\mathrm{arcsin}{\displaystyle \frac{y\pm \frac{x}{2}}{r_\pm }}[{\displaystyle \frac{\pi }{2}},{\displaystyle \frac{\pi }{2}}]`$ (85) leads to $$I(x)=2i𝑑y\left(\frac{x^2}{4}y^2\right)(r_+r_{})^{\alpha ^2}\mathrm{sin}\left(\alpha ^2(\varphi _{}\varphi _+)\right).$$ (86) The flow is dominated by the term decaying most slowly with $`B`$, which corresponds to the large-$`x`$ behavior of $`I(x)`$. Therefore we can approximate $$r_\pm \left\|\frac{x}{2}\right\|\|z\pm 1\|,$$ (87) where $`z=2y/x`$, and find $$I(x)=2i\left\|\frac{x}{2}\right\|^{32\alpha ^2}_{\mathrm{}}^{\mathrm{}}𝑑z(1z^2)\|1z^2\|^{\alpha ^2}\mathrm{sin}(\alpha ^2(\varphi _{}\varphi _+)).$$ (88) With the above approximation one has $`\varphi _\pm \{\frac{\pi }{2},\frac{\pi }{2}\}`$. Thus the only contributions to $`I(x)`$ come from regions with $`\varphi _{}\varphi _+`$ $`\varphi _{}\varphi _+=\pi `$ $``$ $`x>0,1<z<1`$ (89) $`\varphi _{}\varphi _+=\pi `$ $``$ $`x<0,1<z<1`$ (90) leading to $`I(x)`$ $`=`$ $`2i\mathrm{sin}(\alpha ^2\pi )\mathrm{sgn}(x)\left\|{\displaystyle \frac{x}{2}}\right\|^{32\alpha ^2}{\displaystyle _1^1}𝑑z(1z^2)^{1\alpha ^2}`$ (91) $`=`$ $`2i{\displaystyle \frac{\pi ^{3/2}}{\mathrm{\Gamma }(\alpha ^21)\mathrm{\Gamma }(\frac{5}{2}\alpha ^2)}}\mathrm{sgn}(x)\left\|{\displaystyle \frac{x}{2}}\right\|^{32\alpha ^2}.`$ (92) With (54) and (55) the sum over $`p`$ in (82) gives $$\underset{p}{}pu(B;p)u(B;kp)e^{ipax}=i\left(\frac{\stackrel{~}{u}(B)}{4\pi ^2a^2}\right)^2\left(\frac{2\pi a}{L}\right)^{2\alpha ^2}\sqrt{\frac{\pi }{8B}}\left(\frac{ax}{16B}+\frac{ik}{2}\right)\mathrm{exp}\left(\frac{a^2x^2}{32B}\frac{ikax}{2}2Bk^2\right).$$ (93) The final step is to perform the $`x`$–integration in (82). This can be done in closed form leading to hypergeometric functions. However, the flow is determined by the IR–limit $`k0`$ where the integral is simpler $$t_{k=0}=i\frac{32\pi ^3}{\mathrm{\Gamma }(\alpha ^21)}(32B)^{1\alpha ^2}\left(\frac{\stackrel{~}{u}(B)}{4\pi ^2a^2}\right)^2\left(\frac{2\pi a}{L}\right)^{2\alpha ^2}.$$ (94) For the full $`k`$–dependence we write $$t_k=t_{k=0}f(\alpha ^2;k\sqrt{B}).$$ (95) To leading order the only information that we will need about $`f(\alpha ^2;x)`$ is that it falls of rapidly to zero for large arguments $`\|x\|1`$. For example one easily shows $`f(\alpha ^2=1;x)`$ $`=`$ $`e^{4x^2}(18x^2)`$ (96) $`f(\alpha ^2=2;x)`$ $`=`$ $`e^{4x^2}\sqrt{2\pi }x\mathrm{erf}(\sqrt{2}x)e^{2x^2}.`$ (97) These functions are depicted in Fig. 2. Putting everything together the $`R`$–terms in (69) from $`[\eta ,H_{\mathrm{int}}]`$ contribute $$[\eta ,H_{\mathrm{int}}]\frac{32}{a^2}\left(\frac{32B}{a^2}\right)^{1\alpha ^2}\frac{\alpha ^2}{2\mathrm{\Gamma }(\alpha ^21)}\stackrel{~}{u}^2(B)\underset{k0}{}\|k\|f(\alpha ^2;k\sqrt{B})\sigma _1(k)\sigma _2(k).$$ (98) An important observation can be made for $`\alpha =1`$: Due to the divergent $`\mathrm{\Gamma }`$–function in the denominator, the term (98) vanishes for $`\beta _0^2=4\pi `$. This is an immediate consequence of the fact that in this case the vertex operators describe fermions. The interaction term of the sine–Gordon model is then simply a quadratic term in the fermions and no higher order interactions are generated during the flow, therefore according to (63) $`R0`$ for all $`B`$. For $`\beta _0^2=4\pi `$ we will be able to solve the flow equations without any approximations, thereby recovering the equivalence to the noninteracting Thirring model discussed in Sect. II.C. This demonstrates a fundamental difference of our approach to perturbative scaling, where the scaling trajectories go right through the line $`\beta ^2=4\pi `$ (see Fig. 1). #### 4 $`H_{\mathrm{diag}}`$ Let us next look at the fifth and sixth term in (69). The total contribution from $`[\eta ,H_{\mathrm{int}}]`$ to terms of this structure is $`[\eta ,H_{\mathrm{int}}]`$ $``$ $`2i{\displaystyle 𝑑x_1𝑑x_2𝑑y_1𝑑y_2\frac{u(B;y_1)}{y_1}u(B;y_2)}`$ (101) $`\times (s_1^{\alpha ^2}V_2(\alpha ;x_1y_1)V_2(\alpha ;x_2y_2)\overline{s}_1^{\alpha ^2}V_2(\alpha ;x_2y_2)V_2(\alpha ;x_1y_1)`$ $`+s_1^{\alpha ^2}V_2(\alpha ;x_1y_1)V_2(\alpha ;x_2y_2)\overline{s}_1^{\alpha ^2}V_2(\alpha ;x_2y_2)V_2(\alpha ;x_1y_1)).`$ For simplicity we will only look at the first term in this expression since all the other terms can be treated likewise. We first exchange the two vertex operators using (302) as this will lead to a normal–ordered expression below $$2idx_1dx_2dy_1dy_2\frac{u(B;y_1)}{y_1}u(B;y_2)s_1^{\alpha ^2}\frac{\overline{s}_2^{\alpha ^2}}{s_2^{\alpha ^2}}V_2(\alpha ;x_2y_2)V_2(\alpha ;x_1y_1)$$ (102) with $`s_1,s_2`$ from (68). We can rewrite this in terms of Fourier transforms ($`\alpha >0`$) $$V_j(\alpha ;x)\stackrel{\mathrm{def}}{=}\underset{p}{}e^{ipx}V_j(\alpha ;p),V_j(\alpha ;p)\stackrel{\mathrm{def}}{=}\left[V_j(\alpha ;p)\right]^{},$$ (103) substitute $`x_2x_2+x_1`$ and perform the integral over $`x_1`$. This leads to $`4\pi i\left({\displaystyle \frac{L}{2\pi }}\right)^{\alpha ^2}{\displaystyle 𝑑x_2𝑑y_1𝑑y_2\frac{u(B;y_1)}{y_1}u(B;y_2)}`$ (105) $`\times {\displaystyle \underset{k}{}}V_2(\alpha ;k)V_2(\alpha ;k)(ix_2+a)^{\alpha ^2}{\displaystyle \frac{[i(x_2y_2+y_1)+a]^{\alpha ^2}}{[i(x_2y_2+y_1)+a]^{\alpha ^2}}}e^{ik(x_2+y_1y_2)}`$ $`=`$ $`\mathrm{}`$ (106) $`=`$ $`8\pi ^2\left({\displaystyle \frac{L}{2\pi }}\right)^{\alpha ^2}{\displaystyle 𝑑x𝑑y\underset{k,p}{}pu^2(B;p)e^{ikx+ipy}V_2(\alpha ;k)V_2(\alpha ;k)[i(x+y)+a]^{\alpha ^2}\frac{[ix+a]^{\alpha ^2}}{[ix+a]^{\alpha ^2}}},`$ (107) where we have employed (52). Next the $`y`$–integration can be done using $$𝑑y(iy+ix+a)^{\alpha ^2}e^{ipy}=e^{ipx}\|p\|^{\alpha ^21}\frac{2\pi }{\mathrm{\Gamma }(\alpha ^2)}\theta (p),$$ (108) which is valid in the limit $`\|ap\|1`$: This holds except for an (unimportant) initial transient with wavevectors of order $`a^1`$. We find $$\frac{16\pi ^3}{\mathrm{\Gamma }(\alpha ^2)}\left(\frac{L}{2\pi }\right)^{\alpha ^2}𝑑x\underset{k}{}\underset{p>0}{}p^{\alpha ^2}u^2(B;p)e^{i(k+p)x}V_2(\alpha ;k)V_2(\alpha ;k)\frac{[ix+a]^{\alpha ^2}}{[ix+a]^{\alpha ^2}}.$$ (109) Next we can approximate $$\frac{[ix+a]^{\alpha ^2}}{[ix+a]^{\alpha ^2}}\mathrm{cos}(\pi \alpha ^2)+i\mathrm{sin}(\pi \alpha ^2)\mathrm{sgn}(x)$$ (110) using the same reasoning as in (304). This gives $`{\displaystyle \frac{16\pi ^3}{\mathrm{\Gamma }(\alpha ^2)}}\left({\displaystyle \frac{L}{2\pi }}\right)^{\alpha ^2}`$ $`(2\pi \mathrm{cos}(\pi \alpha ^2){\displaystyle \underset{k>0}{}}k^{\alpha ^2}u^2(B;k)V_2(\alpha ;k)V_2(\alpha ;k)`$ (112) $`+i\mathrm{sin}(\pi \alpha ^2){\displaystyle \underset{k}{}}{\displaystyle \underset{p>0}{}}p^{\alpha ^2}u^2(B;p){\displaystyle }dxe^{i(k+p)x}\mathrm{sgn}(x)V_2(\alpha ;k)V_2(\alpha ;k))`$ $`=2a^{\alpha ^24}{\displaystyle \frac{\stackrel{~}{u}^2(B)}{\mathrm{\Gamma }(\alpha ^2)}}\left({\displaystyle \frac{2\pi a}{L}}\right)^{\alpha ^2}`$ $`(\mathrm{cos}(\pi \alpha ^2){\displaystyle \underset{k>0}{}}k^{\alpha ^2}v^2(B;k)V_2(\alpha ;k)V_2(\alpha ;k)`$ (116) $`+{\displaystyle \frac{1}{\pi }}\mathrm{sin}(\pi \alpha ^2){\displaystyle \underset{k}{}}{\displaystyle \underset{\begin{array}{c}p>0\hfill \\ pk\hfill \end{array}}{}}p^{\alpha ^2}v^2(B;p){\displaystyle \frac{1}{pk}}V_2(\alpha ;k)V_2(\alpha ;k))`$ In order to do the sum over $`p`$ in the second term it will be sufficient to use the approximate solution $`v(B;p)=e^{4Bp^2}`$ from (55). Deviations from this approximate solution essentially only occur close to the strong–coupling fixed point $`\alpha ^2=1`$, where the second term vanishes anyway. This $`p`$–summation leads to an integral of the type $`h(\alpha ^2;x)`$ $`=`$ $`\mathrm{P}{\displaystyle _0^{\mathrm{}}}𝑑ye^{y^2}{\displaystyle \frac{y^{\alpha ^2}}{yx}}`$ (117) $`=`$ $`{\displaystyle \frac{1}{2}}e^{x^2}\|x\|^{\alpha ^2}\mathrm{Re}\left\{i^{\alpha ^2}\left[\mathrm{\Gamma }\left(1+{\displaystyle \frac{\alpha ^2}{2}}\right)\mathrm{\Gamma }({\displaystyle \frac{\alpha ^2}{2}},x^2)i\mathrm{sgn}(x)\mathrm{\Gamma }\left({\displaystyle \frac{1+\alpha ^2}{2}}\right)\mathrm{\Gamma }({\displaystyle \frac{1\alpha ^2}{2}},x^2)\right]\right\},`$ (118) where $`\mathrm{\Gamma }(s,z)`$ denotes the incomplete $`\mathrm{\Gamma }`$–function. One easily shows $`h(\alpha ^2;x=0)=\mathrm{\Gamma }(\alpha ^2/2)/2`$ and the asymptotic behavior for $`\|x\|1`$ $$h(\alpha ^2;x)=\frac{\mathrm{\Gamma }\left(\frac{1+\alpha ^2}{2}\right)}{2x}+O(x^2)$$ (119) with a smooth crossover in between. It will be convenient to use the normalized operators $`S_j(\alpha ;k)`$ $`\stackrel{\mathrm{def}}{=}`$ $`\left[{\displaystyle \frac{2\pi }{L}}\mathrm{\Gamma }(\alpha ^2)\left({\displaystyle \frac{L\|k\|}{2\pi }}\right)^{1\alpha ^2}\right]^{1/2}V_j(\alpha ;k)`$ (120) $`S_j^{}(\alpha ;k)`$ $`=`$ $`\left[{\displaystyle \frac{2\pi }{L}}\mathrm{\Gamma }(\alpha ^2)\left({\displaystyle \frac{L\|k\|}{2\pi }}\right)^{1\alpha ^2}\right]^{1/2}V_j(\alpha ;k)`$ (121) with the properties (see (276)) $$S_1^{}(\alpha ;k)\|\mathrm{\Omega }=S_1^{}(\alpha ;k)\|\mathrm{\Omega }=S_2^{}(\alpha ;k)\|\mathrm{\Omega }=S_2^{}(\alpha ;k)\|\mathrm{\Omega }=0k>0$$ (122) and the normalization $`S_1^{}(\alpha ;k)S_1^{}(\alpha ;k^{})=S_2^{}(\alpha ;k)S_2^{}(\alpha ;k^{})`$ $`=`$ $`\delta _{kk^{}}\theta (k)L/2\pi `$ (123) $`S_1^{}(\alpha ;k)S_1^{}(\alpha ;k^{})=S_2^{}(\alpha ;k)S_2^{}(\alpha ;k^{})`$ $`=`$ $`\delta _{kk^{}}\theta (k)L/2\pi `$ (124) for $`\|ak\|,\|ak^{}\|1`$ as follows easily from (290) and (291). We express (116) in terms of these operators and find $`{\displaystyle \frac{2\stackrel{~}{u}^2(B)}{a^3\mathrm{\Gamma }^2(\alpha ^2)}}`$ $`\times `$ $`(\mathrm{cos}(\pi \alpha ^2){\displaystyle \underset{k>0}{}}(ak)^{2\alpha ^21}v^2(B;k)S_2^{}(\alpha ;k)S_2^{}(\alpha ;k)`$ (127) $`+{\displaystyle \frac{1}{\pi }}\mathrm{sin}(\pi \alpha ^2){\displaystyle \underset{k>0}{}}(ak)^{\alpha ^21}(8B/a^2)^{\alpha ^2/2}h(\alpha ^2;\sqrt{8B}k)S_2^{}(\alpha ;k)S_2^{}(\alpha ;k)`$ $`+{\displaystyle \frac{1}{\pi }}\mathrm{sin}(\pi \alpha ^2){\displaystyle \underset{k>0}{}}(ak)^{\alpha ^21}(8B/a^2)^{\alpha ^2/2}h(\alpha ^2;\sqrt{8B}k)S_2^{}(\alpha ;k)S_2^{}(\alpha ;k))`$ In the first two terms we do not need the subtraction operation $``$ anymore since the vacuum is already annihilated by them. The third term does not yet annihilate the vacuum. This can be easily achieved by using (315). However, already the second term will turn out to have hardly any effect, and the third term is again smaller than the second term. In order to simplify our expressions we therefore omit the third term in the sequel, although there would be no problem at all in carrying it along as well. Let us now also collect the other terms from (101) leading to $`[\eta ,H_{\mathrm{int}}]`$ $``$ $`{\displaystyle \frac{4\stackrel{~}{u}^2(B)}{a^3\mathrm{\Gamma }^2(\alpha ^2)}}{\displaystyle \underset{k>0}{}}\left(\mathrm{cos}(\pi \alpha ^2)(ak)^{2\alpha ^21}v^2(B;k)+{\displaystyle \frac{1}{\pi }}\mathrm{sin}(\pi \alpha ^2)(ak)^{\alpha ^21}(8B/a^2)^{\alpha ^2/2}h(\alpha ^2;\sqrt{8B}k)\right)`$ (129) $`\times \left(S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)+S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)+S_2^{}(\alpha ;k)S_2^{}(\alpha ;k)+S_2^{}(\alpha ;k)S_2^{}(\alpha ;k)\right).`$ Since $`\alpha `$ generically flows as a function of $`B`$, this implies that vertex operators with different scaling dimensions contribute to each wavevector $`k`$. However, most of the contribution to a given $`k`$–vector occurs in a narrow range of the flow: We will see in Sect. IV.C that for a given $`k`$ the main contribution occurs when $`BB_k`$ with $`B_k`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \frac{1}{4k^2}}\text{(weak–coupling phase)}`$ (130) $`B_k`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \frac{1}{4k\sqrt{k^2+M^2}}}\text{(strong–coupling phase with mass gap }M\text{)}.`$ (131) In order to simplify our notation we therefore use a single scaling dimension $`\alpha `$ corresponding to each $`k`$ with $`\alpha =\alpha (B_k)`$.<sup>§</sup><sup>§</sup>§ This approximation becomes exact in the low–energy limit, e.g. in the strong–coupling phase for $`\|k\|M`$. The newly generated term in the Hamiltonian can then be written as $$H_{\mathrm{diag}}(B)=\underset{k>0}{}\omega (B;k)\left(P_1^{}(k)P_1^{}(k)+P_1^{}(k)P_1^{}(k)+P_2^{}(k)P_2^{}(k)+P_2^{}(k)P_2^{}(k)\right)$$ (132) with $$P_j^{}(k)\stackrel{\mathrm{def}}{=}S_j^{}(\alpha (B_k);k),P_j^{}(k)=S_j^{}(\alpha (B_k);k)$$ (133) and according to (129) $$\frac{\omega (B;k)}{B}=\frac{4\stackrel{~}{u}^2(B)}{a^3\mathrm{\Gamma }^2(\alpha ^2)}\left(\mathrm{cos}(\pi \alpha ^2)(ak)^{2\alpha ^21}v^2(B;k)+\frac{1}{\pi }\mathrm{sin}(\pi \alpha ^2)(ak)^{\alpha ^21}(8B/a^2)^{\alpha ^2/2}h(\alpha ^2;\sqrt{8B}k)\right)$$ (134) with $`\omega (B=0;k)=0`$. Using Eqs. (47) one can easily check $`[H_0,H_{\mathrm{diag}}(B)]=0`$, therefore $`H_{\mathrm{diag}}(B)`$ can be interpreted as diagonal. #### 5 $`H_{\mathrm{res}}`$ So far we have not discussed the commutators with the structure $`[V_1^{}V_2^{},V_1^{}V_2^{}]`$ and $`[V_1^{}V_2^{},V_1^{}V_2^{}]`$ that are also generated by $`[\eta ,H_{\mathrm{int}}]`$. From (60) one concludes that these commutators contain only the $`R`$–term. The operator product expansion in the $`R`$–term then generates interactions with the structure $`V_1(2\alpha )V_2(2\alpha )`$ etc., that is only terms with larger scaling dimensions. These interactions are neglected in the present approximation, just like the higher–order terms in (76). We formally sum up these neglected terms with larger scaling dimensions in $`H_{\mathrm{res}}`$. ### E Commutator $`[\eta ,H_{\mathrm{diag}}]`$ We also have to study the effect of the infinitesimal unitary transformations on the newly generated terms (132). An overlap exists essentially only for wavevectors of order $`B^{1/2}`$. For notational simplicity we can therefore use the running scaling dimension $`\alpha =\alpha (B)`$ in (132) and arrive at the follwing commutator $`[\eta ,H_{\mathrm{diag}}]`$ $`=`$ $`2i{\displaystyle 𝑑x𝑑y\frac{u(B;y)}{y}\underset{k>0}{}\omega (B;k)}`$ (137) $`\times ([V_1(\alpha ;x)V_2(\alpha ;xy),S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)+S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)`$ $`+S_2^{}(\alpha ;k)S_2^{}(\alpha ;k)+S_2^{}(\alpha ;k)S_2^{}(\alpha ;k)]+\mathrm{h}.\mathrm{c}.).`$ A typical contribution comes e.g. from $$[V_1(\alpha ;x),S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)]=V_1(\alpha ;x)S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)V_1(\alpha ;x)S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)V_1(\alpha ;x).$$ (138) The first and the second term on the rhs lead to normal–ordered interactions with larger scaling dimensions and are therefore neglected (or formally contained in $`H_{\mathrm{res}}`$). The third term on the rhs gives rise to interactions of the type $`H_{\mathrm{int}}`$ leading to $`[\eta ,H_{\mathrm{diag}}]`$ $``$ $`2i{\displaystyle 𝑑x𝑑y\frac{u(B;y)}{y}\underset{k>0}{}\omega (B;k)}`$ (142) $`\times (V_1(\alpha ;x)S_2^{}(\alpha ;k)V_2(\alpha ;xy)S_2^{}(\alpha ;k)V_1(\alpha ;x)S_2^{}(\alpha ;k)S_2^{}(\alpha ;k)V_2(\alpha ;xy)`$ $`+V_1(\alpha ;x)S_1^{}(\alpha ;k)S_1^{}(\alpha ;k)V_2(\alpha ;xy)S_1^{}(\alpha ;k)V_1(\alpha ;x)S_1^{}(\alpha ;k)V_2(\alpha ;xy)`$ $`+\mathrm{h}.\mathrm{c}.)`$ $`=`$ $`2i{\displaystyle 𝑑x𝑑y\frac{u(B;y)}{y}\underset{k>0}{}\omega (B;k)\left[\mathrm{\Gamma }(\alpha ^2)\frac{2\pi }{L}\left(\frac{L\|k\|}{2\pi }\right)^{1\alpha ^2}\right]^{1/2}}`$ (145) $`\times (V_1(\alpha ;x)S_2^{}(\alpha ;k)e^{ik(xy)}V_1(\alpha ;x)S_2^{}(\alpha ;k)e^{ik(xy)}`$ $`+S_1^{}(\alpha ;k)V_2(\alpha ;xy)e^{ikx}S_1^{}(\alpha ;k)V_2(\alpha ;xy)e^{ikx}+\mathrm{h}.\mathrm{c}.)`$ $`=`$ $`4\pi {\displaystyle 𝑑x\underset{k>0}{}\omega (B;k)ku(B;k)\left[\mathrm{\Gamma }(\alpha ^2)\frac{2\pi }{L}\left(\frac{L\|k\|}{2\pi }\right)^{1\alpha ^2}\right]^{1/2}}`$ (148) $`\times (V_1(\alpha ;x)S_2^{}(\alpha ;k)e^{ikx}+V_1(\alpha ;x)S_2^{}(\alpha ;k)e^{ikx}`$ $`+S_1^{}(\alpha ;k)V_2(\alpha ;x)e^{ikx}+S_1^{}(\alpha ;k)V_2(\alpha ;x)e^{ikx}+\mathrm{h}.\mathrm{c}.)`$ $`=`$ $`8\pi L{\displaystyle \underset{k}{}}\omega (B;\|k\|)\|k\|u(B;k){\displaystyle \frac{1}{\mathrm{\Gamma }(\alpha ^2)}}\left({\displaystyle \frac{L\|k\|}{2\pi }}\right)^{\alpha ^21}\left(S_1^{}(\alpha ;k)S_2^{}(\alpha ;k)+S_2^{}(\alpha ;k)S_1^{}(\alpha ;k)\right).`$ (149) This term has to be compared with (42) $`H_{\mathrm{int}}`$ $`=`$ $`{\displaystyle }dxdyu(B;y)(V_1(\alpha ;x)V_2(\alpha ;xy)+\mathrm{h}.\mathrm{c}.)`$ (150) $`=`$ $`2\pi L{\displaystyle \underset{k}{}}u(B;k){\displaystyle \frac{1}{\mathrm{\Gamma }(\alpha ^2)}}\left({\displaystyle \frac{L\|k\|}{2\pi }}\right)^{\alpha ^21}\left(S_1^{}(\alpha ;k)S_2^{}(\alpha ;k)+S_2^{}(\alpha ;k)S_1^{}(\alpha ;k)\right).`$ (151) Eq. (149) therefore gives an additional contribution to the previous flow equation (55) and one finds $$\frac{v(B;k)}{B}=4k^2v(B;k)4\|k\|\omega (B;\|k\|)v(B;k).$$ (152) ### F Flow of $`\beta ^2`$ #### 1 Unitary transformation $`e^U`$ We have seen that the term (98) $$\underset{k}{}w_k(B)\|k\|\sigma _1(k)\sigma _2(k)dB$$ (153) with $$w_k(B)=\frac{32}{a^2}\left(\frac{32B}{a^2}\right)^{1\alpha ^2(B)}\frac{\alpha ^2(B)}{2\mathrm{\Gamma }(\alpha ^2(B)1)}\stackrel{~}{u}^2(B)f(\alpha ^2(B);k\sqrt{B})$$ (154) is generated during the flow. This term is not contained in the original sine–Gordon Hamiltonian. We will now show how this term can be eliminated by an additional unitary transformation $`e^U`$ with $$U=\underset{p>0}{}\psi _p\left(\sigma _1(p)\sigma _2(p)\sigma _1(p)\sigma _2(p)\right)$$ (155) with suitable parameters $`\psi _p`$ . Let us first write down some general properties of this unitary transformation for general $`\psi _p`$. The bosonic fields are transformed according to $`e^U\sigma _1(p)e^U`$ $`=`$ $`\sigma _1(p)\mathrm{cosh}\psi _p+\sigma _2(p)\mathrm{sinh}\psi _p`$ (156) $`e^U\sigma _2(p)e^U`$ $`=`$ $`\sigma _2(p)\mathrm{cosh}\psi _p+\sigma _1(p)\mathrm{sinh}\psi _p`$ (157) and vertex operators as $`e^UV_1(\alpha ;x)e^U`$ $`=`$ $`e^C:\mathrm{exp}\left[\alpha {\displaystyle \underset{p0}{}}\mathrm{cosh}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)\right]::\mathrm{exp}\left[\alpha {\displaystyle \underset{p0}{}}\mathrm{sinh}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _2(p)\right]:`$ (158) $`e^UV_2(\alpha ;x)e^U`$ $`=`$ $`e^C:\mathrm{exp}[\alpha {\displaystyle \underset{p0}{}}\mathrm{cosh}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _2(p)]::\mathrm{exp}[\alpha {\displaystyle \underset{p0}{}}\mathrm{sinh}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)]:,`$ (159) where $$C=\alpha ^2\underset{p>0}{}\mathrm{sinh}\psi _p\mathrm{sinh}\psi _p\frac{e^{a\|p\|}}{p}.$$ (160) #### 2 $`e^UHe^U`$ We take the point of view that (153) has been generated infinitesimally by integrating the flow equations from $`B`$ to $`B+dB`$. We therefore apply the above unitary transformation (155) to $`H`$ $$H(B+dB)e^UH(B+dB)e^U.$$ (161) Let us first investigate the effect on $`H_0`$: For the choice $$\psi _p=\frac{w_p(B)}{2}dB$$ (162) the transformation $`e^UH_0e^U`$ reproduces $`H_0`$ and generates an additional term that annihilates (153). This can be shown easily by using the transformation rules (157). Notice that terms of order $`\psi _p^2`$ and higher can be neglected since $`\psi _p`$ is of order $`dB`$. Next we have to find the effect of this transformation on $`H_{\mathrm{int}}(B)`$. Using (158) one finds $`e^UV_1(\alpha ;x)V_2(\alpha ;y)e^U`$ $`=`$ $`V_1(\alpha ;x):\mathrm{exp}\left[\alpha {\displaystyle \underset{p0}{}}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipy}\sigma _1(p)\right]:`$ (164) $`\times V_2(\alpha ;y):\mathrm{exp}\left[\alpha {\displaystyle \underset{p0}{}}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _2(p)\right]:,`$ where we have again used that $`\psi _p`$ is of order $`dB`$. Using an OPE, we can combine the first two terms into a vertex operator with a modified scaling dimension, and likewise for the second two terms. The calculation proceeds along similar lines as in (296): $`V_1(\alpha ;x):\mathrm{exp}\left[\alpha {\displaystyle \underset{p0}{}}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipy}\sigma _1(p)\right]:`$ (165) $`=`$ $`\mathrm{exp}\left[\alpha {\displaystyle \underset{p>0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)\right]\mathrm{exp}\left[\alpha {\displaystyle \underset{p<0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)\right]`$ (167) $`\times \mathrm{exp}\left[\alpha {\displaystyle \underset{p>0}{}}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipy}\sigma _1(p)\right]\mathrm{exp}\left[\alpha {\displaystyle \underset{p<0}{}}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipy}\sigma _1(p)\right]`$ $`=`$ $`e^C\times \mathrm{exp}\left[\alpha {\displaystyle \underset{p>0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}(1+\psi _pe^{ip(yx)})\sigma _1(p)\right]`$ (169) $`\times \mathrm{exp}\left[\alpha {\displaystyle \underset{p<0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}(1+\psi _pe^{ip(yx)})\sigma _1(p)\right]`$ where $$C=\alpha ^2\underset{p>0}{}\psi _p\frac{\mathrm{exp}(apip(xy))}{p}.$$ (170) From (97) we know that $`\psi _p`$ falls of rapidly on an energy scale of order $`B^{1/2}`$. The leading behavior of the sum is therefore (except for an uninteresting initial transient) $$C=\alpha ^2\psi _{p=0}\underset{p>0}{}\frac{\mathrm{exp}(\stackrel{~}{a}pip(xy))}{p}$$ (171) leading to (compare Eq. (286)) $$e^C=\left(\frac{2\pi (\stackrel{~}{a}+i(xy))}{L}\right)^{\alpha ^2\psi _{p=0}}.$$ (172) Here $`\stackrel{~}{a}=s\sqrt{B}`$ takes into account the UV–cutoff generated by the decay of the functions $`f(\alpha ^2;x)`$ from Fig. 2. The actual value of the proportionality constant $`s`$ will only affect our results in next to leading order as will be shown later. Still we give its value here for latter comparison with the three loop scaling results: For $`\alpha ^2=2`$ one finds in a somehow lengthy calculation $$\mathrm{ln}s=\mathrm{ln}2+\frac{\pi }{4}\frac{\gamma }{2},$$ (173) where $`\gamma 0.577`$ is Euler’s constant. $`s`$ depends only weakly on $`\alpha ^2`$, therefore we will use this value throughout. Next $`\mathrm{exp}\left[\alpha {\displaystyle \underset{p>0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}(1+\psi _pe^{ip(yx)})\sigma _1(p)\right]`$ (174) $`=`$ $`\mathrm{exp}\left[\alpha (1+\psi _{p=0}){\displaystyle \underset{p>0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)\right]`$ (176) $`\times \mathrm{exp}\left[i(xy)\alpha \psi _{p=0}{\displaystyle \underset{p>0}{}}\sqrt{\|p\|}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)+O((xy)^2)\right]`$ $`=`$ $`V_1(\alpha (1+\psi _{p=0});x)\times \left(1+i(xy)\alpha \psi _{p=0}{\displaystyle \underset{p>0}{}}\sqrt{\|p\|}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)+O((xy)^2)\right)`$ (177) $`=`$ $`V_1(\alpha (1+\psi _{p=0});x)+\text{more irrelevant terms}`$ (178) up to less singular (more irrelevant) terms. We formally include these more irrelevant terms (they have the structure of products of vertex operators multiplied by derivatives of the bosonic field) into $`H_{\mathrm{res}}`$ and neglect them from now on. It should also be noted that we have approximated the flow of the scaling dimension in the vertex operator by restricting ourselves to the flow in the IR–limit $`\alpha \alpha (1+\psi _{p=0})`$. At first sight this seems to be a problematic approximation since the functions $`f(\alpha ^2;x)`$ in (97) are nontrivial. However, the effect of the unitary transformations is only cumulative on the low–energy scale where $`f(\alpha ^2;x)\stackrel{x0}{=}1`$. Hence it is possible to restrict ourselves to the IR–limit. Putting everything together we find $`H_{\mathrm{int}}(B+dB)`$ $``$ $`e^UH_{\mathrm{int}}(B+dB)e^U`$ (179) $`=`$ $`\left({\displaystyle \frac{2\pi s\sqrt{B}}{L}}\right)^{\alpha ^2(B)w_{p=0}(B)dB}{\displaystyle 𝑑x𝑑yu(B;y)}`$ (181) $`\times [V_1(\alpha (1{\displaystyle \frac{w_{p=0}(B)}{2}}dB);x)V_2(\alpha (1{\displaystyle \frac{w_{p=0}(B)}{2}}dB;xy)+\mathrm{h}.\mathrm{c}.]`$ Similar to the OPE above we have set $`xy=0`$ in (172), that is we have neglected less singular terms. Finally, we have to investigate the effect of the additional unitary transformation (161) on $`H_{\mathrm{diag}}`$ from (132). This is similar to the action on the interaction term, except that here we find terms of the structure $`e^UV_1(\alpha ;x)V_1(\alpha ;y)e^U`$ $`=`$ $`V_1(\alpha ;x)V_1(\alpha ;y):\mathrm{exp}\left[\alpha {\displaystyle \underset{p0}{}}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _2(p)\right]:`$ (183) $`\times :\mathrm{exp}[\alpha {\displaystyle \underset{p0}{}}\psi _p{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipy}\sigma _2(p)]:`$ and likewise for $`j=2`$. The third and fourth term can be combined using an OPE (297) and one easily notices that the only surviving term is a constant 1 since all other terms are of order $`\psi _p^2=O(dB^2)`$. Hence $$H_{\mathrm{diag}}(B+dB)e^UH_{\mathrm{diag}}(B+dB)e^U=H_{\mathrm{diag}}(B+dB).$$ (184) Summing up, we have looked at another infinitesimal unitary transformation (161) that acts on the Hamiltonian during the flow equation procedure in addition to the generator (49). An alternative viewpoint is to say that the full generator of the flow now takes the structure $`\eta _{\mathrm{new}}(B)`$ $`=`$ $`2i{\displaystyle }dxdy{\displaystyle \frac{u(B;y)}{y}}(V_1(\alpha ;x)V_2(\alpha ;xy)+\mathrm{h}.\mathrm{c}.)`$ (186) $`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{p}{}}w_p(B)\left(\sigma _1(p)\sigma _2(p)\sigma _1(p)\sigma _2(p)\right).`$ #### 3 Flow equation for $`\beta ^2`$ From (179) we can read of that the additional infinitesimal unitary transformation generates a flow in the scaling dimension of the vertex operators $`\alpha (B)`$ $``$ $`\alpha (B)\left(1{\displaystyle \frac{w_{p=0}(B)}{2}}dB\right)`$ (187) $`{\displaystyle \frac{d\alpha ^2}{dB}}`$ $`=`$ $`w_{p=0}(B)\alpha ^2(B)`$ (188) $`=`$ $`{\displaystyle \frac{32}{a^2}}\left({\displaystyle \frac{32B}{a^2}}\right)^{1\alpha ^2(B)}{\displaystyle \frac{\alpha ^4(B)}{2\mathrm{\Gamma }(\alpha ^2(B)1)}}\stackrel{~}{u}^2(B).`$ (189) We can already see that $`\alpha ^2=1`$ is a fixed point of the flow equation approach due to the diverging $`\mathrm{\Gamma }`$–function in the denominator of (188). This will be one of the key results of our new approach. According to (179) this flow of the scaling dimension now induces a flow of the coupling constant $`u(B;y)`$ $`u(B;y)`$ $``$ $`u(B;y)\left({\displaystyle \frac{2\pi s\sqrt{B}}{L}}\right)^{\alpha ^2(B)w_{p=0}(B)dB}`$ (190) $`=`$ $`u(B;y)\left(1\alpha ^2(B)w_{p=0}(B)dB\mathrm{ln}\left({\displaystyle \frac{2\pi s\sqrt{B}}{L}}\right)\right)`$ (191) $`{\displaystyle \frac{du(B;y)}{dB}}`$ $`=`$ $`u(B;y){\displaystyle \frac{d\alpha ^2}{dB}}\mathrm{ln}\left({\displaystyle \frac{2\pi s\sqrt{B}}{L}}\right).`$ (192) The solution is straightforward $`u(B;y)`$ $`=`$ $`u(B=0;y)\mathrm{exp}\left({\displaystyle _0^B}𝑑B^{}{\displaystyle \frac{d\alpha ^2}{dB^{}}}\mathrm{ln}\left({\displaystyle \frac{2\pi s\sqrt{B^{}}}{L}}\right)\right)`$ (193) $`=`$ $`u(B=0;y)\mathrm{exp}\left({\displaystyle _0^B}𝑑B^{}{\displaystyle \frac{d\alpha ^2}{dB^{}}}\mathrm{ln}\left({\displaystyle \frac{2\pi a}{L}}\sqrt{{\displaystyle \frac{32B^{}}{a^2}}}{\displaystyle \frac{s}{\sqrt{32}}}\right)\right)`$ (194) $`=`$ $`u(B=0;y)\left({\displaystyle \frac{2\pi a}{L}}\right)^{\alpha ^2(B)\alpha ^2(0)}\left({\displaystyle \frac{s}{\sqrt{32}}}\right)^{\alpha ^2(B)\alpha ^2(0)}\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle _0^B}𝑑B^{}{\displaystyle \frac{d\alpha ^2}{dB^{}}}\mathrm{ln}\left({\displaystyle \frac{32B^{}}{a^2}}\right)\right).`$ (195) Using the parametrization (54) $$u(B;p)=\frac{\stackrel{~}{u}(B)}{4\pi ^2a^2}\left(\frac{2\pi a}{L}\right)^{\alpha ^2(B)}v(B;p)$$ (196) we see that this can be most conveniently expressed as a flow equation for the running coupling constant $`\stackrel{~}{u}(B)`$ in (196) $$\stackrel{~}{u}(B)=u_0\left(\frac{s}{\sqrt{32}}\right)^{\alpha ^2(B)\alpha ^2(0)}\mathrm{exp}\left(\frac{1}{2}_0^B𝑑B^{}\frac{d\alpha ^2}{dB^{}}\mathrm{ln}\left(\frac{32B^{}}{a^2}\right)\right).$$ (197) Introducing the dimensionless logarithmic flow parameter $`\mathrm{}`$ $$\mathrm{}\stackrel{\mathrm{def}}{=}\frac{1}{2}\mathrm{ln}\left(\frac{32B}{a^2}\right),$$ (198) one can show by partial integration $$\frac{1}{2}_0^B𝑑B^{}\frac{d\alpha ^2}{dB^{}}\mathrm{ln}\left(\frac{32B^{}}{a^2}\right)=_0^{\mathrm{}}𝑑\mathrm{}^{}\alpha ^2(\mathrm{}^{})+\alpha ^2(\mathrm{})\mathrm{}.$$ (199) Using this we can sum up the results of this section in the following two equations: $`\stackrel{~}{u}(\mathrm{})`$ $`=`$ $`u_0\left({\displaystyle \frac{s}{\sqrt{32}}}\right)^{\alpha ^2(\mathrm{})\alpha ^2(0)}\mathrm{exp}\left({\displaystyle _0^{\mathrm{}}}𝑑\mathrm{}^{}\alpha ^2(\mathrm{}^{})+\alpha ^2(\mathrm{})\mathrm{}\right)`$ (200) $`{\displaystyle \frac{d\alpha ^2}{d\mathrm{}}}`$ $`=`$ $`u_0^2\left({\displaystyle \frac{s^2}{32}}\right)^{\alpha ^2(\mathrm{})\alpha ^2(0)}{\displaystyle \frac{\alpha ^4(\mathrm{})}{\mathrm{\Gamma }(\alpha ^2(\mathrm{})1)}}\mathrm{exp}\left(4\mathrm{}2{\displaystyle _0^{\mathrm{}}}𝑑\mathrm{}^{}\alpha ^2(\mathrm{}^{})\right).`$ (201) These two equations constitute the key results of this work. Eq. (201) describes the flow of the scaling dimension under the flow equation procedure, and from Eq. (200) it follows how this induces the flow of the running coupling constant $`\stackrel{~}{u}(\mathrm{})`$. Therefore these equations will serve as a generalization of the scaling equations derived in perturbative renormalization theory in Sect. II.B. The value of the constant $`s`$ in these equations will turn out to affect our results only in next to leading order. ## IV Solution of the flow equations ### A Summary of the flow equations In this section we will sum up the results for the flow of the sine–Gordon Hamiltonian under the effect of the infinitesimal unitary transformation (186) as derived above. For general $`B`$ the sequence $`H(B)`$ of unitarily equivalent Hamiltonians takes the form $$H(B)=H_0+H_{\mathrm{int}}(B)+H_{\mathrm{diag}}(B)+H_{\mathrm{res}}(B).$$ (202) Here $`H_0`$ $`=`$ $`{\displaystyle 𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2\right)}`$ (203) $`H_{\mathrm{int}}(B)`$ $`=`$ $`{\displaystyle }dxdyu(B;y)(V_1(\alpha (B);x)V_2(\alpha (B);xy)+\mathrm{h}.\mathrm{c}.)`$ (204) $`H_{\mathrm{diag}}(B)`$ $`=`$ $`{\displaystyle \underset{k>0}{}}\omega (B;k)\left(P_1^{}(k)P_1^{}(k)+P_1^{}(k)P_1^{}(k)+P_2^{}(k)P_2^{}(k)+P_2^{}(k)P_2^{}(k)\right),`$ (205) with $`P_j(k)`$ given by (133) $$P_j(k)=\left[\frac{\mathrm{\Gamma }(\alpha ^2(B_k))}{2\pi L}\left(\frac{L\|k\|}{2\pi }\right)^{1\alpha ^2(B_k)}\right]^{1/2}𝑑xe^{ikx}V_j(\alpha (B_k);x)$$ (206) and $`B_k`$ from (131). The operators $`P_j^{}(k),P_j^{}(k)`$ will turn out to be the soliton creation and annihilation operators. They are normalized according to (124). Notice $`H_{\mathrm{diag}}(B)\|\mathrm{\Omega }=0`$. $`H_{\mathrm{res}}`$ contains the neglected terms and will from now on be omitted in our analysis. At any given $`B`$–scale these neglected terms have a larger scaling dimension than the interaction term $`H_{\mathrm{int}}(B)`$: They are more irrelevant by at least two spatial derivatives. Notice that $`H_{\mathrm{res}}`$ vanishes for $`\beta ^2=4\pi `$ since then no approximations are made. Summing up the differential flow equations for the parameters in $`H(B)`$ we have $`{\displaystyle \frac{v(B;k)}{B}}`$ $`=`$ $`4k^2v(B;k)4k\omega (B;k)v(B;k)`$ (207) $`{\displaystyle \frac{(a\omega (B;k))}{(B/a^2)}}`$ $`=`$ $`{\displaystyle \frac{4\stackrel{~}{u}^2(B)}{\mathrm{\Gamma }^2(\alpha ^2)}}\left(\mathrm{cos}(\pi \alpha ^2)(ak)^{2\alpha ^21}v^2(B;k)+{\displaystyle \frac{1}{\pi }}\mathrm{sin}(\pi \alpha ^2)(ak)^{\alpha ^21}(8B/a^2)^{\alpha ^2/2}h(\alpha ^2;\sqrt{8B}k)\right)`$ (208) for $`k>0`$, and $`v(B;k)=v(B;k)`$ symmetric in $`k`$. For notational convenience we have written $`\alpha (B)`$ without its argument $`B`$ in these equations. $`h(\alpha ^2;x)`$ has been defined in (118) and $$u(B;y)=\frac{\stackrel{~}{u}(B)}{4\pi ^2a^2}\left(\frac{2\pi a}{L}\right)^{\alpha ^2(B)}\underset{p}{}v(B;p)e^{ipy}.$$ (209) The initial conditions are $`H_{\mathrm{res}}(B=0)=0`$, $`\omega (B=0;k)=0`$ and $`v(B=0;k)=1`$, whereupon (202) takes the form (19) of our original sine–Gordon Hamiltonian. For $`\alpha _0^2<1`$ one should in fact be more cautious and start the integration only at $`B=a^2`$: The above equations hold only for $`Ba^2`$ since we have used $`\|ak\|1`$ in our calculation. One can easily verify that there is no flow of the parameters for $`Ba^2`$, and the simplest way to take this into account is to pose the initial conditions at $`B=a^2`$. As we will see later the flow of the parameters is such that $`v(B=\mathrm{};k)=0`$ for $`\beta _0^2>2\pi `$. In this parameter region the Hamiltonian $`H(B)`$ from (202) therefore becomes diagonal with $`H_{\mathrm{int}}(B=\mathrm{})=0`$ in the limit $`B\mathrm{}`$ as expected under the flow equation procedure $$H(B=\mathrm{})=H_0+H_{\mathrm{diag}}(B=\mathrm{}),$$ (210) notice $`[H_0,H_{\mathrm{diag}}(B)]=0`$. The equations (207) and (208) have to supplemented with the differential equation governing the flow of the scaling dimension (201) and the thereby induced flow of the running coupling constant (200). It is possible to rewrite (201) in a more conventional form for comparison with the perturbative scaling approach. Introducing a new function $$u(\mathrm{})\stackrel{\mathrm{def}}{=}u_0\left(\frac{s}{\sqrt{32}}\right)^{\alpha ^2(\mathrm{})\alpha ^2(0)}\mathrm{exp}\left(2\mathrm{}_0^{\mathrm{}}𝑑\mathrm{}^{}\alpha ^2(\mathrm{}^{})\right)$$ (211) we can rewrite (201) as a set of two coupled differential equations $`{\displaystyle \frac{d\alpha ^2}{d\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{\alpha ^4(\mathrm{})}{\mathrm{\Gamma }(\alpha ^2(\mathrm{})1)}}u^2(\mathrm{})`$ (212) $`{\displaystyle \frac{du}{d\mathrm{}}}`$ $`=`$ $`\left(2\alpha ^2(\mathrm{})\right)u(\mathrm{})+\left({\displaystyle \frac{\pi }{4}}{\displaystyle \frac{\gamma }{2}}{\displaystyle \frac{1}{2}}\mathrm{ln}8\right){\displaystyle \frac{d\alpha ^2}{d\mathrm{}}}u(\mathrm{})`$ (213) with the initial conditions $`u(\mathrm{}=0)=u_0`$ and $`\alpha (\mathrm{}=0)=\beta _0/\sqrt{4\pi }`$. For convenience we have used the dimensionless flow parameter $`\mathrm{}=\frac{1}{2}\mathrm{ln}(32B/a^2)`$ from (198) in these equations. Notice that the second term in (213) is of order $`u^3`$ and does therefore not contribute to the leading behavior for small $`u_0`$. Since our present flow equation expansion has not taken all the terms in order $`u^3`$ into account anyway, we can omit this term and arrive at $`{\displaystyle \frac{d\alpha ^2}{d\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{\alpha ^4(\mathrm{})}{\mathrm{\Gamma }(\alpha ^2(\mathrm{})1)}}u^2(\mathrm{})`$ (214) $`{\displaystyle \frac{du}{d\mathrm{}}}`$ $`=`$ $`\left(2\alpha ^2(\mathrm{})\right)u(\mathrm{}).`$ (215) The running coupling constant $`\stackrel{~}{u}(\mathrm{})`$ from (200) can also be expressed as $$\stackrel{~}{u}(\mathrm{})=u(\mathrm{})\mathrm{exp}\left((\alpha ^2(\mathrm{})2)\mathrm{}\right).$$ (216) In terms of the sine–Gordon parameter $`\beta `$ these equations take the equivalent form $`{\displaystyle \frac{d\beta ^2(\mathrm{})}{d\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi \mathrm{\Gamma }\left(1+\beta ^2(\mathrm{})/4\pi \right)}}u^2(\mathrm{})`$ (217) $`{\displaystyle \frac{du}{d\mathrm{}}}`$ $`=`$ $`\left(2{\displaystyle \frac{\beta ^2(\mathrm{})}{4\pi }}\right)u(\mathrm{})`$ (218) and $$\stackrel{~}{u}(\mathrm{})=u(\mathrm{})\mathrm{exp}\left(\left(\frac{\beta ^2(\mathrm{})}{4\pi }2\right)\mathrm{}\right).$$ (219) Finally it is of some interest to express $`H_{\mathrm{int}}(B)`$ directly in terms of the bosonic field $`\varphi (x)`$ and its dual $`\mathrm{\Theta }(x)`$ (see Eq. (12)). After a short calculation one finds $`H_{\mathrm{int}}(B)`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{u}(B)}{\pi a^2}}{\displaystyle 𝑑x𝑑yv(B;y)\mathrm{cos}\left(\beta (B)𝑑ϵc(ϵ)\frac{1}{2}\left(\varphi (x+ϵ)+\varphi (xy+ϵ)+\mathrm{\Theta }(x+ϵ)\mathrm{\Theta }(xy+ϵ)\right)\right)}.`$ (220) Since $`v(B;y)`$ becomes more and more nonlocal during the flow, one sees that the interaction term of the sine–Gordon model evolves from the original $`\mathrm{cos}(\beta \varphi (x))`$–structure to a nonlocal interaction term with the structure $$\mathrm{cos}\left((\beta /2)(\varphi (x)+\varphi (xy)+\mathrm{\Theta }(x)\mathrm{\Theta }(xy))\right).$$ (221) It is also possible to express $`H_{\mathrm{diag}}(B)`$ in terms of $`\varphi (x)`$ and $`\mathrm{\Theta }(x)`$, however, this expression does not lead to new insights. ### B Strong–coupling phase #### 1 Fixed points and phase structure We will now work on the explicit solution of the flow equations (207), (208) and (218). First we focus on the solution of (218), since knowledge of the flow of $`\beta (B)`$ and $`\stackrel{~}{u}(B)`$ is necessary for solving the system of equations (207) and (208) later on.Implicitly an assumption about the behavior of $`v(B;k)`$ has been made when deriving (218) in Sect. III.D. One can check that this assumption is self–consistently justified by analyzing the whole system of equations. From (218) one concludes that there are two possible kinds of asymptotic behavior: Either $`\beta ^2(\mathrm{})=4\pi `$ or $`\beta ^2(\mathrm{})8\pi `$. $`\beta ^2(\mathrm{})=4\pi `$ will turn out to be the attractive strong–coupling fixed point and values $`\beta ^2(\mathrm{})8\pi `$ correspond to the gapless weak–coupling phase. These flows are depicted in Fig. 3. Notice that the fundamental difference from the perturbative scaling equations (23) is the $`\mathrm{\Gamma }`$–function in the denominator of our flow equation for $`\beta ^2(\mathrm{})`$. Therefore $`\beta ^2=4\pi `$ is a fixed point in our approach which will be the main difference as compared to perturbative RG. This does not come as a surprise since the flow of $`\beta ^2`$ followed from higher order terms in the commutator $`[\eta (B),H_{\mathrm{int}}(B)]`$ in Sect. III.D. However, for $`\beta ^2=4\pi `$ our interaction term is quadratic if considered as an interaction term for Thirring fermions (see Sect. II.C), and naturally no higher order terms can be generated in $`[\eta (B),H_{\mathrm{int}}(B)]`$. Another way of saying this is that the flow in $`\beta ^2`$ is due to approximations in the flow equation scheme when higher order terms are generated. No flow of $`\beta ^2`$ can occur if the flow equations close exactly. In both phases $`\stackrel{~}{u}(\mathrm{})`$ remains finite for $`\mathrm{}\mathrm{}`$. On the other hand, one easily checks that $`u(\mathrm{})`$ diverges in the strong–coupling phase and vanishes asymptotically in the weak–coupling phase. Since $`u(\mathrm{})`$ has so far only been introduced for rewriting (201), its divergence in the strong–coupling phase need not worry us here. The question of the true expansion parameter of our approach will be discussed below in Sect. IV.D, and we will see that this expansion parameter is not $`u(\mathrm{})`$. In order to analyze the phase boundaries we can expand the $`\mathrm{\Gamma }`$–function in (218) around $`\beta ^2=8\pi `$. In leading order this reproduces the perturbative scaling equations (23) $`{\displaystyle \frac{d\beta ^2}{d\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{u^2}{4\pi }}`$ (222) $`{\displaystyle \frac{du}{d\mathrm{}}}`$ $`=`$ $`\left(2{\displaystyle \frac{\beta ^2}{4\pi }}\right)u.`$ (223) This approximation eventually breaks down in the strong–coupling phase as $`\beta ^2(\mathrm{})`$ flows to $`4\pi `$: Then (223) is not a good approximation for the true flow equation (218) anymore. Notice the sign difference from (23) because $`\mathrm{}`$ from (198) corresponds to $`\mathrm{ln}\mathrm{\Lambda }`$, therefore $`\mathrm{}`$ is integrated from 0 to $`\mathrm{}`$. Our flow equation approach therefore reproduces the conventional two–loop scaling equations if we expand around $`\beta ^2=8\pi `$. In this way we also reproduce the hidden SU(2)–symmetry of the sine–Gordon model for $`\beta _0^2=8\pi (1\pm u_0)`$ mentioned in Sect. II.C, although our approximation scheme does not manifestly respect this symmetry. Besides showing the consistency of our new approach with the conventional perturbative RG scheme at $`\beta ^28\pi `$, we also see immediately that our flow equations reproduce the Kosterlitz–Thouless phase diagram Fig. 1 of the sine–Gordon model established with RG: In the limit of small initial $`u_0`$ and $`\beta _0^2>8\pi (1+u_0+O(u_0^2))`$ we flow to a weak–coupling fixed point, for $`\beta _0^2<8\pi (1+u_0+O(u_0^2))`$ to the strong–coupling fixed point $`\beta ^2(\mathrm{})=4\pi `$. These two phases are again separated by a Kosterlitz–Thouless type transition along $`\beta _0^2=8\pi (1+u_0+O(u_0^2))`$. For the rest of this section we will focus on the strong–coupling phase. We will return to the weak–coupling phase in Sect. IV.C. #### 2 Low–energy effective Hamiltonian One advantage of the flow equation scheme is that we can easily analyze the behavior of our model with the final diagonal Hamiltonian $`H(B=\mathrm{})`$. However, a simpler kind of analysis is possible by identifying a low–energy effective Hamiltonian and analyzing this effective Hamiltonian. This will be done in this subsection. Our results will be confirmed by the analysis of $`H(B=\mathrm{})`$ in Sect. IV.C later on, but the identification of a low–energy effective Hamiltonian allows us to make contact with the conventional scaling picture, which is very useful for providing a simple coherent description of the flow equation approach in the strong–coupling phase. Let us look at $`H(B)`$ for large $`B`$ such that $`\|\beta ^2(B)4\pi \|1`$. Then we can approximately set $`\alpha (B)=1`$ in $`H_{\mathrm{int}}(B)`$ from (204) and rewrite (202) $$H(B)=H_{\mathrm{eff}}(B)+H_{\mathrm{diag}}(B),$$ (224) where $$H_{\mathrm{eff}}(B)=𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2\right)+\frac{\stackrel{~}{u}(B)}{a}\underset{k}{}v(B;k)\left(P_1^{}(k)P_2^{}(k)+P_2^{}(k)P_1^{}(k)\right).$$ (225) Now for $`\alpha =1`$ the creation and annihilation operators $`P_j^{}(k)`$, $`P_j^{}(k)`$ from (133) obey fermionic anticommutation relations (300) $$\{P_j^{}(k),P_j(k^{})\}=\delta _{kk^{}}\frac{L}{2\pi },\{P_j^{}(k),P_j^{}(k^{})\}=\{P_j(k),P_j(k^{})\}=0.$$ (226) One could also rewrite the kinetic term $`H_0`$ in terms of these fermions and would then arrive at a noninteracting Thirring model (28) with a nonlocal mass term as the low–energy effective Hamiltonian in the strong–coupling phase. However, we can also analyze the spectrum of our low–energy effective Hamiltonian $`H_{\mathrm{eff}}(B)`$ directly by working out the following commutators $`[H_{\mathrm{eff}}(B),P_1^{}(k)]`$ $`=`$ $`kP_1^{}(k)+{\displaystyle \frac{\stackrel{~}{u}(B)}{a}}v(B;k)P_2^{}(k)`$ (227) $`[H_{\mathrm{eff}}(B),P_2^{}(k)]`$ $`=`$ $`kP_2^{}(k)+{\displaystyle \frac{\stackrel{~}{u}(B)}{a}}v(B;k)P_1^{}(k)`$ (228) leading to the dispersion relation $$E_k=\sqrt{k^2+\left(\frac{\stackrel{~}{u}(B)}{a}v(B;k)\right)^2}.$$ (229) This dispersion relation describes the single–particle/hole excitation spectrum of the full Hamiltonian $`H(B)`$ for momenta $`\|k\|1/\sqrt{B}`$: According to (208) the terms in $`H_{\mathrm{diag}}(B)`$ corresponding to such momenta are only generated for even larger $`B`$, therefore we can neglect the effect of $`H_{\mathrm{diag}}(B)`$ for excitations with $`\|k\|1/\sqrt{B}`$. In this limit we also find the initial value $`v(B;k)=1`$ unchanged according to (207). Summing up, in the limit $`k0`$ the dispersion relation for single–particle/hole excitations in the strong–coupling phase has the form $`\pm E_k`$ with $$E_k=\sqrt{k^2+M^2}$$ (230) and the mass $$M=\frac{\stackrel{~}{u}(\mathrm{}=\mathrm{})}{a}.$$ (231) Eq. (230) is also the form expected from exact methods using integrability . Equation (231) is a key result in this work since it describes the relation between the running coupling and the generated mass term in the strong–coupling regime. We observe that the finiteness of the running coupling $`\stackrel{~}{u}(\mathrm{})`$ in (225) in the limit $`\mathrm{}\mathrm{}`$ is of fundamental importance in the flow equation scheme since $`\stackrel{~}{u}(\mathrm{})`$ sets the generated mass gap in the spectrum. Via (219) we can also establish the following expression for the “usual” running coupling $`u(\mathrm{})`$ in the language of the scaling equations (218). In the limit of $`\mathrm{}\mathrm{}`$ one finds $$u(\mathrm{})=aMe^{\mathrm{}}.$$ (232) Since $`\mathrm{\Lambda }_BB^{1/2}`$ plays the role of an effective UV–cutoff generated by the flow equations, this means that the dimensionless parameter $`u(\mathrm{\Lambda }_B)`$ diverges simply as the mass $`M`$ divided by the effective cutoff $`\mathrm{\Lambda }_B`$ $$u(B)\frac{M}{\mathrm{\Lambda }_B},$$ (233) which allows a simple physical picture of the diverging $`u(B)`$ in the strong–coupling phase. Again, this does not imply the breakdown of our approach since $`u(B)`$ is not the expansion parameter in our approach (see Sect. IV.D below). #### 3 Scaling behavior of the mass gap We will now analyze the behavior of the single–particle/hole mass $`M`$ in various regimes and compare this with results obtained with other methods. First we concentrate on the scaling behavior in the limit $`u_00`$ as this can be derived analytically: According to (232) it amounts to finding the scaling invariant $$I(\mathrm{})=e^{\mathrm{}}f(u(\mathrm{}),\beta ^2(\mathrm{}))$$ (234) with some suitable function $`f(u,\beta ^2)`$ such that $$\frac{dI(\mathrm{})}{d\mathrm{}}=0$$ (235) along the flow generated by (218). Like in the conventional scaling analysis the mass $`M`$ is then given by $$Mf(u_0,\beta _0^2)/a.$$ (236) One finds the same behavior as in two–loop order in Sect.II.B: For example for fixed $`\beta _0^2<8\pi `$ and $`u_00`$ one can easily check that $`f(u,\beta ^2)=u^{1/(2\beta ^2/4\pi )}`$ gives a scaling invariant up to terms in second order $$\frac{dI(\mathrm{})}{d\mathrm{}}=I(\mathrm{})\times O\left(\left(\frac{u(\mathrm{})}{2\beta ^2(\mathrm{})/4\pi }\right)^2\right)$$ (237) and therefore $$Mu_0^{1/(2\beta _0^2/4\pi )}/a$$ (238) in agreement with (24). Likewise one also finds the scaling behavior (25) and (26) since we could reproduce the perturbative RG–equations (223) in the vicinity of $`\beta ^2=8\pi `$ within our flow equation approach. Since we find agreement in two–loop order, it is of some interest to also compare with higher loop calculations . For simplicity we focus on the mass gap along $`S_{}`$ in Fig. 1, that is for $`\beta _0^2=8\pi (1u_0)`$ in the limit $`u_00`$. We write $`\beta ^2(\mathrm{})=8\pi (1v(\mathrm{}))`$ and expand (212) and (213) up to third order in $`u`$ and $`v`$ $`{\displaystyle \frac{dv}{d\mathrm{}}}`$ $`=`$ $`2u^24(1+\gamma )u^2v`$ (239) $`{\displaystyle \frac{du}{d\mathrm{}}}`$ $`=`$ $`2uv4\left({\displaystyle \frac{\pi }{4}}{\displaystyle \frac{\gamma }{2}}{\displaystyle \frac{1}{2}}\mathrm{ln}8\right)u^3`$ (240) where $`\gamma 0.577`$ is Euler’s constant. It is straightforward to derive the scaling invariant from these equations and one finds $$Mu_0^\tau \mathrm{exp}\left(\frac{1}{2u_0}\right)/a$$ (241) with $$\tau =\frac{1\mathrm{ln}8+\pi /2}{3}0.16.$$ (242) This should be compared with the three loop result (27 with the correct exponent $`\tau _{\mathrm{RG}}=1/2`$.The deviation from the result for $`\tau `$ in Ref. occurs because the term in order $`u^3`$ in (213) was neglected there. We see that the present order of our flow equation expansion is correct up to two loop order and deviates if compared with three loop RG. This is not surprising since we have not systematically taken all the terms in order $`\stackrel{~}{u}^3(\mathrm{})`$ into account in the present order of our flow equation scheme and there are contributions in order $`u^3`$ missing on the rhs of (240). In general no closed analytical solution for the set of differential equations (218) could be found. Some numerical solutions of (218) are depicted in Fig. 4. Of course the scaling behavior (238) is reproduced by these numerical solutions as can be seen in Fig. 4. Finally for $`\beta _0^2=4\pi `$ one can easily prove $`M=u_0/a`$ exactly using the above flow equations as should be expected since our scheme becomes exact in this case. ### C Properties of $`𝐇(𝐁=\mathrm{})`$ Since the flow equation procedure diagonalizes the sine–Gordon Hamiltonian, we can not only learn something about the mass gap in the spectrum, but analyze the entire dispersion relation throughout the crossover region. The final Hamiltonian takes the structure $$H(B=\mathrm{})=H_0+H_{\mathrm{diag}}(B=\mathrm{})$$ (243) with $`H_0`$ $`=`$ $`{\displaystyle 𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2\right)}`$ (244) $`H_{\mathrm{diag}}(B=\mathrm{})`$ $`=`$ $`{\displaystyle \underset{k>0}{}}\omega (B=\mathrm{};k)\left(P_1^{}(k)P_1^{}(k)+P_1^{}(k)P_1^{}(k)+P_2^{}(k)P_2^{}(k)+P_2^{}(k)P_2^{}(k)\right)`$ (245) and $`P_j(k)`$ defined in (133). Notice that $`[H_0,H_{\mathrm{diag}}(B=\mathrm{})]=0`$ and $`H(B=\mathrm{})\|\mathrm{\Omega }=0`$. Using (315) and (122) (see also the reasoning below (127)) and the normalization (124) one finds the following single–particle/hole excitation spectrum of $`H(B=\mathrm{})`$: * Soliton (particle) excitations with excitation energy $`E_k`$: + $`P_1^{}(k)\|\mathrm{\Omega }`$ for $`k>0`$ + $`P_2^{}(k)\|\mathrm{\Omega }`$ for $`k<0`$ * Antisoliton (hole) excitations with excitation energy $`E_k`$: + $`P_1^{}(k)\|\mathrm{\Omega }`$ for $`k<0`$ + $`P_2^{}(k)\|\mathrm{\Omega }`$ for $`k>0`$ The dispersion relation $`E_k`$ is given by $$E_k=\|k\|+\omega (B=\mathrm{};\|k\|).$$ (246) In order to find $`\omega (B=\mathrm{};k)`$ we next have to solve the system of equations (207) and (208). A closed analytical solution has not been possible except for the trivial case $`\beta _0^2=4\pi `$ where one reproduces (33) exactly. However, we will see below that the flow of $`\omega (B;k)`$ from its initial value 0 to $`\omega (B=\mathrm{};k)`$ occurs on the $`B`$–scale $`B_k`$ (131) and is negligible for $`BB_k`$ or $`BB_k`$. We can therefore to a good approximation replace $`\alpha (B)`$ and $`\stackrel{~}{u}(B)`$ in (208) by its values for $`B=B_k`$ and consider them as constants. Also one can verify numerically that the term in the differential equation proportional to $`\mathrm{sin}(\pi \alpha ^2)`$ changes the dispersion relation (246) only in relative order $`1\%`$ for $`\beta _0^24\pi `$ (for $`\|k/M\|<5`$ it e.g. affects $`\omega (B=\mathrm{};k)`$ less then $`2\%`$). In order to gain some first analytical insight we can neglect it. Notice that this approximation becomes exact in the low–energy limit in the strong–coupling phase since $`\alpha ^2(B)\stackrel{B\mathrm{}}{}1`$. We arrive at ($`k>0`$) $`{\displaystyle \frac{v(B;k)}{B}}`$ $`=`$ $`4k^2v(B;k)4k\omega (B;k)v(B;k)`$ (247) $`{\displaystyle \frac{\omega (B;k)}{B}}`$ $`=`$ $`4kv^2(B;k)c_k`$ (248) with $$c_k=\frac{\mathrm{cos}(\pi \alpha ^2(B_k))}{\mathrm{\Gamma }^2(\alpha ^2(B_k))}\frac{\stackrel{~}{u}^2(B_k)}{a^2}\|ak\|^{2\alpha ^2(B_k)2}.$$ (249) Now this approximated system of flow equations can be solved easily. One finds for $`0c_k>k^2`$ $$\omega (B;k)=k+\sqrt{k^2+c_k}\mathrm{coth}\left(4k\sqrt{k^2+c_k}B+\mathrm{arccoth}\left(\frac{k}{\sqrt{k^2+c_k}}\right)\right),$$ (250) and for $`c_k0`$ $$\omega (B;k)=k+\sqrt{k^2+c_k}\mathrm{tanh}\left(4k\sqrt{k^2+c_k}B+\mathrm{arctanh}\left(\frac{k}{\sqrt{k^2+c_k}}\right)\right).$$ (251) In both cases one obtains $$\omega (B=\mathrm{};k)=k+\sqrt{k^2+c_k}.$$ (252) In the case $`c_kk^2`$ the solution for $`\omega (B;k)`$ diverges as $`B\mathrm{}`$. For small initial couplings $`u_0`$ this scenario can according to (249) only occur for $`\alpha ^2(B=0)<1/2`$. This just defines our permissible range of paramters $`\beta _0^22\pi `$ as mentioned above. From (250) and (251) we can also read of the justification for our above approximation in the system of differential equations: Nearly all the flow from $`\omega (B=0;k)=0`$ to $`\omega (B=\mathrm{};k)`$ occurs on the scale $`B[4k\sqrt{k^2+c_k}]^1B_k`$ with $`B_k`$ from (131). One can verify numerically that for $`\beta _0^24\pi `$ and $`\|k\|<3\stackrel{~}{u}(\mathrm{})/a`$ the solution (252) of the approximated flow equations agrees to within 20% with the full numerical solution of the equations (207) and (208) and becomes exact for $`\|k\|\stackrel{~}{u}(\mathrm{})/a`$. Putting everything together the dispersion relation is $$E_k=\sqrt{k^2\mathrm{cos}(\pi \alpha ^2(B_k))\left(\frac{1}{\mathrm{\Gamma }(\alpha ^2(B_k))}\frac{\stackrel{~}{u}(B_k)}{a}\|ak\|^{\alpha ^2(B_k)1}\right)^2}.$$ (253) In the low–energy limit the dispersion relation (253) takes different forms in the weak– and strong–coupling phases: * Weak–coupling phase: Here $`\alpha ^2(B=\mathrm{})2`$ and we find for $`\|k\|\stackrel{~}{u}(\mathrm{})/a`$ $$E_k=\|k\|\sqrt{1\mathrm{cos}(\pi \alpha ^2(\mathrm{}))\left(\frac{1}{\mathrm{\Gamma }(\alpha ^2(\mathrm{}))}\stackrel{~}{u}(\mathrm{})\|ak\|^{\alpha ^2(\mathrm{})2}\right)^2},$$ (254) that is a gapless spectrum with $`E_k=\|k\|`$ for $`k0`$. * Strong–coupling phase: Here $`\alpha ^2(B=\mathrm{})=1`$ and we find for $`\|k\|\stackrel{~}{u}(\mathrm{})/a`$ $$E_k=\sqrt{k^2+\left(\frac{\stackrel{~}{u}(\mathrm{})}{a}\right)^2}.$$ (255) This agrees with the result (230) obtained from the effective Hamiltonian analysis in Sect. IV.B.2, that is we find a gapped spectrum with the mass $`M=\stackrel{~}{u}(\mathrm{})/a`$. One can verify numerically that in the strong–coupling phase for $`\beta _0^24\pi `$ the full dispersion relation obtained by solving (207) and (208) is very accurately described by $`\sqrt{k^2+M^2}`$ even in the crossover region: In the small coupling limit $`\|u_0\|1`$ there are $`\beta _0`$–dependent universal corrections in the crossover region $`k=O(M)`$ that vanish for $`\beta _0^24\pi `$ and reach at most $`3\%`$ (for $`\beta _0^2=8\pi `$).<sup>\</sup><sup>\</sup>\**Such small corrections might be expected from the exact results for the one–dimensional spin–1/2 XYZ–chain in Ref. , where the form (255) holds exactly also in the crossover region. However, a strict comparison with is difficult since there are nontrivial renormalization subtleties in mapping the continuum sine–Gordon model to the disrete spin chain (in this context see e.g. Ref. ). The respective scaling forms of the dispersion relation are depicted in Fig. 5. Notice that according to (133) the scaling dimension of our single–particle/hole excitations varies continuously along these dispersion curves, see Figs. 5 and 6: In the strong–coupling phase one finds the initial scaling dimension $`\alpha (0)`$ for excitations with large momenta $`\|k\|M`$, and the low–energy effective Thirring fermions with $`\alpha (\mathrm{})=1`$ for small momenta $`\|k\|M`$. This is consistent with the exact $`S`$–matrix results by Zamolodchikov discussed in Sect. II.C. Also notice that our elementary excitations in this section are described with respect to a transformed basis since $`H(B=\mathrm{})`$ and $`H(B=0)`$ are related by a complicated unitary transformation. Finally let us look at the solution of the equations (207) and (208) for initial parameters $`\beta _0^2<4\pi `$. We have seen above that the differential equation for $`\omega (B;k)`$ leads to divergences for $`k0`$ if $`\alpha _0^2<1/2`$.<sup>††</sup><sup>††</sup>†† In fact the $`\mathrm{sin}(\pi \alpha ^2)`$–term in (208) already leads to IR–problems for $`\alpha _0^2<1`$. The source of this problem is related to the breakdown at $`\alpha _0^2=1/2`$ and can be resolved in a likewise manner as will be shown in a subsequent publication: The $`\mathrm{sin}(\pi \alpha ^2)`$–terms turn out to be generally unimportant and we can safely neglect them in our present discussion also for $`\alpha _0^2<1`$. In other words the $`\mathrm{cos}(\beta _0\varphi (x))`$–term becomes too relevant for our flow equation approach in its present form for $`\beta _0^2<2\pi `$. This is the reason why the parameter space of the sine–Gordon model that we can deal with in this paper is restricted to $`\beta _0^22\pi `$. In the interval $`2\pi \beta _0^2<4\pi `$ our approximations become exact in the limit $`\beta _0^24\pi `$ and, according to the observations above, eventually break down for $`\beta _0^2<2\pi `$. We will investigate the accuracy of our approximations in more detail in the following section. ### D Approximations and expansion parameter Both during the flow and for the analysis of the final diagonal Hamiltonian we have neglected the terms in $`H_{\mathrm{res}}(B)`$ in (202). For any given $`B`$–scale these terms are more irrelevant by at least two spatial derivatives than the interaction term $`H_{\mathrm{int}}(B)`$. Our approximations are therefore similar to the expansions in renormalization approaches. To judge the accuracy of our approximations in more detail it is important to study the prefactors of these terms, that is to find the expansion parameter of our approach. Since $`H_{\mathrm{int}}(B)`$ is treated as the perturbing term, we can do this easily by comparing it with $`H_0`$: E.g. from the dispersion relation (253) one sees that on the momentum scale $`k`$ the effect of $`H_{\mathrm{int}}(B)`$ as compared to $`H_0`$ is $`{\displaystyle \frac{\|ak\|^{\alpha ^2(B_k)}\stackrel{~}{u}(B_k)/a}{\|k\|}}`$ $``$ $`\stackrel{~}{u}(B_k)\left({\displaystyle \frac{\sqrt{B_k}}{a}}\right)^{2\alpha ^2(B_k)}`$ (256) $`=`$ $`u(B_k),`$ (257) where we have used (216). Not surprisingly the expansion parameter seems to be the running coupling constant $`u(B)`$ like in perturbative RG. According to the systems of differential equations (218) this is good news in the weak–coupling phase since $`u(B)`$ decays to zero on small energy scales ($`B\mathrm{}`$). However, in the strong–coupling phase $`u(B)`$ diverges according to (232). Still our method is a systematic approximation even in the strong–coupling phase since our solution becomes exact for $`\beta ^2=4\pi `$. As we have seen in Sect. III, our system of flow equations closes for $`\beta ^2=4\pi `$ ($`\alpha ^2=1`$) and no higher order interactions are generated during the flow. Therefore $`H_{\mathrm{res}}(B)`$ vanishes identically on this line. This remarkable observation is our main difference from perturbative RG. It is in fact even a trivial observation since according to Sect. II.C the sine–Gordon Hamiltonian becomes equivalent to a noninteracting Thirring model for $`\beta ^2=4\pi `$. And a quadratic Hamiltonian can easily be solved exactly with our scheme of unitary transformations. Therefore the true expansion parameter of our method is necessarily some product of $`u(B)`$ and $`(\alpha ^2(B)1)`$. One can check explicitly that the leading terms in $`H_{\mathrm{res}}(B)`$ contribute like $$g(B)\stackrel{\mathrm{def}}{=}u^2(B)(\alpha ^2(B)1).$$ (258) In this context also compare our analysis in Sect. III.D.3 where we have seen in (98) that the $`R`$–term (that we had to treat approximately) vanishes linearly like $`(\alpha ^21)`$ as $`\alpha ^21`$. This dimensionless combination $`g(B)`$ is therefore the expansion parameter of our scheme<sup>‡‡</sup><sup>‡‡</sup>‡‡Also the combination $`g(B)u(B)`$ appears, but this does not make any difference for our analysis. and has to remain small throughout the entire flow in order to have a systematic and controllable approximation. In order to verify this let us first investigate how $`(\alpha ^2(B)1)`$ vanishes for $`B\mathrm{}`$ in the strong–coupling phase: We define $`ϵ(\mathrm{})\stackrel{\mathrm{def}}{=}\alpha ^2(\mathrm{})1`$ and approximate (212) for large $`\mathrm{}`$ yielding $$\frac{dϵ(\mathrm{})}{d\mathrm{}}=ϵ(\mathrm{})a^2M^2e^2\mathrm{},$$ (259) which can be solved easily leading to $$ϵ(\mathrm{})=ϵ(\mathrm{}_0)\mathrm{exp}\left[\frac{(aM)^2}{2}(e^2\mathrm{}e^{2\mathrm{}_0})\right].$$ (260) This very fast decay for $`\mathrm{}\mathrm{}`$ avoids a divergence in our expansion parameter $`g(\mathrm{})`$ $$g(\mathrm{})(aM)^2\mathrm{exp}\left[2\mathrm{}\frac{(aM)^2}{2}e^2\mathrm{}\right]\stackrel{\mathrm{}\mathrm{}}{}0.$$ (261) For the full flow of $`g(\mathrm{})`$ one finds universal curves in the small coupling limit $`u_00`$ that depend on $`\beta _0^2`$. Numerical solutions for these universal curves are shown in Fig. 7. The fact that one finds universal nonvanishing curves for $`\beta _0^24\pi `$ means that there are nonzero corrections to our present approximations even in the limit $`u_00`$. This is not surprising due to the strong–coupling nature of our model. A precise statement about the actual size of the errors is difficult since this depends on unknown prefactors with which $`g(\mathrm{})`$ enters. However, it is encouraging to observe that $`\|g(\mathrm{})\|`$ remains relatively small throughout the entire flow for all $`\beta _0^24\pi `$. Higher order corrections can therefore be expected to be small and are systematically obtainable in an expansion that takes more and more irrelevant terms into account in our flow equation procedure. Finally, for $`\beta _0^2<4\pi `$ the maximum in $`\|g(\mathrm{})\|`$ grows rapidly as $`\beta _0^2`$ becomes smaller as can be seen in Fig. 7. The error of our approximation therefore becomes larger as $`\beta _0^20`$, which agrees with the observation in Sect. IV.C that our present scheme actually breaks down for $`\beta _0^2<2\pi `$. This is not unexpected since the perturbing $`\mathrm{cos}(\beta \varphi (x))`$–term becomes more and more relevant for small $`\beta _0`$. ## V Conclusions In this paper we have developed a new approach for solving the quantum sine–Gordon model (21) $$H(B=0)=𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2+\frac{u}{\pi a^2}\mathrm{cos}\left[\beta 𝑑ϵc(ϵ)\varphi (x+ϵ)\right]\right)$$ (262) by means of infinitesimal unitary transformations as introduced by Wegner and Głazek and Wilson . Within an approximation that neglected operators with larger scaling dimensions (more irrelevant terms) we obtained a flow that unitarily linked the initial Hamiltonian (262) to a diagonal Hamiltonian $`H(B=\mathrm{})`$ $`=`$ $`{\displaystyle 𝑑x\left(\frac{1}{2}\mathrm{\Pi }^2(x)+\frac{1}{2}\left(\frac{\varphi }{x}\right)^2\right)}`$ (264) $`+{\displaystyle \underset{k>0}{}}\omega (B=\mathrm{};k)\left(P_1^{}(k)P_1^{}(k)+P_1^{}(k)P_1^{}(k)+P_2^{}(k)P_2^{}(k)+P_2^{}(k)P_2^{}(k)\right).`$ Here the $`P_j(k)`$ are soliton and antisoliton creation and annihilation operators as defined in (133). Their dispersion relation $`\pm E_k`$ was calculated in (253). In the small coupling limit $`\|u\|1`$ we found $`E_k=\sqrt{k^2+M^2}`$ (with very small deviations from this form) in the strong–coupling phase, and a gapless spectrum $`E_k=\|k\|`$ in the weak–coupling phase. In the strong–coupling phase, our low–energy solitons and antisolitons are fermionic (compare Fig. 6) as known from the exact $`S`$–matrix solution . Within our approach their mass $`M`$ can be obtained by the solution of the flow equations (218) $`{\displaystyle \frac{d\beta ^2(\mathrm{})}{d\mathrm{}}}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi \mathrm{\Gamma }\left(1+\beta ^2(\mathrm{})/4\pi \right)}}u^2(\mathrm{})`$ (265) $`{\displaystyle \frac{du}{d\mathrm{}}}`$ $`=`$ $`\left(2{\displaystyle \frac{\beta ^2(\mathrm{})}{4\pi }}\right)u(\mathrm{})`$ (266) via the relation $`M=e^{\mathrm{}}u(\mathrm{})/a`$ in the limit $`\mathrm{}\mathrm{}`$. The above equations have to be integrated from their initial values for $`\mathrm{}=0`$ to $`\mathrm{}=\mathrm{}`$. Our results for the scaling behavior of the mass agree with exact methods and the two–loop scaling analysis (compare Fig. 4). We also reproduce the phase diagram Fig. 1 of the sine–Gordon model in our approach. Our equations (265) and (266) are similar to the two–loop RG equations (23) except for the $`\mathrm{\Gamma }`$–function in the denominator of (265) that makes $`\beta ^2=4\pi `$ an attractive strong–coupling fixed point of the flow equation method (see Fig. 3). This is the main difference between our approach and perturbative RG. Since the sine–Gordon model for $`\beta ^2=4\pi `$ can be interpreted as a noninteracting Thirring model, our flow equation procedure becomes exact at this point and diagonalizes the ensuing quadratic Hamiltonian easily. The expansion parameter of our approach is therefore not $`u(\mathrm{})`$ (notice that $`u(\mathrm{})`$ diverges in the strong–coupling phase as in perturbative scaling), but according to (258) the product $`g(\mathrm{})=u^2(\mathrm{})(1+\beta ^2(\mathrm{})/4\pi )`$. $`g(\mathrm{})`$ remains small throughout the entire flow for $`\beta ^24\pi `$ (compare Fig. 7). This allows a systematic improvement of our present approximations by successively taking terms with larger scaling dimensions into account in our flow equation procedure. Furthermore, it allows us to conclude that our present approximation already provides a good description of the crossover region, which is notoriously difficult to study with other techniques. For example, we worked out the dispersion relation of the single–particle/hole excitations for all momenta in Fig. 5. Notice that higher order terms in our expansion cannot endanger the stability of the strong–coupling fixed point. As can be deduced from Fig. 7, our approximations become less accurate for $`\beta ^2<4\pi `$ and eventually our present approach breaks down for $`\beta ^2<2\pi `$ (Sect. IV.C). In addition, the bound states present in the spectrum for $`\beta ^2<4\pi `$ according to the exact solution are absent in our solution. These bound states will be generated by interactions in $`H_{\mathrm{res}}`$ that are not included in our present approximation. Work on these issues for $`\beta ^2<4\pi `$ is in progress. To summarize, we have obtained an explicit approximate relation between the strong–coupling problem (262) and its diagonalized form (264) without using the integrable structure of the model. The method presented here provides a theoretical tool that is capable of achieving this in a systematic expansion throughout the crossover region. We have been able to carry out this program completely for $`\beta ^24\pi `$, and obtained first results for $`\beta ^2<4\pi `$ (where more work remains to be done e.g. regarding the bound states). Our approach is conceptually simple since a small parameter is identified and used as an expansion parameter. It is therefore possible to study nonintegrable perturbations with our approach, and the calculation of correlation functions also seems feasible. Finally, there are various other onedimensional strong–coupling problems, as for example the Kondo model, where the present approach should be useful . The author acknowledges many valuable discussions with W. Hofstetter and D. S. Fisher. This work was supported by the Deutsche Forschungsgemeinschaft (DFG), by the SFB 484 of the DFG and by the National Science Foundation (NSF) under grants DMR 9630064, DMR 9976621 and DMR 9981283. ## Properties of vertex operators ### 1 Operator product expansion This Appendix compiles some important properties of vertex operators. For a review of these properties see also Ref. . We first want to establish a relation between the normal–ordered and the non–normal–ordered vertex operators. By definition $`V_1(\alpha ;x)`$ $`\stackrel{\mathrm{def}}{=}`$ $`:\mathrm{exp}\left(\alpha {\displaystyle \underset{p0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)\right):`$ (267) $`=`$ $`\mathrm{exp}\left(\alpha {\displaystyle \underset{p>0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)\right)\mathrm{exp}\left(\alpha {\displaystyle \underset{p<0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)\right)`$ (268) and next we can use the formula $$e^Ae^B=e^{C/2}e^{A+B}$$ (269) with $`C=[A,B]`$ since $`C`$ is a number and commutes with both $`A`$ and $`B`$: $`C`$ $`=`$ $`[\alpha {\displaystyle \underset{p>0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p),\alpha {\displaystyle \underset{q<0}{}}{\displaystyle \frac{\sqrt{\|q\|}}{q}}e^{\frac{a}{2}\|q\|iqx}\sigma _1(q)]`$ (270) $`=`$ $`\alpha ^2{\displaystyle \underset{p>0}{}}{\displaystyle \frac{e^{a\|p\|}}{p}}.`$ (271) The sum over $`q`$ yields $$C=\alpha ^2\mathrm{ln}\left(1e^{2\pi a/L}\right)$$ (272) and in the thermodynamic limit $`L\mathrm{}`$ $$C=\alpha ^2\mathrm{ln}\left(\frac{2\pi a}{L}\right).$$ (273) Therefore $`V_j(\alpha ;x)`$ $`\stackrel{\mathrm{def}}{=}`$ $`:\mathrm{exp}(\pm \alpha {\displaystyle \underset{p0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _j(p)):`$ (274) $`=`$ $`\left({\displaystyle \frac{L}{2\pi a}}\right)^{\alpha ^2/2}\mathrm{exp}\left(\pm \alpha {\displaystyle \underset{p0}{}}{\displaystyle \frac{\sqrt{\|p\|}}{p}}e^{\frac{a}{2}\|p\|ipx}\sigma _j(p)\right).`$ (275) An important property of the Fourier–transformed vertex operators (103) is their action on the vacuum $$V_1(\alpha ;k)\|\mathrm{\Omega }=V_1(\alpha ;k)\|\mathrm{\Omega }=V_2(\alpha ;k)\|\mathrm{\Omega }=V_2(\alpha ;k)\|\mathrm{\Omega }=0k>0.$$ (276) This is shown easily, e.g. $`V_1(\alpha ;k)\|\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑xe^{ikx}V_1(\alpha ;x)\|\mathrm{\Omega }}`$ (277) $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑xe^{ikx}\mathrm{exp}\left(\alpha \underset{p>0}{}\frac{\sqrt{\|p\|}}{p}e^{\frac{a}{2}\|p\|ipx}\sigma _1(p)\right)\|\mathrm{\Omega }}.`$ (278) One expands the second exponential and arrives at terms with the structure $$𝑑xe^{ix(k+p_1+\mathrm{}+p_n)}=0$$ (279) since $`k,p_1,\mathrm{},p_n>0`$. Therefore $$V_1(\alpha ;k)\|\mathrm{\Omega }=0$$ (280) for $`k>0`$ and the analysis for the rest of (276) proceeds likewise. Next we want to want to evaluate expectation values for products of vertex operators. We look at the product $`V_1(\alpha ;x)V_1(\alpha ;y)`$ $`=`$ $`\mathrm{exp}\left(\alpha {\displaystyle \underset{p>0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|p\|ipx)}{\sqrt{\|p\|}}}\sigma _1(p)\right)\mathrm{exp}\left(\alpha {\displaystyle \underset{p<0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|p\|ipx)}{\sqrt{\|p\|}}}\sigma _1(p)\right)`$ (282) $`\times \mathrm{exp}\left(\alpha {\displaystyle \underset{q>0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|q\|iqy)}{\sqrt{\|q\|}}}\sigma _1(q)\right)\mathrm{exp}\left(\alpha {\displaystyle \underset{q<0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|q\|iqy)}{\sqrt{\|q\|}}}\sigma _1(q)\right)`$ and commute the second and third exponentials using the formula $$e^Ae^B=e^Ce^Be^A$$ (283) with $`C=[A,B]`$ since again $`C`$ is a number and commutes with both $`A`$ and $`B`$: $`C`$ $`=`$ $`[\alpha {\displaystyle \underset{p<0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|p\|ipx)}{\sqrt{\|p\|}}}\sigma _1(p),\alpha {\displaystyle \underset{q>0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|q\|iqy)}{\sqrt{\|q\|}}}\sigma _1(q)]`$ (284) $`=`$ $`\alpha ^2{\displaystyle \underset{q>0}{}}{\displaystyle \frac{\mathrm{exp}(a\|q\|+iq(xy))}{q}}`$ (285) $`=`$ $`\alpha ^2\mathrm{ln}\left(1e^{i\frac{2\pi }{L}(xy+ia)}\right).`$ (286) In the thermodynamic limit $`L\mathrm{}`$ this gives $$e^C=\left(\frac{L/2\pi }{i(yx)+a}\right)^{\alpha ^2}.$$ (287) Therefore $`V_1(\alpha ;x)V_1(\alpha ;y)`$ $`=`$ $`\left({\displaystyle \frac{L/2\pi }{i(yx)+a}}\right)^{\alpha ^2}\mathrm{exp}\left(\alpha {\displaystyle \underset{p>0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|p\|ipx)}{\sqrt{\|p\|}}}\sigma _1(p)\right)\mathrm{exp}\left(\alpha {\displaystyle \underset{q>0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|q\|iqy)}{\sqrt{\|q\|}}}\sigma _1(q)\right)`$ (289) $`\times \mathrm{exp}\left(\alpha {\displaystyle \underset{p<0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|p\|ipx)}{\sqrt{\|p\|}}}\sigma _1(p)\right)\mathrm{exp}\left(\alpha {\displaystyle \underset{q<0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|q\|iqy)}{\sqrt{\|q\|}}}\sigma _1(q)\right).`$ From this we obtain the expectation value $$V_1(\alpha ;x)V_1(\alpha ;y)=\left(\frac{L/2\pi }{i(yx)+a}\right)^{\alpha ^2}.$$ (290) A similar calculation can be done for $`V_2(\alpha ;x)`$ and the only difference is the exchange of $`x`$ and $`y`$ in the denominator $$V_2(\alpha ;x)V_2(\alpha ;y)=\left(\frac{L/2\pi }{i(xy)+a}\right)^{\alpha ^2}.$$ (291) The operator product expansion (OPE) for vertex operators can be deduced from (289) $`V_1(\alpha ;x)V_1(\alpha ;y)`$ $`=`$ $`\left({\displaystyle \frac{L/2\pi }{i(yx)+a}}\right)^{\alpha ^2}\mathrm{exp}\left(\alpha {\displaystyle \underset{p>0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|p\|)}{\sqrt{\|p\|}}}(e^{ipx}e^{ipy})\sigma _1(p)\right)`$ (293) $`\times \mathrm{exp}\left(\alpha {\displaystyle \underset{p<0}{}}{\displaystyle \frac{\mathrm{exp}(\frac{a}{2}\|p\|)}{\sqrt{\|p\|}}}(e^{ipx}e^{ipy})\sigma _1(p)\right)`$ $`=`$ $`\left({\displaystyle \frac{L/2\pi }{i(yx)+a}}\right)^{\alpha ^2}\mathrm{exp}\left(i\alpha (xy){\displaystyle \underset{p>0}{}}e^{\frac{a}{2}\|p\|ipx}\sqrt{\|p\|}\sigma _1(p)+O((xy)^2)\right)`$ (295) $`\times \mathrm{exp}\left(i\alpha (xy){\displaystyle \underset{p<0}{}}e^{\frac{a}{2}\|p\|ipx}\sqrt{\|p\|}\sigma _1(p)+O((xy)^2)\right)`$ $`=`$ $`\left({\displaystyle \frac{L/2\pi }{i(yx)+a}}\right)^{\alpha ^2}\left(1+i\alpha (yx){\displaystyle \underset{p0}{}}e^{\frac{a}{2}\|p\|ipx}\sqrt{\|p\|}\sigma _1(p)+O((xy)^2)\right).`$ (296) Higher order terms in the OPE can easily be deduced using the above scheme. These terms are less singular as $`xy`$, or, in the language of renormalization theory, they have a larger scaling dimension (are more irrelevant) and can be expressed as spatial derivatives of the bosonic field. A similar calculation for $`V_2(\alpha ;x)`$ gives $$V_2(\alpha ;x)V_2(\alpha ;y)=\left(\frac{L/2\pi }{i(xy)+a}\right)^{\alpha ^2}\left(1+i\alpha (xy)\underset{p0}{}e^{\frac{a}{2}\|p\|ipx}\sqrt{\|p\|}\sigma _2(p)+O((xy)^2)\right),$$ (297) again the only difference is the exchange of $`x`$ and $`y`$. It is well–known that for $`\alpha =\pm 1`$ the vertex operators describe fermion creation and annihilation operators. This can be checked easily by using the OPE (296) in the anticommutator for the special case $`\alpha =1`$ $`\{V_1(1;x),V_1(1;y)\}`$ $`=`$ $`\left({\displaystyle \frac{L/2\pi }{i(yx)+a}}+{\displaystyle \frac{L/2\pi }{i(xy)+a}}\right)\left(1+i\alpha (yx){\displaystyle \underset{p0}{}}e^{\frac{a}{2}\|p\|ipx}\sqrt{\|p\|}\sigma _1(p)+O((xy)^2)\right)`$ (298) $`\stackrel{a0}{=}`$ $`L\delta (xy)\left(1+i\alpha (yx){\displaystyle \underset{p0}{}}e^{\frac{a}{2}\|p\|ipx}\sqrt{\|p\|}\sigma _1(p)+O((xy)^2)\right)`$ (299) $`=`$ $`L\delta (xy)`$ (300) in the limit $`a0`$. All higher order terms in the OPE vanish in this limit. Likewise one finds $$\{V_1(1;x),V_1(1;y)\}=\{V_1(1;x),V_1(1;y)\}\stackrel{a0}{=}0$$ (301) and the same relations for $`V_2(\pm 1;x)`$. ### 2 Exchange relations Let us look at the commutation relation of vertex operators in momentum space. These are worked out in the following for general $`\alpha `$. For simplicity we only consider $`V_1(\pm \alpha ;x)`$, the calculation for $`V_2(\pm \alpha ;x)`$ proceeds along similar lines. From the definition of the vertex operators one easily verifies the following relation $$V_1(\alpha ;x)V_1(\alpha ;y)=\frac{[i(yx)+a]^{\alpha ^2}}{[i(yx)+a]^{\alpha ^2}}V_1(\alpha ;y)V_1(\alpha ;x).$$ (302) For large distances $`\|xy\|a`$ the coefficient can be well approximated by $$\mathrm{cos}(\pi \alpha ^2)+i\mathrm{sgn}(xy)\mathrm{sin}(\pi \alpha ^2),$$ (303) whereas for small distances it becomes equal to 1. Now for small distances the operator product expansion for the two vertex operators can be used and it is then possible to write generally for all $`xy`$ $$V_1(\alpha ;x)V_1(\alpha ;y)\mathrm{OPE}(xy)=\left(\mathrm{cos}(\pi \alpha ^2)+i\mathrm{sgn}(xy)\mathrm{sin}(\pi \alpha ^2)\right)\left(V_1(\alpha ;y)V_1(\alpha ;x)\mathrm{OPE}(xy)\right)$$ (304) For our purposes here it will be sufficient to look only at the leading $`c`$–number term in the OPE (higher orders can easily be taken into account if necessary). This is equivalent to subtracting the ground state expectation value on both sides (56). Notice that our relation “closes” $`V_1(\alpha ;x)V_1(\alpha ;y)`$ $`=`$ $`(\mathrm{cos}(\pi \alpha ^2)+i\mathrm{sgn}(xy)\mathrm{sin}(\pi \alpha ^2))V_1(\alpha ;y)V_1(\alpha ;x)`$ (305) $`=`$ $`\left(\mathrm{cos}(\pi \alpha ^2)+i\mathrm{sgn}(xy)\mathrm{sin}(\pi \alpha ^2)\right)`$ (307) $`\times (\mathrm{cos}(\pi \alpha ^2)+i\mathrm{sgn}(yx)\mathrm{sin}(\pi \alpha ^2))V_1(\alpha ;x)V_1(\alpha ;y)`$ $`=`$ $`V_1(\alpha ;x)V_1(\alpha ;y)`$ (308) In terms of Fourier components (103) this reads ($`\alpha >0`$) $`V_1(\alpha ;k_1)V_1(\alpha ;k_2)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{p}{}}V_1(\alpha ;k_2+p)V_1(\alpha ;k_1+p){\displaystyle 𝑑xe^{ipx}\left(\mathrm{cos}(\pi \alpha ^2)+i\mathrm{sgn}(x)\mathrm{sin}(\pi \alpha ^2)\right)}`$ (309) $`=`$ $`\mathrm{cos}(\pi \alpha ^2)V_1(\alpha ;k_2)V_1(\alpha ;k_1)`$ (312) $`+\mathrm{sin}(\pi \alpha ^2){\displaystyle \frac{2}{\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{2n+1}}[V_1(\alpha ;k_2{\displaystyle \frac{2\pi }{L}}(2n+1))V_1(\alpha ;k_1{\displaystyle \frac{2\pi }{L}}(2n+1))`$ $`V_1(\alpha ;k_2+{\displaystyle \frac{2\pi }{L}}(2n+1))V_1(\alpha ;k_1+{\displaystyle \frac{2\pi }{L}}(2n+1))]`$ $`=`$ $`\mathrm{cos}(\pi \alpha ^2)V_1(\alpha ;k_2)V_1(\alpha ;k_1)`$ (314) $`{\displaystyle \frac{1}{\pi }}\mathrm{sin}(\pi \alpha ^2)\mathrm{Re}{\displaystyle }dq{\displaystyle \frac{1}{q+iϵ}}V_1(\alpha ;k_2+q)V_1(\alpha ;k_1+q)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}[e^{i\pi \alpha ^2}{\displaystyle }dq{\displaystyle \frac{1}{q+iϵ}}V_1(\alpha ;k_2+q)V_1(\alpha ;k_1+q)]`$ (315) in the thermodynamic limit with $`lim_{ϵ0}`$ being understood. Likewise $$V_1(\alpha ;k_1)V_1(\alpha ;k_2)=\frac{1}{\pi }\mathrm{Im}[e^{i\pi \alpha ^2}dq\frac{1}{q+iϵ}V_1(\alpha ;k_2+q)V_1(\alpha ;k_1+q)].$$ (316) Notice that these relations close in $`k`$–space in the same sense as (308). They are only “simple” for integer $`\alpha ^2`$ when the vertex operators behave as bosons or fermions: Then the term proportional to $`\mathrm{sin}(\pi \alpha ^2)`$ with its summation over wavevectors vanishes.
warning/0006/hep-ph0006120.html	ar5iv	text	# Additional Isospin-Breaking Effects in ϵ'/ϵ ## Introduction The recent measurement of a non-zero value of $`\mathrm{Re}(ϵ^{}/ϵ)`$ ktev establishes the existence of CP violation in direct decay and thus provides an important first check of the mechanism of CP violation in the Standard Model (SM). Nevertheless, the world average which emerges is $`\mathrm{Re}(ϵ^{}/ϵ)=(19.3\pm 2.4)10^4`$ world , which is larger than the “central” SM prediction of $`7.010^4`$ burasr ; buraslh by nearly a factor of three. This compels us to scrutinize the SM prediction in further detail: we study isospin-violating effects arising from the $`u`$-$`d`$ quark mass difference. Isospin violation plays an important role in the analysis of $`ϵ^{}/ϵ`$, for the latter is predicated by the difference of the imaginary to real part ratios in the $`\|\mathrm{\Delta }I\|=1/2`$ and $`\|\mathrm{\Delta }I\|=3/2`$ $`K\pi \pi `$ amplitudes. The differing charges of the $`u`$ and $`d`$ quarks engender $`\|\mathrm{\Delta }I\|=3/2`$ electroweak penguin contributions, whereas $`\pi ^0`$-$`\eta ,\eta ^{}`$mixing, driven by the $`u`$-$`d`$ quark mass difference, modifies the relative contribution of the $`\|\mathrm{\Delta }I\|=1/2`$ and $`\|\mathrm{\Delta }I\|=3/2`$ amplitudes in a significant way. Here we describe isospin-breaking effects in the matrix elements of the gluonic penguin operators sggv , such as $`Q_6`$. These operators have always been thought to induce exclusively $`\|\mathrm{\Delta }I\|=1/2`$ transitions, but this is true only in the limit of isospin symmetry. The difference in the up and down quark masses effectively distinguishes the interaction of gluons with up and down quarks, so that the $`\pi \pi \|Q_6\|K`$ matrix element possesses a $`\|\mathrm{\Delta }I\|=3/2`$ component as well bpipi . Let us begin by showing why the numerical prediction of $`ϵ^{}/ϵ`$is sensitive to the presence of isospin violation. The value of $`ϵ^{}/ϵ`$is inferred from a ratio of ratios, namely $$\mathrm{Re}\left(\frac{\epsilon ^{}}{ϵ}\right)=\frac{1}{6}\left[\left\|\frac{\eta _+}{\eta _{00}}\right\|^21\right],$$ (1) where $$\eta _+\frac{𝒜(K_L\pi ^+\pi ^{})}{𝒜(K_S\pi ^+\pi ^{})}ϵ+ϵ^{};\eta _{00}\frac{𝒜(K_L\pi ^0\pi ^0)}{𝒜(K_S\pi ^0\pi ^0)}ϵ2ϵ^{}.$$ (2) In the isospin-perfect limit, the two independent amplitudes present in $`K\pi \pi `$ decay are distinguished by the isospin of the final-state pions, namely $`A_I𝒜(K(\pi \pi )_I)`$ with $`I=0,2`$. $`ϵ^{}/ϵ`$can thus be written $$\frac{ϵ^{}}{ϵ}=\frac{\omega }{\sqrt{2}\|ϵ\|}\xi (1\mathrm{\Omega }),$$ (3) with $$\omega \frac{\mathrm{Re}A_2}{\mathrm{Re}A_0};\xi \frac{\mathrm{Im}A_0}{\mathrm{Re}A_0};\mathrm{\Omega }\frac{\mathrm{Im}A_2}{\omega \mathrm{Im}A_0}.$$ (4) In standard practice, $`\omega 1/22`$ and $`\mathrm{Re}A_0`$ are taken from experiment whereas $`\mathrm{Im}A_I`$ is computed using the operator-product expansion burasr ; buraslh , that is, via $$_{\mathrm{eff}}(\|\mathrm{\Delta }S\|=1)=4\frac{G_\mathrm{F}}{\sqrt{2}}V_{us}^{}V_{ud}\underset{i=1}{\overset{10}{}}C_i(\mu )Q_i(\mu )+\mathrm{h}.\mathrm{c}.$$ (5) The numerical value of $`ϵ^{}/ϵ`$is driven by the matrix elements of the QCD penguin operator $`Q_6`$ and the electroweak penguin operator $`Q_8`$ buras93 . Writing $`Q_i_I`$ as $`(\pi \pi )_I\|Q_i\|KB_i^{((I+1)/2)}(\pi \pi )_I\|Q_i\|K^{(\mathrm{vac})}`$, where “vac” indicates the use of the vacuum saturation approximation, one recovers the schematic formula burasr $$\frac{ϵ^{}}{ϵ}=13\mathrm{Im}\lambda _t\left[B_6^{(1/2)}(1\mathrm{\Omega }_{\eta +\eta ^{}})0.4B_8^{(3/2)}\right].$$ (6) Using $`B_6^{(1/2)}=1.0`$, $`B_8^{(3/2)}=0.8`$, and $`\mathrm{\Omega }_{\eta +\eta ^{}}=0.25`$ yields the “central” SM value of $`ϵ^{}/ϵ7.010^4`$ burasr , roughly a factor of three smaller than the measured value. Larger estimates of $`B_6^{(1/2)}`$, and hence of $`ϵ^{}/ϵ`$, exist bertolini ; hambye ; bijpra ; we investigate sources of $`\mathrm{\Omega }_{\eta +\eta ^{}}`$. Note that under $`\mathrm{\Omega }_{\eta +\eta ^{}}\mathrm{\Omega }_{\eta +\eta ^{}}`$, $`\mathrm{Re}(ϵ^{}/ϵ)`$$``$ 2.2 $`\mathrm{Re}(ϵ^{}/ϵ)`$. Were $`\|B_6^{(1/2)}\|\|B_8^{(3/2)}\|`$, flipping the sign of $`\mathrm{\Omega }_{\eta +\eta ^{}}`$ would increase $`ϵ^{}/ϵ`$by a factor of 1.7. Let us consider possible sources of $`\mathrm{\Omega }_{\eta +\eta ^{}}`$. We replace $`\mathrm{\Omega }_{\eta +\eta ^{}}`$ by $`\mathrm{\Omega }_{\mathrm{IB}}`$, where $$\mathrm{\Omega }_{\mathrm{IB}}=\left(\frac{\sqrt{2}}{3\omega }\right)\frac{\mathrm{Im}(A_\mathrm{P}(K^0\pi ^+\pi ^{})A_\mathrm{P}(K^0\pi ^0\pi ^0))}{\mathrm{Im}A_\mathrm{P}(K^0\pi \pi )}$$ (7) and $`\mathrm{Im}A_\mathrm{P}(K^0\pi \pi )=(\mathrm{Im}A_\mathrm{P}(K^0\pi ^+\pi ^{})+\mathrm{Im}A_\mathrm{P}(K^0\pi ^0\pi ^0))/2`$. “$`A_\mathrm{P}`$” denotes an amplitude induced by $`(8_L,1_R)(\mathrm{e}.\mathrm{g}.,Q_6)`$ operators — the empirical $`\|\mathrm{\Delta }I\|=1/2`$ rule suggests such operators dominate the isospin-violating effects. $`\mathrm{\Omega }_{\mathrm{IB}}`$ vanishes in the absence of isospin violation, i.e., if $`m_u=m_d`$, $`e_u=e_d`$. It can be generated by both strong-interaction and electromagnetic effects, mediated by $`m_dm_u`$ and $`e_ue_d`$ cirigli , respectively. We focus on $`m_dm_u`$ effects. The latter include $`\pi ^0`$-$`\eta ,\eta ^{}`$mixing bijwi ; dght ; buge ; in $`𝒪(p^2,1/N_c)`$ this yields $`\mathrm{\Omega }_{\eta +\eta ^{}}=0.25\pm 0.05`$ dght ; buge , used in the analysis of Ref. burasr . However, $`m_um_d`$ effects can also spawn a $`\|\mathrm{\Delta }I\|=3/2`$ component in the matrix elements of the gluonic penguin operators bpipi , as illustrated in Fig. 1. We turn to a chiral Lagrangian analysis in order to estimate the size of this effect sggv . ## Chiral Lagrangian Analysis The weak chiral Lagrangian for $`K\pi \pi `$ decay is written in terms of the unitary matrix $`U=\mathrm{exp}(i\varphi /f)`$ and the function $`\chi `$, both of which transform as $`URUL^{}`$ under the chiral group $`SU(3)_L\times SU(3)_R`$. The function $`\varphi `$ represents the octet of pseudo-Goldstone bosons, i.e., $`\varphi =_{a=1,\mathrm{},8}\lambda _a\varphi _a`$. In the absence of external fields, $`\chi =2B_0M`$ with $`M=\mathrm{diag}(m_u,m_d,m_s)`$ and $`B_0\overline{q}q`$. The leading-order, $`𝒪(p^2)`$, weak chiral Lagrangian contains no mass-dependent terms cronin , so that $`m_dm_u`$ effects in the hadronization of the gluonic penguin operators first appear in $`𝒪(p^4)`$. This is illustrated in Fig. 2. Let us enumerate the possible isospin-violating effects which occur in $`𝒪(m_dm_u)`$ and $`𝒪(p^4)`$: * $`\pi ^0`$-$`\eta `$ mixing realized from the $`𝒪(p^2)`$strong chiral Lagrangian, in concert with the $`𝒪(p^2)`$weak chiral Lagrangian, computed to one-loop order. * $`\pi ^0`$-$`\eta `$ mixing, realized from the $`𝒪(p^2)`$strong chiral Lagrangian, combined with the isospin-conserving vertices of the $`𝒪(p^4)`$weak chiral Lagrangian. * $`\pi ^0`$-$`\eta `$ mixing as realized from the strong chiral Lagrangian in $`𝒪(p^4)`$, combined with the $`𝒪(p^2)`$weak chiral Lagrangian. The $`\pi ^0`$-$`\eta ^{}`$ mixing effects included in Refs. dght ; buge ape this effect. * Isospin violation in the vertices of the $`𝒪(p^4)`$weak chiral Lagrangian. This serves as our focus here, for it contains the qualitatively new effects we argue. We use the octet terms in the $`𝒪(p^4)`$, CP-odd weak chiral Lagrangian of Ref. kambor . Collecting the $`\chi `$-dependent terms as per iv), working to $`𝒪(m_dm_u)`$, and dropping terms suppressed by $`M_\pi ^2/M_K^2`$, we find $$\mathrm{\Omega }_\mathrm{P}=\frac{2\sqrt{2}}{3\omega }\frac{M_{K^0}^2}{M_{K^0}^2M_\pi ^2}\frac{B_0(m_dm_u)}{c_2^{}}\stackrel{~}{E}^{}\frac{0.12\mathrm{GeV}^2}{c_2^{}}\stackrel{~}{E}^{}$$ (8) with $`\stackrel{~}{E}^{}=2E_1^{}2E_3^{}4E_4^{}E_{10}^{}E_{11}^{}4E_{12}^{}E_{15}^{}`$. Note that $`c_2^{}`$ is the low-energy constant associated with the $`𝒪(p^2)`$, $`(8_L,1_R)`$ weak chiral Lagrangian kambor . As per ii), $`\pi ^0`$-$`\eta `$ mixing in $`𝒪(p^2)`$also enters when combined with the isospin-conserving vertices of the $`𝒪(p^4)`$weak chiral Lagrangian. Including the $`\chi `$-dependent octet terms, we find $$\mathrm{\Omega }_{\eta +\eta ^{}}^{(4)}=\frac{2\sqrt{2}}{3\omega }\frac{M_{K^0}^2}{M_{K^0}^2M_\pi ^2}\frac{B_0(m_dm_u)}{c_2^{}}E_{\eta +\eta ^{}}^{}\frac{0.12\mathrm{GeV}^2}{c_2^{}}E_{\eta +\eta ^{}}^{}$$ (9) with $`E_{\eta +\eta ^{}}^{}=2(E_3^{}+E_4^{}E_5^{})+(E_{10}^{}E_{11}^{})/22E_{12}^{}+E_{14}^{}+3E_{15}^{}/2`$ — so that no manifest cancellation with the terms of Eq. (8) occurs. The low-energy constants $`E_i^{}`$ are unknown, so that we turn to the factorization approximation to proceed. The construction relevant to $`(8_L,1_R)`$ transitions in $`K^0\pi \pi `$ decay is sekhar $`_\mathrm{P}=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{us}^{}V_{ud}C_6\left(8(\overline{s}_Lq_R)(\overline{q}_Rd_L)\right)+\mathrm{h}.\mathrm{c}.`$ (10) $``$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{us}^{}V_{ud}C_6\mathrm{\hspace{0.17em}32}B_0^2{\displaystyle \frac{\delta _{\mathrm{str}}}{\delta \chi _{3i}^{}}}{\displaystyle \frac{\delta _{\mathrm{str}}}{\delta \chi _{i2}}}+\mathrm{h}.\mathrm{c}.,`$ (11) where $`_{\mathrm{str}}`$ is the strong chiral Lagrangian. To generate terms of $`𝒪(p^4)`$in $`_\mathrm{P}`$ requires terms of both $`𝒪(p^4)`$ gl and $`𝒪(p^6)`$ fearing in $`_{\mathrm{str}}`$. Unfortunately, the low-energy constants of the latter are also unknown; the use of “resonance saturation” allows us to estimate some of them. We explicitly consider the scalar nonet of resonances as per Ref. bijnens . An example of the manner in which the scalar resonances can generate contributions to the $`E_i^{}`$ is illustrated in Fig. 3. Integrating out the scalar resonances for $`p^2M_S^2`$, we find two terms which contribute to the scalar densities in the bosonization of $`Q_6`$ bijnens , $$_S^{(6)}=\frac{d_mc_m^2}{2M_S^4}\chi _+^3+\frac{c_dc_md_m}{M_S^4}\chi _+^2L^2,$$ (12) yielding contributions to $`E_1^{}`$ and $`E_{10}^{}`$ in terms of $`d_m`$, $`c_m`$, $`c_d`$, and $`M_S`$. The parameter $`d_m`$ is ill-known; we find $`d_m2.4(0.76)`$. The sign of $`d_m`$ and thus of $`\mathrm{\Omega }_\mathrm{P}`$ in our model results from the mass of the lowest-lying strange scalar being greater than that of the lowest-lying isovector scalar. As per our earlier classification, $`\mathrm{\Omega }_{\mathrm{IB}}^{(4)}=\mathrm{\Omega }_{\mathrm{IB}}^{(4),i}+\mathrm{\Omega }_{\mathrm{IB}}^{(4),ii}+\mathrm{\Omega }_{\mathrm{IB}}^{(4),iii}+\mathrm{\Omega }_{\mathrm{IB}}^{(4),iv}`$, so that with $`d_m=2.4(0.76)`$, we have $`\mathrm{\Omega }_{\mathrm{IB}}^{(4),iv}=0.79(0.21)`$. Estimating $`\mathrm{\Omega }_{\mathrm{IB}}^{(4),ii}`$ using the $`\chi `$-dependent $`E_i^{}`$ yields $`\mathrm{\Omega }_{\mathrm{IB}}^{(4),ii}=0.12(0.03)`$. $`\mathrm{\Omega }_{\mathrm{IB}}^{(4),iii}`$ has been partially determined through the inclusion of $`\pi ^0`$-$`\eta ^{}`$ mixing in $`\mathrm{\Omega }_{\eta +\eta ^{}}=0.25\pm 0.05`$ dght ; buge . Using the result $`\mathrm{\Omega }_{\mathrm{IB}}^{(2)}+\mathrm{\Omega }_{\mathrm{IB}}^{(4),iii}=0.16\pm 0.03`$ ecker and neglecting $`\mathrm{\Omega }_{\mathrm{IB}}^{(4),i}`$, as the ill-known $`E_i^{}`$ do not warrant such a calculation, we estimate, finally, that $`\mathrm{\Omega }_{\mathrm{IB}}=\mathrm{\Omega }_{\mathrm{IB}}^{(2)}+\mathrm{\Omega }_{\mathrm{IB}}^{(4)}0.050.78`$. For reference, note that $`\mathrm{\Omega }_{\mathrm{IB}}^{(2)}0.13`$. The large value of $`\mathrm{\Omega }_{\mathrm{IB}}^{(4)}`$ is driven by the numerical prefactor of Eqs. (8,9) — the contributions in $`\mathrm{\Omega }_{\mathrm{IB}}^{(4)}`$ are “naturally” of the same size as $`\mathrm{\Omega }_{\mathrm{IB}}^{(2)}`$. Thus we find a very large correction to the value of $`\mathrm{\Omega }_{\eta +\eta ^{}}=0.25\pm 0.05`$, used in “central value” of $`ϵ^{}/ϵ`$. The large negative change in $`\mathrm{\Omega }_{\mathrm{IB}}`$ found in $`𝒪(p^4)`$generates a substantial increase in $`ϵ^{}/ϵ`$. The $`\mathrm{\Omega }_{\mathrm{IB}}`$ we calculate impacts $`ϵ^{}/ϵ`$in a significant manner. Our estimate of $`\mathrm{\Omega }_{\mathrm{IB}}`$ from the specific $`m_dm_u`$ effects we consider ranges from $`0.050.78`$; this range exceeds the central value, $`\mathrm{\Omega }_{\eta +\eta ^{}}=0.25\pm 0.05`$, used in earlier analyses and reflects a variation in $`ϵ^{}/ϵ`$of more than a factor of two. The presence of unknown low-energy constants implies that we lack a reliable way to calculate the effects we consider. Such limitations, however, underscore the need for a larger uncertainty in the Standard Model prediction of $`ϵ^{}/ϵ`$. The collaboration of G. Valencia is gratefully acknowledged. The work of S.G. is supported in part by the U.S. DOE under contract number DE-FG02-96ER40989.
warning/0006/hep-th0006048.html	ar5iv	text	# 1 Introduction ## 1 Introduction The duality between weakly coupled string theory on $`AdS_5\times S_5`$ and large $`N`$, N=4 super-Yang-Mills (SYM) theory is by now well established. The correspondence, that excitations in the string/SUGRA theory act as sources for operators in the field theory at the boundary, suggests that mass terms in the field theory will correspond to VEVs of scalar fields in the $`AdS`$ space. Several authors have explored SUGRA duals to theories which in the IR correspond to softly broken N=4 in this fashion -. Many of those results have been obtained by reduction to five dimensions where the problem can be studied in terms of a theory of scalars coupled to gravity, and we will pursue this approach further in this paper. The first field theories studied in this way had both IR and UV fixed points which in SUGRA corresponded to solutions of the relevant scalar equations of motion that flowed between two fixed points of the scalar potential in the radial direction in AdS ($`y`$) \- ,. Following the interpretations in the functional dependence on $`y`$ corresponds to RG flow in the field theory. Flows that at large $`y`$ approach the N=4 fixed point at the origin of the scalar potential but at small $`y`$ flow down the potential to a singularity have also been studied. They are interpreted as RG flows in a theory with no IR fixed point (at low scales the mass terms grow without bound). In this fashion SUGRA duals of strongly coupled relatives of N=2 and N=1 theories have been studied ,-. It must always be remembered that introducing mass term perturbations to the N=4 field theory does not decouple the massive fields from the strong dynamics because the N=4 theory is conformal and strongly coupled above the breaking scale. Supersymmetry hopefully ensures that the resulting theories do lie in the same universality class as their cousins with weak coupling at the breaking scale though. The singularities in the 5$`d`$ SUGRA approach to non-conformal field theories show that in the deep IR a fuller stringy picture of the dynamics is required and a number of papers have begun to study these descriptions (the singularities hide such objects as the enhancon singularity in N=2 and the Myers’ D-brane polarisation effect in N=1 ). Nevertheless as we will see the SUGRA description still contains much of the physics of the field theories. A set of flows have been identified in the SUGRA dual to N=2 SYM and some criteria must be used to distinguish between different flows to identify the physical ones. In it has been proposed that the potential evaluated along physical flows should remain bounded by the asymptotic value of the potential at the origin. Imposing this condition, as we discuss further below, picks out a set of flows that reasonably correspond to the RG flows in the field theory associated with different choices of position on the scalar moduli space. The extremum flows are the natural candidates to be identified with the singular points on the N=2 theory’s moduli space. In flows corresponding to N=1 SYM were investigated. With the introduction of an appropriate mass term (scalar) a single flow was found (up to the ability to rescale the RG parameter $`y`$ and neglecting the gaugino condensate). In this paper we want to explore the set of models that lie between these two extremes. The field theory of the N=2 model is known to have a quantum mechanical moduli space which at a general point has a $`U(1)^{N1}`$ gauge symmetry in the IR . The couplings of the U(1) gauge fields are determined by the periods of the Seiberg-Witten curve $$y^2=\underset{i=1}{\overset{N}{}}(x\varphi _i)^2\mathrm{\Lambda }^{2N},$$ (1) The singular points at $`tr\varphi _i^2=\mathrm{\Lambda }^2`$ correspond to places where the U(1) couplings diverge and there are massless dyons charged under the U(1)s. When the theory is perturbed by the addition of a mass term for the scalar multiplet the resulting potential pins the theory at the singular points. Holomorphy determines that as the mass terms are increased the vacuum must smoothly deform to the vacuum of N=1 SYM with the scalar VEV approaching zero. On the field theory side, the breaking N=4 $``$ N=2 $``$ $`N=1`$ has been studied in . In the analysis below we will present the SUGRA dual of this pinning on the moduli space and evolution to the N=1 SUGRA flow with no field theory scalar VEV. ## 2 Mass Terms and SUGRA Scalar Potentials Our starting point is the N=4 SYM theory which, in N=1 language, has three, adjoint chiral multiplets ($`\mathrm{\Phi }_1,\mathrm{\Phi }_2,\mathrm{\Phi }_3`$) in addition to the gauge multiplet. Previously deformations of the N=4 theory with an equal mass term for two of the gauge chiral multiplets and equal mass terms for all the three N=1 chiral multiplets have been considered. Here we will consider a generalisation with different masses for the chiral superfields. More precisely, we give equal mass, $`m`$ to two of the three superfields, and mass $`M`$ to the other one. In N=1 notations, this corresponds to $$W=m\underset{i=1,2}{}\mathrm{\Phi }_i^2+M\mathrm{\Phi }_3^2$$ (2) where $`m`$ and $`M`$ are complex. For $`m=M`$, this deformation corresponds to N=4 SYM softly broken to N=1 , while for $`M=0`$ one recovers N=4 SYM softly broken to N=2 . For $`M<<m`$, our N=1 solution should correspond to the soft breaking of N=2 to N=1 described above. As we increase $`M/m`$ we should smoothly return to the N=1 solutions. In the SUGRA description of the N=4 theory the mass terms (sources) in the field theory correspond to VEVs of the SUGRA fields. The precise fields have been identified from their symmetry properties under the conformal group (which indicates they are scalars) and the $`SU(4)_R`$ global symmetry of the N=4 theory. We will now identify the appropriate scalars and construct the five-dimensional supergravity solution corresponding to the deformation (2). The scalars of N=8 gauged supergravity are in the coset $`E_6/USp(8)`$ . The elements of the coset are $`27\times 27`$ matrices, $`U`$, transforming in the fundamental representation of $`E_6`$. In a unitary gauge, $`U`$ can be written as $`U=e^X,X=_a\lambda _aT_a`$, where $`T_a`$ are the generators of $`E_6`$ that do not belong to $`USp(8)`$. This matrix is parametrised by 42 real scalars, which are the physical scalars of the theory. They transform as the $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$, $`\mathrm{𝟐𝟎}^{}`$, and $`\mathrm{𝟏}_𝐜`$ of the gauge group $`SO(6)`$. The singlet is associated with the marginal deformation corresponding to a shift in the coupling constant of the N=4 theory. The mode in the $`\mathrm{𝟐𝟎}^{}`$ is associated with a symmetric traceless mass term for the scalars $$\text{Tr }(X_IX_J)\frac{1}{6}\delta _{IJ}Tr(X_LX_L),I,J=1,\mathrm{},6,$$ (3) while the $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$ correspond to the fermion mass term $$\text{Tr }\lambda _A\lambda _B+h.c.,A,B=1,\mathrm{}4.$$ (4) A generic supersymmetric mass term for the three chiral multiplets corresponds to turning on scalars both in the $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟎}}`$ and $`\mathrm{𝟐𝟎}^{}`$. Let us consider first the fermion mass term: $`m_{ij}\text{Tr }(\lambda _i\lambda _j)`$ with $`i=1,2,3`$. $`m_{ij}`$ is a complex, symmetric matrix that transforms as the $`\mathrm{𝟔}`$ of $`SU(3)SU(4)`$ ($`SO(6)`$ $`SU(4)`$). The corresponding supergravity mode appears in the decomposition of the $`\mathrm{𝟏𝟎}\mathrm{𝟏}+\mathrm{𝟔}+\mathrm{𝟑}`$ of $`SU(4)`$ under $`SU(3)\times U(1)`$. The singlet in this decomposition corresponds to the scalar $`\sigma `$ dual to the gaugino condensate of the N=1 SYM. In the analysis to follow we will set this field to zero in order to simplify the computation of the potential and reduce the dimension of the parameter space of flows. In principle it should be present and is presumably non-zero though, as we will see, neglecting it does not appear to disrupt the physical interpretation of flows in the remaining fields. In principle, a non-zero VEV for $`m_{ij}`$ will induce non-zero VEVs for other scalars as well, due to the existence of linear couplings of $`m`$ to other fields in the potential. A simple group theory analysis shows that the only couplings of the $`\mathrm{𝟔}`$ that give rise to a singlet of the symmetry group are of the form $`\mathrm{𝟖}\times (\overline{\mathrm{𝟔}}\times \mathrm{𝟔})^k`$ where $`k`$ is an integer number and the $`\mathrm{𝟖}`$ appears in the decomposition of the $`\mathrm{𝟐𝟎}^{}\mathrm{𝟖}+\mathrm{𝟔}+\overline{\mathrm{𝟔}}`$ of $`SU(4)`$ under $`SU(3)\times U(1)`$. The $`\mathrm{𝟖}\mathrm{𝟐𝟎}^{}`$ is then the only other mode that has to be considered, and all the remaining fields can be consistently set to zero. Notice that the $`\mathrm{𝟖}`$ corresponds exactly to the scalar mass term one would expect on the field theory side. Indeed by supersymmetry the mass term for the scalars is the square of the fermionic one, $`\mathrm{\Lambda }_{ij}=m_{ij}m_{jl}^{}`$ $``$ $`\mathrm{𝟏}+\mathrm{𝟖}\mathrm{𝟔}\times \overline{\mathrm{𝟔}}`$. The singlet, which amounts to the trace of the scalar mass terms, is associated with a massive stringy state and has no dual SUGRA description (the SUGRA only contains the massless string sector). An added complication is that a scalar VEV in the field theory has the same symmetry properties as the scalar mass term and is therefore also described by the $`\mathrm{𝟐𝟎}^{}`$. The deformation of eq.(2) corresponds to taking the matrices $`m_{ij}`$ and $`\mathrm{\Lambda }_{ij}`$ diagonal. In the complex basis of , they read<sup>1</sup><sup>1</sup>1The factors of $`\sqrt{2}`$ and $`\sqrt{6}`$ are required in supergravity in order for the fields to have canonical kinetic term.: $$(m_{ij})=\text{diag}(\frac{m}{\sqrt{2}},\frac{m}{\sqrt{2}},M)\text{and}(\mathrm{\Lambda }_{ij})=\frac{\rho }{\sqrt{6}}\text{diag}(1,1,2).$$ (5) If we fix $`m_{ij}`$ then for $`\mathrm{\Lambda }_{ij}`$ to correspond to the supersymmetric mass terms we must assume that the massive stringy mode diag(1,1,1) has also developed a VEV though this is not explicit in the SUGRA (the ability to describe these solutions must be implicitly present in the SUGRA). For a given choice of $`\rho `$ we assume the the stringy mode VEV brings the first two elements in line with supersymmetric requirements. We interpret the discrepancy in the third element from $`M^2`$ as a VEV for $`tr\varphi _3^2`$. Thus changing $`\rho `$, with fixed $`m_{ij}`$, allows us to explore the N=2 theory’s moduli space. Note that $`tr\varphi _3^2`$ is complex whilst $`\rho `$ is real, so we will only be able to explore the moduli space along a single radial direction. In the large $`N`$ limit though the $`Z_N`$ symmetry is restored to a U(1) so the radial direction is sufficient. The lengthy computation of the potential and kinetic terms is performed along the lines of . We refer to previous papers for an extensive description of these kinds of calculation. The 5-dimensional action for the scalars $`m,M`$ and $`\rho `$ is<sup>2</sup><sup>2</sup>2In the following we will always set the coupling constant $`g`$ equal to $`2`$, so that the scalar potential in the $`N=8`$ vacuum, where all scalars have zero VEV, is normalised as $`V(N=8)=3`$. $`L=\sqrt{g}\{{\displaystyle \frac{R}{4}}+{\displaystyle \frac{1}{2}}(m)^2+{\displaystyle \frac{1}{2}}(M)^2+{\displaystyle \frac{1}{2}}(\rho )^2+V\},`$ (6) where the potential $`V`$ is given by $`V`$ $`=`$ $`{\displaystyle \frac{1}{8}}e^{\frac{4\rho }{\sqrt{6}}}\left[5+\mathrm{cosh}(4M)4\mathrm{cosh}(2M)\right]e^{\frac{2\rho }{\sqrt{6}}}\mathrm{cosh}(\sqrt{2}m)\left[\mathrm{cosh}(2M)+1\right]+`$ (7) $`+{\displaystyle \frac{1}{16}}e^{\frac{8\rho }{\sqrt{6}}}\left[3+2\mathrm{cosh}(2\sqrt{2}m)+\mathrm{cosh}(4M)\right].`$ The potential above contains as special case the examples previously studied in the literature: for $`M=0`$ and $`M=\sqrt{2}m`$ it gives the N=2 and N=1 potentials, respectively, while for $`m=0`$ it reduces to the potential for the flow to the N=1 supersymmetric fixed point described in . In ref. the conditions for a supersymmetric flow were found. For a supersymmetric solution, the potential $`V`$ can be written in terms of a superpotential $`W`$ as $$V=\frac{1}{8}\underset{a=1}{\overset{n}{}}\left\|\frac{W}{\lambda _a}\right\|^2\frac{1}{3}\left\|W\right\|^2.$$ (8) In terms of the computation of the potential as described in $`W`$ is one of the eigenvalues of the tensor $`W_{ab}`$. Moreover, for such a supersymmetric solution, since the fermionic shifts vanish, the second order equations reduce to first order ones $`\dot{\lambda }_a`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{W}{\lambda _a}},`$ (9) $`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{1}{3}}W.`$ (10) Here $`\varphi `$ is the scale factor in the 5-dimensional metric $$ds^2=dy^2+e^{2\varphi (y)}dx^\mu dx_\mu ,\mu =0,1,2,3,$$ (11) where $`y`$ is the fifth coordinate of $`AdS_5`$, which we interpret as an energy scale : $`y\mathrm{}`$ corresponds to the UV regime while $`y\mathrm{}`$ to the IR. The dot in eq.(10) indicates derivative with respect to $`y`$. The $`W_{ab}`$ tensor of , has two different eigenvalues, both with degeneracy 2, that satisfy the condition in eq.(8), $`W_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\frac{4\rho }{\sqrt{6}}}\left(2\mathrm{cosh}(\sqrt{2}m)+\mathrm{cosh}(2M)1\right)e^{\frac{2\rho }{\sqrt{6}}}\left(\mathrm{cosh}(2M)+1\right)`$ (12) $`W_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{\frac{4\rho }{\sqrt{6}}}\left(2\mathrm{cosh}(\sqrt{2}m)\mathrm{cosh}(2M)+1\right)e^{\frac{2\rho }{\sqrt{6}}}\left(\mathrm{cosh}(2M)+1\right),`$ (13) indicating that there is more than one superpotential that generates the potential. We then expect to have two N=1 supersymmetric flows, with different field theoretic interpretations depending on the asymptotic behaviour of the fields for $`y\mathrm{}`$ . These are obtained by substituting into the equations of motion the linearised expressions for the eigenvalues (10) in the neighbour of the N=4 fixed point ($`m=M=\rho =0`$): $`W_13m^23M^22\rho ^2\frac{4}{\sqrt{6}}\rho m^2`$ and $`W_23m^2M^22\rho ^2\frac{4}{\sqrt{6}}\rho \left(m^22M^2\right)`$. Given those expressions, it is easy to check that while $`m`$ always scales as a mass deformation ($`me^y`$) , the behaviour of $`M`$ depends on the choice of $`W_i`$: a scalar VEV $`Me^{3y}`$ for $`W_1`$ or a mass deformation $`Me^y`$ for $`W_2`$. Thus flows of softly broken N=2 SYM should correspond to solutions of the equation of motion with $`W=W_2`$. Notice also that the same eigenvalue $`W_2`$ gives the equations of motions for the N=1 and N=2 cases mentioned above. ## 3 Properties of the solutions As discussed in the previous section, the SUGRA duals of RG flows of softly broken N=2 super Yang-Mills should be given by solutions of the equations $`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{1}{6}}e^{\frac{4\rho }{\sqrt{6}}}\left(2\mathrm{cosh}(\sqrt{2}m)\mathrm{cosh}(2M)+1\right)+{\displaystyle \frac{1}{3}}e^{\frac{2\rho }{\sqrt{6}}}\left(\mathrm{cosh}(2M)+1\right),`$ (14) $`\dot{\rho }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}\left[e^{\frac{4\rho }{\sqrt{6}}}\left(2\mathrm{cosh}(\sqrt{2}m)\mathrm{cosh}(2M)+1\right)e^{\frac{2\rho }{\sqrt{6}}}\left(\mathrm{cosh}(2M)+1\right)\right],`$ (15) $`\dot{m}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}e^{\frac{4\rho }{\sqrt{6}}}\mathrm{sinh}(\sqrt{2}m),`$ (16) $`\dot{M}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(e^{\frac{4\rho }{\sqrt{6}}}2e^{\frac{2\rho }{\sqrt{6}}}\right)\mathrm{sinh}(2M),`$ (17) with the boundary conditions that the scalars $`m,M`$ and $`\rho `$ vanish and $`\varphi y`$ for $`y\mathrm{}`$ . Contrary to the N=1 and N=2 cases, these equations do not seem to be analytically solvable, and we must rely upon numerical results. The solutions will depend on a certain number of parameters which correspond to the value of the fields at a given UV scale. The asymptotic behaviours of the fields for $`y\mathrm{}`$ are $`m`$ $``$ $`m_0e^y,`$ (18) $`M`$ $``$ $`M_0e^y,`$ (19) $`\rho `$ $``$ $`\sqrt{2/3}\left(m_0^2M_0^2\right)e^{2y}y+\rho _0e^{2y},`$ (20) $`\varphi `$ $``$ $`\varphi _0+y.`$ (21) The coefficients $`m_0`$ and $`M_0`$ corresponds to the UV field theory masses for the fermions. As expected, we can distinguish two contributions in the asymptotics of $`\rho `$: the first one corresponds to a mass term for the scalar components of $`\mathrm{\Phi }_{1,2}`$ and indeed it has the coefficient fixed by supersymmetry in terms of the fermion masses, the second term is associated to a VEV for the scalar of $`\mathrm{\Phi }_3`$ and the arbitrary coefficient represents the freedom in moving along the moduli space of the theory. In our numerical analysis we will not distinguish between these two different contributions, and we will just indicate the UV values of the fields as $`M(y_{UV})`$, etc… We first consider the N=2 ($`M=0`$) case for which an analytical solution has been given . This provides us with an interesting check of the numerical analysis, that we can then extend to the N=1 solution. ### 3.1 N=2 Flows Setting $`M=0`$ we find the N=2 flows described by the equations $`\dot{\rho }`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{6}}}\left(e^{4\rho /\sqrt{6}}\mathrm{cosh}\sqrt{2}me^{2\rho /\sqrt{6}}\right),`$ (22) $`\dot{m}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}e^{4\rho /\sqrt{6}}\mathrm{sinh}\sqrt{2}m,`$ (23) from which we find $$\frac{\rho }{m}=\frac{2}{\sqrt{3}\mathrm{sinh}\sqrt{2}m}\left(\mathrm{cosh}\sqrt{2}me^{6\rho /\sqrt{6}}\right).$$ (24) The solution to this equation was given in $$e^{\sqrt{6}\rho }=\mathrm{cosh}\sqrt{2}m+\mathrm{sinh}^2\sqrt{2}m\left[k+\mathrm{log}\left(\frac{\mathrm{sinh}m/\sqrt{2}}{\mathrm{cosh}m/\sqrt{2}}\right)\right].$$ (25) The constant $`k`$ that parametrises the solutions is related to the UV values of the fields by $`\left(2k+1\mathrm{log}2\right)\sqrt{6}=\rho _0/m_0^2\sqrt{2/3}\mathrm{log}m_0`$. We plot these flows over the superpotential in Fig. 1. In order to interpret these flows in terms of the N=2 gauge theory one requires some criteria for distinguishing between acceptable and unacceptable flows. In it has been proposed that the flows may be distinguished based on the behaviour of the potential and superpotential along the flows. In particular, extrapolating from finite temperature situations where there are requirements for consistent black hole solutions, the author has proposed that only flows where the scalar potential is bounded above all along the flow are physical. This seems intuitive since these are flows that begin at the origin of the $`\rho m`$ plane at large $`y`$ (that is they look like the N=4 theory in the UV) and then flow away from the origin, in directions where the potential falls, to large VEVs at small $`y`$ (the field theory IR). To discuss this criteria for the flows of Fig.1 it is sensible to formulate boundary conditions that are easily interpretable in the field theory. In the field theory it is natural to start at some UV scale with a fixed mass and look for RG flows that correspond to different positions on the moduli space of the theory. Thus in the SUGRA we should fix $`m`$ at some $`y=y_{UV}`$ and look at flows with varying $`\rho `$. We are therefore taking initial conditions on a vertical slice through Fig 1. We are then interested in the behaviour of the potential along the flow to lower $`y`$. We plot the evolution of the potential (and superpotential) for several values of $`y`$ along such flows using numerical solutions of the field equations in Fig.2. The behaviour is that all flows which start with $`\rho (y_{UV})`$ greater than the value on the $`k=0`$ “ridge flow” eventually meet a positive potential wall, and hence by the conditions in any flows with $`k>0`$ are to be considered unphysical. That the $`k=0`$ curve is the critical curve where this behaviour ends is clear from Fig 1 since only for the $`k>0`$ curves do the flows reach the singular behaviour of the superpotential which corresponds to that of the potential. It seems likely therefore that the flows with $`k0`$ correspond to the RG flows of different points of the field theory moduli space. There is some evidence to suggest this is true. For the curves $`k<0`$, as can be seen in Fig 1, the flows asymptotically approach $`\rho \mathrm{}`$ and $`m`$ constant. Evaluating the potential eq.(7) in this limit shows that the first term in the potential dominates and the potential becomes independent of $`m`$. Thus these flows asymptotically see the same potential suggesting a moduli space. That the different flows approach the asymptotic form of the potential at different rates with respect to $`y`$ is perhaps an indication that in the field theory different points on the moduli space have a different scale $`tr\varphi _3^2`$ at which the gauge symmetry is broken, $`SU(N)U(1)^{N1}`$. The singular point on the N=2 moduli space, where the U(1) couplings diverge, should naturally be equated with a special or extremal flow. The obvious such flow is the case $`k=0`$ which follows the crest of the ridge on the superpotential . These identifications can only be tentative based on the analysis so far. In further evidence was provided by lifting the 5$`d`$ SUGRA solution to a 10$`d`$ SUGRA solution. This allows the authors to evaluate the gauge coupling (corresponding to the VEVs of the singlet scalars amongst the 42 scalars in the 5$`d`$ SUGRA theory - in the 5$`d`$ theory they do not enter the potential and so their VEVs can not be determined). They find that the coupling diverges on the $`k=0`$ flow but runs to a constant elsewhere which seems in accord with the field theory although the functional dependence of the coupling on the moduli space has not been reproduced. We next move on to consider breaking the N=2 theory to the N=1 theory by the inclusion of the mass term $`M`$. In the field theory we expect to be pinned at the singular point and indeed we will see that the SUGRA flows pick out the $`k=0`$ flow providing further evidence for the above identification of the flows with the N=2 moduli space. ### 3.2 N=1 Flows Flows with non-zero $`M`$ are expected to correspond to N=1 super-Yang-Mills theories. We begin by looking at theories which are N=2 SYM plus a small mass $`M`$ in the UV that breaks the supersymmetry to N=1. We again, in the SUGRA, fix $`m`$ and $`M`$ at some $`y=y_{UV}`$ and look at flows with varying $`\rho (y_{UV})`$. We are then interested in the behaviour of the potential along the flow to lower $`y`$. In Fig.3 we plot the evolution of the potential (and superpotential) for several values of $`y`$ along such flows using numerical solutions of the field equations with $`M(y_{UV})=0.1m(y_{UV})`$. From Fig.3 it is apparent that as we proceed towards the IR all the flows except the flow along the ridge of the superpotential eventually meet a positive “barrier” in the potential. Thus only the single ridge flow is allowed by the criteria for distinguishing flows discussed above. This matches our expectations from field theory where we expected the theory to become pinned at the singular point of the N=2 theory previously identified with that ridge flow. The existence of a ridge flow seems to be confirmed by the analysis of the numerical solutions of the equations of motion. One can indeed distinguish two sets of solutions with radically different behaviours. The change of behaviour seems to take place for a negative value of the $`\rho (y_{UV})`$, which should correspond to the position of the ridge flow. This is consistent with what shown in Fig 3. In Fig. 4 we plot $`M,m`$ and $`\rho `$ as functions of $`y`$ for two values of $`\rho (y_{UV})`$ on either sides of the ridge flow. In analogy with the N=2 theory case, we will call the ridge flow $`k=0`$, while $`k<0`$ and $`k>0`$ will indicate the flows on the left and on the right of it, respectively. Notice however that here the parameter $`k`$ is not related with any precise form of the solutions. Unfortunately it is not very easy to reproduce the behaviour of the ridge flow. Numerically, it is very difficult to pick out the precise value of $`\rho (y_{UV})`$ corresponding to the ridge flow, and analytically, it is not easy to solve the equations of motion even in the IR limit. It is even hard to guess the correct asymptotic behaviour of the fields. For example, one would have hoped to recover a logarithmic behaviour like in the N=1 or N=2 case, but this possibility seems to be ruled out. This is confirmed by the plots in Fig. 5 below, where the behaviour of the ratios of the various fields as a function of $`y`$ are shown. The existence of, and distinction between, the ridge flow and the two classes of flows to either side are clear, but note that close to the ridge flow there is no linear relation between the fields. From Fig. 4, one can see that our solutions are all singular in the IR. It is then natural to ask whether they correspond to physical or unphysical flows. The expectation is that only the ridge flow should be physical. Unfortunately we can explicitly check only the case $`k<0`$, where we can extract information from the IR asymptotic behaviours $`m`$ $``$ $`\text{const},`$ (26) $`M`$ $``$ $`{\displaystyle \frac{3}{7}}\mathrm{log}\|yy_0\|+\text{const},`$ (27) $`\rho `$ $``$ $`{\displaystyle \frac{\sqrt{6}}{14}}\mathrm{log}\|yy_0\|+\text{const},`$ (28) $`\varphi `$ $``$ $`{\displaystyle \frac{1}{7}}\mathrm{log}\|yy_0\|+\text{const}.`$ (29) Notice in particular the logarithmic divergence of the field $`\varphi `$: $`\varphi A\mathrm{log}\|yy_0\|`$. Indeed for logarithmically divergent flows, the criterion proposed in to select physically sensible solutions seems to pick up flows with $`A1/4`$ . This rules out the $`k<0`$ flows for which $`A=1/7`$. For the ridge and the $`k>0`$ flows we can rely only on the numerical behaviour of the potential and the superpotential, that, as observed above, seem to allow only for one physical flow, the ridge. Finally let us consider the situation where the two masses $`m`$ and $`M`$ are of the same order. In field theory we expect to recover pure N=1 SYM. In particular the vacuum state should evolve into the N=1 solution of . From the supergravity point of view the N=1 and the N=2 theories can be distinguished by the scalars one has to turn on. As discussed in , in the N=1 solution the scalar $`\rho `$ is set to zero. In this case, the supergravity mode corresponding to the mass term for the scalar fields corresponds to the stringy mode only and it is therefore not present in the 5$`d`$ supergravity Lagrangian. Thus we expect that as we increase $`M(y_{UV})`$ the ridge line flow should smoothly move to the value $`\rho =0`$. This is indeed the case, as can be seen in Fig. 6 where we have plotted the potential and superpotential as a function of the UV values of $`\rho `$, in the deep IR (small $`y`$), for three different values of the ratio $`M(y_{UV})/m(y_{UV})`$. $`M(y_{UV})=m(y_{UV})/\sqrt{2}`$ corresponds to pure N=1 SYM and the superpotential ridge has indeed moved to $`\rho =0`$. ## 4 Discussion We have calculated the 5$`d`$ SUGRA scalar potential corresponding to fields dual to scalar and fermion masses and a scalar VEV for the lightest chiral multiplet in N=4 SYM softly broken in a cascade to N=2 and then N=1 SYM. We have found numerical solutions of the resulting SUGRA equations of motion and used the criteria that the potential must fall all along the flow, as suggested in , to distinguish physical flows. Interpreting these flows in terms of the dual field theory leads to a pleasing picture of the transition between the N=2 flows of and the N=1 flows of . The N=2 theory has a moduli space and thus there is a class of flows in SUGRA which are physically acceptable. The extremal curve has been identified with the singular point on the moduli space. When a small mass is introduced to break the theory to N=1 SYM all these flows except the extremal flow have their potential lifted in the IR which, according to the criteria above, indicates they are unphysical. We thus see the unique vacuum of the N=1 theory emerging from the N=2 theory and, as expected from the field theory, the N=1 theory is pinned at the singular point on the N=2 moduli space. As the two perturbing masses (the one that sets the breaking N=4 to N=2 and the one that further breaks to N=1) become of the same order, leaving N=1 SYM in the IR, the extremal flow is seen to smoothly move to the origin of the N=2 theory’s moduli space, again in agreement with field theory expectations. A number of issues remain for future investigation. The N=1 theory is expected to have a gaugino condensate, represented in SUGRA by a further scalar we have neglected in our discussion. Although the potential is easy to compute, including this scalar would raise the dimension of the parameter space of flows making analysis harder, so we prefer to leave it for future analysis. From the field theory we also expect to see monopole condensation and confinement (in the field theory of the softly broken N=2 theory there are $`N^21`$ different string tensions one might hope to see a SUGRA dual of ). The solutions we present have an IR singularity so we could hope that the ridge flow solution exhibits confinement. However we have no means to check the behaviour without explicit solutions. The usual methods one can apply in supergravity to check confinement, such as the computation of a Wilson loop or the the evaluation of the spectrum of scalars , all require at least the knowledge of the IR asymptotic behaviour of the solution. In this respect it may be profitable to study the lift of the 5$`d`$ SUGRA to the full 10$`d`$ theory where explicit solutions may be more readily obtainable. The 10$`d`$ solutions would also allow study of the running of the gauge coupling which would provide further checks of our interpretation. AcknowledgementsM.P. would like to thank A. Zaffaroni and G. Salam for very useful discussions and comments. M.P is partially supported by INFN and MURST, and by the European Commission TMR program ERBFMRX-CT96-0045, wherein she is associated to Imperial College, London, and by the PPARC SPG grant 613. N.E is grateful for the support of a PPARC Advanced Fellowship.
warning/0006/math0006115.html	ar5iv	text	# Geometric Interpretations of Quandle Homology ## 1 Introduction A quandle is a set with a self-distributive binary operation (defined below) whose definition was motivated from knot theory. A (co)homology theory was defined in for quandles, which is a modification of rack (co)homology defined in . State-sum invariants using quandle cocycles as weights are defined and computed for important families of classical knots and knotted surfaces . Quandle homomorphisms and virtual knots are applied to this homology theory . The invariants were applied to study knots, for example, in detecting non-invertible knotted surfaces . On the other hand, knot diagrams colored by quandles can be used to study quandle homology groups. This view point was developed in for rack homology and homotopy. In this paper, we generalize this approach to quandle homology theory. For this purpose, we generalize abstract knot diagrams , and represent cycles and boundaries geometrically by colored abstract knot diagrams. Equivalence moves are given in terms of colored knot diagrams. The paper is organized as follows. In Section 2 the definition of quandle homology is recalled. Abstract knot diagrams are generalized in Section 3. Colorings of generalized diagrams are defined (Section 4), and used to represent cycles of quandle homology (Section 5). Equivalence moves for colored diagrams are given in Section 6. In Section 7, the fundamental quandles for abstract knots are discussed. The boundary homomorphisms in an exact sequence are shown to be trivial (Section 8). Examples appear throughout the paper. ## 2 Definitions of Quandle (Co)Homology ###### 2.1 Definition. A quandle, $`X`$, is a set with a binary operation $``$ such that (I. idempotency) for any $`aX`$, $`aa=a`$, (II. right-invertibility) for any $`a,bX`$, there is a unique $`cX`$ such that $`a=cb`$, and (III. self-distributivity) for any $`a,b,cX`$, we have $`(ab)c=(ac)(bc)`$. A rack is a set with a binary operation that satisfies (II) and (III). Racks and quandles have been studied in, for example, ,,,, and . A map $`f:XY`$ between two quandles (resp. racks) $`X,Y`$ is called a quandle (resp. rack) homomorphism if $`f(ab)=f(a)f(b)`$ for any $`a,bX`$. A (quandle or rack) homomorphism is a (quandle or rack) isomorphism if it is bijective. An isomorphism between the same quandle (or rack) is an automorphism. ###### 2.2 Examples. Any set $`X`$ with the operation $`xy=x`$ for any $`x,yX`$ is a quandle called the trivial quandle. The trivial quandle of $`n`$ elements is denoted by $`T_n`$. Any group $`G`$ is a quandle by conjugation as operation: $`ab=b^1ab`$ for $`a,bG`$. Any subset of $`G`$ that is closed under conjugation is also a quandle. For example, the set, $`QS(5)`$, of non-identity elements of the symmetric group on three letters is a quandle. Similarly the set $$QS(6)=\{(1234)=a,(1243)=b,(1324)=c,(1342)=B,(1423)=C,(1432)=A\}$$ of 4-cycles in the symmetric group on four letters is a quandle. Let $`n`$ be a positive integer. For elements $`i,j\{0,1,\mathrm{},n1\}`$, define $`ij=2ji`$ where the sum on the right is reduced mod $`n`$. Then $``$ defines a quandle structure called the dihedral quandle, $`R_n`$. This set can be identified with the set of reflections of a regular $`n`$-gon with conjugation as the quandle operation. We also represent the elements of $`R_3`$ by $`\alpha ,\beta ,`$ and $`\gamma `$, where the quandle multiplication is given by $`xy=z`$ where $`zx,y`$ when $`xy`$ and $`xx=x`$, for $`x,y,z\{\alpha ,\beta ,\gamma \}`$. As an exercise the reader may check that there is a quandle homomorphism $`p:QS(6)R_3`$ given by $`p(a)=p(A)=\alpha `$, $`p(b)=p(B)=\beta `$, and $`p(c)=p(C)=\gamma `$. Any $`\mathrm{\Lambda }=𝐙[T,T^1]`$-module $`M`$ is a quandle with $`ab=Ta+(1T)b`$, $`a,bM`$, called an Alexander quandle. Furthermore for a positive integer $`n`$, a mod-$`n`$ Alexander quandle $`𝐙_n[T,T^1]/(h(T))`$ is a quandle for a Laurent polynomial $`h(T)`$. The mod-$`n`$ Alexander quandle is finite if the coefficients of the highest and lowest degree terms of $`h`$ are $`\pm 1`$. See , , , or for further examples. ###### 2.3 Remark. Let $`X`$ denote a quandle. From Axiom II, each element $`bX`$ defines a bijection $`S(b):XX`$ with $`aS(b)=ab`$. The bijection is an automorphism by Axiom III. For a word $`w=b_1^{ϵ_1}\mathrm{}b_n^{ϵ_n}`$ where $`b_1,\mathrm{},b_nX;ϵ_1,\mathrm{},ϵ_n\{\pm 1\}`$, we define $`aw=aS(w)`$ by $`aS(b_1)^{ϵ_1}\mathrm{}S(b_n)^{ϵ_n}`$. An automorphism of $`X`$ is called an inner-automorphism of $`X`$ if it is $`S(w)`$ for a word $`w`$. (The notation $`S(b)`$ follows Joyce’s paper and $`aw`$ ($`=a^w`$) follows Fenn-Rourke .) Let $`C_n^\mathrm{R}(X)`$ be the free abelian group generated by $`n`$-tuples $`(x_1,\mathrm{},x_n)`$ of elements of a quandle $`X`$. Define a homomorphism $`_n:C_n^\mathrm{R}(X)C_{n1}^\mathrm{R}(X)`$ by $`_n(x_1,x_2,\mathrm{},x_n)`$ (1) $`={\displaystyle \underset{i=2}{\overset{n}{}}}(1)^i[(x_1,x_2,\mathrm{},x_{i1},x_{i+1},\mathrm{},x_n)`$ $`(x_1x_i,x_2x_i,\mathrm{},x_{i1}x_i,x_{i+1},\mathrm{},x_n)]`$ for $`n2`$ and $`_n=0`$ for $`n1`$. Then $`C_{}^\mathrm{R}(X)=\{C_n^\mathrm{R}(X),_n\}`$ is a chain complex. Let $`C_n^\mathrm{D}(X)`$ be the subset of $`C_n^\mathrm{R}(X)`$ generated by $`n`$-tuples $`(x_1,\mathrm{},x_n)`$ with $`x_i=x_{i+1}`$ for some $`i\{1,\mathrm{},n1\}`$ if $`n2`$; otherwise let $`C_n^\mathrm{D}(X)=0`$. If $`X`$ is a quandle, then $`_n(C_n^\mathrm{D}(X))C_{n1}^\mathrm{D}(X)`$ and $`C_{}^\mathrm{D}(X)=\{C_n^\mathrm{D}(X),_n\}`$ is a sub-complex of $`C_{}^\mathrm{R}(X)`$. Put $`C_n^\mathrm{Q}(X)=C_n^\mathrm{R}(X)/C_n^\mathrm{D}(X)`$ and $`C_{}^\mathrm{Q}(X)=\{C_n^\mathrm{Q}(X),_n^{}\}`$, where $`_n^{}`$ is the induced homomorphism. Henceforth, all boundary maps will be denoted by $`_n`$. For an abelian group $`G`$, define the chain and cochain complexes $`C_{}^\mathrm{W}(X;G)=C_{}^\mathrm{W}(X)G,`$ $`=\mathrm{id};`$ (2) $`C_\mathrm{W}^{}(X;G)=\mathrm{Hom}(C_{}^\mathrm{W}(X),G),`$ $`\delta =\mathrm{Hom}(,\mathrm{id})`$ (3) in the usual way, where $`\mathrm{W}`$ $`=\mathrm{D}`$, $`\mathrm{R}`$, $`\mathrm{Q}`$. ###### 2.4 Definition. The $`n`$th rack homology group and the $`n`$th rack cohomology group of a rack/quandle $`X`$ with coefficient group $`G`$ are $`H_n^\mathrm{R}(X;G)=H_n(C_{}^\mathrm{R}(X;G)),H_\mathrm{R}^n(X;G)=H^n(C_\mathrm{R}^{}(X;G)).`$ (4) The $`n`$th degeneration homology group and the $`n`$th degeneration cohomology group of a quandle $`X`$ with coefficient group $`G`$ are $`H_n^\mathrm{D}(X;G)=H_n(C_{}^\mathrm{D}(X;G)),H_\mathrm{D}^n(X;G)=H^n(C_\mathrm{D}^{}(X;G)).`$ (5) The $`n`$th quandle homology group and the $`n`$th quandle cohomology group of a quandle $`X`$ with coefficient group $`G`$ are $`H_n^\mathrm{Q}(X;A)=H_n(C_{}^\mathrm{Q}(X;G)),H_\mathrm{Q}^n(X;A)=H^n(C_\mathrm{Q}^{}(X;G)).`$ (6) The homology group of a rack in the sense of is $`H_n^\mathrm{R}(X;G)`$ and the cohomology of a quandle used in is $`H_\mathrm{Q}^n(X;A)`$. Refer to , , , for some calculations and applications of the rack homology groups, and to , for those of quandle cohomology groups. The cycle and boundary groups (resp. cocycle and coboundary groups) are denoted by $`Z_n^\mathrm{W}(X;G)`$ and $`B_n^\mathrm{W}(X;G)`$ (resp. $`Z_\mathrm{W}^n(X;G)`$ and $`B_\mathrm{W}^n(X;G)`$), so that $$H_n^\mathrm{W}(X;G)=Z_n^\mathrm{W}(X;G)/B_n^\mathrm{W}(X;G),H_\mathrm{W}^n(X;G)=Z_\mathrm{W}^n(X;G)/B_\mathrm{W}^n(X;G)$$ where $`\mathrm{W}`$ is one of $`\mathrm{D}`$, $`\mathrm{R}`$, $`\mathrm{Q}`$. We will omit the coefficient group $`G`$ if $`G=𝐙`$ as usual. Here we are almost exclusively interested in quandle homology or cohomology. ###### 2.5 Example. For $`QS(6)`$ (defined in 2.2), we have computed using Mathematica that $`H_3^\mathrm{Q}(QS(6);𝐙)=𝐙_{24}.`$ Similarly we have $`H_3^\mathrm{Q}(R_3;𝐙)=𝐙_3.`$ We will illustrate in Example 7.5 that the homomorphism $`p:QS(6)R_3`$ induces a surjection $`𝐙_{24}𝐙_3`$. Many other calculations are found in . ## 3 Generalized Knot Diagrams In a general framework for colored diagrams was sketched. Such diagrams are useful in representing rack homology classes and homotopy classes of maps into the classifying space of a rack. Here we review, amplify, and generalize some of their constructions to the case of quandle homology classes. Abstract knot diagrams are defined in and the relation to virtual knots is established in . We extend their definition to include arcs and surfaces with boundary. First we generalize crossings of classical knot diagram to higher dimensions. In the top row and the bottom three rows of Fig. 1, $`k`$-crossing $`1`$-diagrams are indicated for $`k=0,1,2`$. In Fig. 2 $`k`$-crossing $`2`$-diagrams are indicated for $`k=0,1,2,`$ and $`3`$. We define such diagrams in general. ###### 3.1 Definition. Fix a non-negative integer $`n`$. Let $`E`$ be the $`(n+1)`$-disk $`[1,1]^{n+1}=\{(x_1,\mathrm{},x_{n+1})\|1x_i1,i=1,\mathrm{},n+1\}`$, and for $`j\{1,\mathrm{},n+1\}`$, let $`E_j`$ be the $`n`$-disk in $`E`$ determined by $`x_j=0`$. A $`0`$-crossing $`n`$-diagram is the disk $`E`$. A $`1`$-crossing $`n`$-diagram (or simply 1-crossing diagram) $`\mathrm{\Sigma }_1`$ is the pair of $`E`$ and the $`n`$-disk $`E_1`$. Let $`k`$ be an integer in the range $`\{1,\mathrm{},n+1\}`$. A $`k`$-crossing $`n`$-diagram (or simply $`k`$-crossing diagram) $`\mathrm{\Sigma }_k`$ is the pair of $`E`$ and the union of the disks $`E_1\mathrm{}E_k`$ with crossing information that is indicated by removing from $`E_j`$ (for $`j=2,\mathrm{},k`$) an open $`\epsilon `$-neighborhood of $`(E_1\mathrm{}E_{j1})E_j`$. For example, when $`n=1`$ a $`2`$-crossing $`1`$-diagram is the standard depiction of a classical knot crossing. When $`n=2`$, a $`2`$-crossing $`2`$-diagram is the broken surface depiction of pair of planes in $`𝐑^4`$ that project to intersect along a double curve as depicted in Fig. 2 third entry from top. A $`3`$-crossing $`2`$-diagram is a broken diagram of a triple point as depicted in Fig. 2 bottom two entries. The removal of the $`\epsilon `$-neighborhood is a convention for indicating relative height in a projection from $`(n+2)`$-dimensional space. Thus in $`E\times [0,1]`$, we re-embed $`E_k`$ as $`\stackrel{~}{E_k}=E_k\times \{\frac{n+2k}{n+2}\}`$. Then the lift, $`\stackrel{~}{E_k}`$, projects to $`E_k`$ under the projection $`E\times [0,1]E`$. The relative position in the $`[0,1]`$ factor then is determined by the number of regions on the sheet $`E_k`$. Such diagrams, $`(E,E_j)`$, are oriented. For example, when $`n=1`$, orientations of arcs are indicated by arrows along the arc. The orientation of $`E`$ is indicated by a short arrow normal to the arc $`E_j`$ where tangent followed by normal is the standard orientation of the disk. The normal arrow may be omitted from the figure when the orientation of the disk $`E`$ is that of the plane of the paper. When $`n=2`$, we assume that the disks $`E_j`$ are oriented in a counterclockwise fashion and only indicate the normal arrow. In higher dimensions, signs of $`(n+1)`$-crossing $`n`$-diagrams are determined by normal vectors as follows. Let the normal to $`E_j`$ be denoted by $`\nu _j`$. Then the $`(n+1)`$-tuple $`(\nu _1,\mathrm{},\nu _{n+1})`$ determines an orientation of the disk $`E`$. The $`(n+1)`$-crossing point is positive if and only if this orientation agrees with the standard orientation. Thus if the normal to $`E_j`$ is the standard basis vector $`e_j=(0,\mathrm{},1,\mathrm{},0)`$ (where the non-zero entry is in the $`j`$th coordinate), then the $`(n+1)`$-crossing point is positive. ###### 3.2 Definition. There are $`2^k1=1+2+\mathrm{}+2^{k1}`$ components in $`\mathrm{\Sigma }_k`$. We call them the $`n`$-regions of $`\mathrm{\Sigma }_k`$. The disk $`E_j`$ is called the level $`j`$ sheet of $`\mathrm{\Sigma }_k`$. The $`k`$-crossing point set (or simply $`k`$-crossing) in $`E`$ is the intersection $`E_1E_2\mathrm{}E_k`$, which is a disk of dimension $`n+1k`$ (or of codimension $`k`$). The complement of $`E_1\mathrm{}E_k`$ in $`E`$ has $`2^k`$ components called $`(n+1)`$-regions. ###### 3.3 Remark. Observe that in the level $`j`$ sheet at a $`k`$-crossing we have a $`(j1)`$-crossing diagram; this diagram is the orthogonal projection of the previous levels onto the $`j`$th level. Next we introduce endpoint, branch point, and hem diagrams. These exceptional diagrams will be useful to us in describing degenerate chains in quandle homology. ###### 3.4 Defintion. An endpoint diagram is a $`2`$-disk with an arc embedded; one end of the arc is in the interior of the disk, the other end is on the boundary, and the arc intersects the boundary transversely. In our illustrations, we indicate the interior point as a fat vertex. See Fig. 1 the second and third pictures from top. ###### 3.5 Definition. A branch point diagram consists of a neighborhood of a branch point (also called Whitney umbrella, a generic singularity of surface maps) in a $`3`$-ball in which over and under crossing information is indicated by removing a thin open triangle along the double point arc as depicted in Figure 3, the top and second top pictures. The vertex of the triangle is at the branch point which is in the interior of the 3-ball. The boundary of the branch point neighborhood on the boundary of the $`3`$-ball is the diagram of the unknot with one crossing. ###### 3.6 Definition. A hemmed $`1`$-crossing diagram is a $`2`$-dimensional half-disk embedded in a $`3`$-ball as depicted in Fig. 3 the second bottom pictures. A part of the boundary of the disk is in the interior of the $`3`$-ball, called the hem. ###### 3.7 Definition. A hemmed $`2`$-crossing diagram consists of two disks immersed (with crossing information) in the 3-ball with the boundary of one of these (the under-sheet) contained in the interior of the $`3`$-ball. Crossing information is depicted in the bottom of Fig. 3 by removing a thin rectangular neighborhood of the double point arc in the under-sheet — the horizontal sheet in the illustration. The under-sheet of a hem $`2`$-crossing diagram is the sheet that has an interior hem boundary (depicted by a thick line segment). It is important to note that at a hemmed $`2`$-crossing diagram it is the hemmed surface that is the under-sheet. An endpoint diagram, a branch point diagram, or a hemmed $`1`$\- or $`2`$-crossing diagram is also called an exceptional diagram. In Figs.1, 2, and 3, the collection of endpoint diagrams, branch point diagrams, crossing diagrams in low dimensions, and the handles that they correspond to are indicated. ###### 3.8 Definition. A generalized abstract $`0`$-knot diagram consists of a collection of vertices (points) embedded in a closed $`1`$-manifold. Let $`n=1`$ or $`n=2`$. A generalized abstract $`n`$-knot diagram is the image of a generic map $`f:MN`$ with crossing information indicated where: 1. $`f(M)N=\mathrm{}`$; 2. Local crossing information indicated as follows: Each point of $`N`$ has a neighborhood $`B`$, such that $`Bf(M)`$ is a $`k`$-crossing $`n`$-diagram ($`k=0,\mathrm{},n+1`$), or an exceptional diagram. The image $`f(M)`$ is called the universe of the diagram. If the manifold $`M`$ is closed, then the diagram is said to be closed. Thus when $`n=2`$, a closed diagram may have branch points but does not have hems. Abstract diagrams will often be denoted by $`K=[f:MN]`$. This is a slight abuse of notation since we are considering the image $`f(M)N`$ with crossing information. Abstract diagrams are depicted in Figs. 4, 5, 7, and 11. ###### 3.9 Remark. In abstract link diagrams were introduced in which the image $`f(M)`$ is a deformation retract of the ambient space $`N`$ (see also ). Here we do not need to use this strong condition. Moreover, a generalized abstract diagram (of a closed manifold $`M`$) in the current sense gives rise to an abstract diagram in the sense of by taking a regular neighborhood of the image. Figure 5 is a strict generalization of abstract diagrams defined in , and Fig. 4 is an abstract diagram in any sense. Throughout the sequel, generalized abstract diagrams are called abstract diagrams for simplicity. Abstract diagrams can be constructed via handle theoretic techniques. In Figs. 1, 2, and 3 crossing diagrams and exceptional diagrams are classified as handles, and attaching regions are indicated. ###### 3.10 Remark. Let $`K`$ be an abstract $`1`$-knot diagram in $`N`$. Then $`K`$ determines an embedded arc $`\widehat{K}`$ in $`N\times [0,1]`$ by lifting each over-arc to a higher level (in the $`[0,1]`$ factor) than the under-arc at each double point of $`K`$, as is typically done in classical knot theory. In the present convention endpoints are embedded in $`N\times \{0\}`$. ###### 3.11 Example. An abstract diagram of a closed curve is indicated in Fig. 4. The $`1`$-manifold is $`S^1`$, and the manifold $`N`$ is a thrice punctured torus. There are three $`2`$-crossing $`1`$-diagrams as its subdiagrams, the top two are positive and the bottom one is negative. The $`0`$-handles of $`N`$ are these $`2`$-crossings depicted as lightly shaded squares. The unshaded bands are $`1`$-handles that are $`1`$-crossing $`1`$-diagrams. ###### 3.12 Example. An abstract diagram of an arc in a planar surface is indicated in Fig. 5. The diagram has two positive crossing points and two endpoints. ###### 3.13 Remark. An abstract closed $`2`$-knot diagram is constructed from copies of $`k`$-crossing $`2`$-diagrams ($`k=1,2,3`$) and branch point diagrams as follows. The crossing diagrams are identified along their boundaries so that the arcs of double points either form closed loops or end at $`3`$-crossing diagrams or branch point diagrams. Thus in the boundary of $`N`$, there is, at most, a collection of simple closed curves. These then are capped-off by adding $`2`$-handles of the form $`1`$-crossing $`2`$-diagrams. Moreover, when crossing diagrams agree at points of their boundaries the broken and unbroken sheets match broken and unbroken sheets, respectively (see Fig. 6, Fig. 2 bottom two pictures, and Fig. 19). It is also required that orientations match when crossing diagrams are glued. ###### 3.14 Remark. An abstract $`2`$-knot diagram represents a surface $`M`$ embedded in $`N\times [0,1]`$ by filling in the surface along the broken sheets in the $`[0,1]`$ direction (see for details on filling in broken surface diagrams). In this way, the surface represented by the diagram is embedded in $`N\times [0,1]`$. ###### 3.15 Example. In Fig. 7, a part of a closed abstract $`2`$-knot diagram is indicated (the relation to the arc diagram in the lower right will be explained subsequently). The diagram has two $`3`$-crossing points (triple points) and two branch points. The arcs of double points that end in the front of the diagram are to be glued to the back of the diagram via two $`1`$-handles each of which is a $`2`$-crossing diagram. In this way, the $`3`$-manifold that contains the diagram is a handle-body, and the boundary contains 5 simple closed curves. These are attaching regions for $`2`$-handles that are $`1`$-crossing diagrams. ###### 3.16 Remark. Abstract diagrams can be defined similarly in higher dimensions using $`k`$-crossing $`n`$-diagrams, as in . However, the branch point set is more subtle in higher dimensions. ## 4 Coloring Abstract Diagrams Abstract diagrams, when colored by quandles, represent homology classes. Colored abstract $`1`$\- and $`2`$-knot diagrams are of particular interest. ###### 4.1 Definition. For $`n=0,1,2`$ and $`kn+1`$, a coloring of a $`k`$-crossing $`n`$-diagram is defined as an assignment of quandle elements to the points, arcs or regions of the diagram such that the following color condition at the crossing is satisfied. The condition is depicted in Fig. 8. In this figure, an under-arc (or lower sheet of a $`2`$-dimensional region) is colored with a quandle element, say $`q_1`$. The arc (or region) that is so colored is the arc away from which the normal to the over-arc (or region) points. The over-arc is colored with $`q_2`$. And the remaining under arc (or region) is colored with $`q_1q_2`$. In case $`n=0`$ only the point is colored with $`q_1`$. In case the crossing point is a triple point, then the middle sheet has two colors ($`q_2`$ and $`q_2q_3`$) as described above, and the lowest sheet has four colors ($`q_1`$, $`q_1q_3`$, $`q_1q_2`$, and $`(q_1q_2)q_3`$). A branch point diagram is colored by assigning a single quandle element to the 2-dimensional region of the diagram. Let $`n=1,2`$. A (quandle) coloring of a closed abstract $`n`$-knot diagram is an assignment of quandle elements to the connected regions of the diagram such that the restriction to each branch point (when $`n=2`$) or $`k`$-crossing diagram is a coloring. More precisely, let $``$ be the set of arcs (for $`n=1`$) or regions (for $`n=2`$) of an abstract $`n`$-knot diagram ($`n=1,2`$). A coloring is a map $`𝒞:X`$, where $`X`$ is a quandle, such that $`\{𝒞(r)\|r\}`$ satisfies the above mentioned conditions at each $`k`$-crossing $`n`$-diagram in the given abstract $`n`$-knot diagram. ###### 4.2 Examples. Figure 4 indicates a coloring of an abstract $`1`$-knot diagram by $`R_3`$. Figure 7 indicates a coloring of an abstract $`2`$-knot diagram by $`R_3`$. In Fig. 9 shadow colorings of $`k`$-crossing $`n`$-diagrams for all $`n=0,1,2,`$ and $`1kn+1`$ are depicted. We define these case by case. ###### 4.3 Definition. For any $`n`$, a shadow coloring of a $`0`$-crossing diagram is an assignment of a quandle element to the $`(n+1)`$-dimensional region of the diagram. A shadow coloring of a $`1`$-crossing diagram is an assignment of three quandle elements to the constituents of the diagram, as follows. The $`(n+1)`$-dimensional region away from which the normal of the $`1`$-crossing recieves a quandle element $`q_0`$, the $`n`$-dimensional region ($`1`$-crossing) receives the color $`q_1`$, and the $`(n+1)`$-dimensional region into which the normal points receives color $`q_0q_1`$. In general, a shadow coloring of a $`k`$-crossing $`n`$-diagram is an assigment of colors to the $`n`$ and $`(n+1)`$-dimensional regions that satisfies the following conditions. Along the $`n`$-dimensional regions of the diagram, it is a coloring of the $`k`$-crossing diagram. Each point in an $`(n+1)`$-dimensional region of the diagram is assigned a quandle element, and if two such regions are separated by an $`n`$-dimensional region, then any point in the $`n`$-dimensional region has a neighborhood homeomorphic to a $`1`$-crossing $`n`$-diagram. In this case, the coloring on the $`(n+1)`$-dimensional regions satisfies the condition that the restriction to such a neighborhood is a shadow coloring of the $`1`$-crossing diagram. A shadow coloring of an endpoint diagram is the assignment of the same color to the arc and the adjacent region. See Fig. 10 (1). A shadow coloring of a branch point diagram is the assignment of a single quandle element, say $`b`$, to the $`2`$-dimensional faces of the diagram. There are three $`3`$-dimensional regions in a neighborhood of the branch point. One such region is assigned the color $`a`$; the other two are assigned the color $`ab`$ or $`(ab)b`$, and the assignment is determined by the condition that normal vectors point into regions that receive quandle products. See Fig. 10 (2) and (3). A shadow coloring of a hem $`1`$-crossing diagram is an assignment of a single quandle element to the region and the disk. See Fig. 10 (4). A shadow coloring of a hem $`2`$-crossing diagram is an assignment of an element $`a`$ on the disk and the region away from which the normal of the disk dividing the hem, the dividing disk receives $`b`$, and then, the region and the disk into which the normal points receive the element $`ab`$. See Fig. 10 (5). Let $`n=1,2`$. A shadow coloring of an abstract $`n`$-knot diagram is an assignment of quandle elements to the $`n`$ and $`(n+1)`$-dimensional regions of the diagram such that the restriction to each endpoint, branch point, hem or $`k`$-crossing diagram is a shadow coloring. A shadow coloring is, as before, a map $`𝒞:X`$ where $``$ is the set of arcs ($`n=1`$) or regions ($`n=2`$) and complementary $`(n+1)`$-dimensional regions. ###### 4.4 Example. Figure 11 indicates a shadow quandle coloring by $`R_3`$ of an abstract $`0`$-knot diagram. The small flag-like arrows at the 1-crossing points indicate the orientation of the vertices. When this arrow coincides with the tangent direction of the circle (as in the figure), the crossing is positive; otherwise the crossing is negative. ###### 4.5 Example. Figure 12 indicates a shadow quandle coloring by $`R_3`$ of the abstract $`1`$-knot diagram that was given in Fig. 5. The boundary of the surface is abbreviated and is not drawn in Fig. 12 for simplicity. In the figure, colors on regions are indicated by letters in squares. ###### 4.6 Remark. Shadow colorings were defined in and used in to prove that the right and left trefoil knots are distinct. The notions of coloring and shadow coloring of $`k`$-crossing $`n`$-diagrams extend for all $`n=1,2,\mathrm{},`$ and for all $`k=0,\mathrm{},n+1`$ (see ). To see that shadow colorings exist, we prove the following lemmas. ###### 4.7 Lemma. Let $`𝒞`$ be a coloring by a quandle $`X`$ of an oriented $`k`$-crossing $`n`$-diagram, $`(kn)`$, which is obtained from an oriented $`(k+1)`$-crossing $`n`$-diagram by ingoring the level $`k+1`$ sheet, $`E_{k+1}`$. Let $`D`$ be an $`n`$-region on $`E_{k+1}`$, and let $`xX`$. Then there is a unique coloring $`𝒞^{}`$ by $`X`$ of the $`(k+1)`$-crossing diagram which restricts to $`𝒞`$ and such that $`𝒞^{}(D)=x`$. Proof. Let $`D_j`$, $`j=1,\mathrm{},2^{k+1}`$, be the $`n`$-regions on $`E_{k+1}`$. Let $`D=D_1`$, and pick points $`p_jD_j`$. For any fixed $`j`$, let $`\gamma `$ be a path in the level $`k+1`$ sheet (the image of a continuous map from $`[0,1]`$ to the level $`k+1`$ sheet) from $`p_1`$ to $`p_j`$. Assume without loss of generality that the path is in general position with the level $`g`$ sheets for all $`g`$, so that it meets the sheets in finitely many points. Let $`s_1,\mathrm{},s_r`$ be the intersection points, and let $`ϵ_i`$, $`i=1,\mathrm{},r`$ be $`+1`$ (resp. $`1`$) if the path goes in the same (resp. the opposite) direction as the normal of the sheet at $`p_i`$. Let $`𝒞(p_i)=c_i`$ (precisely speaking, the color of the disk in which the point $`p_i`$ lies). Then define the color $`𝒞^{}(D_j)=xw`$ where $`w`$ is the word $`c_1^{ϵ_1}c_2^{ϵ_2}\mathrm{}c_n^{ϵ_r}.`$ This definition is made in such a way that the condition of coloring is satisfied at each intersection point along the path $`\gamma `$, and in remains to be proved that the color thus defined does not depend on the choice of the path $`\gamma `$. Let $`\gamma _i`$, $`i=0,1`$, be such two paths. Since the level $`k+1`$ sheet is simply connected, there is a homotopy between them. Such a homotopy is a map from a $`2`$-disk to the level $`k+1`$ sheet, whose image is denoted by $`U`$. Assume without loss of generality that $`U`$ is in general position with the level $`g`$ sheets for all $`g(<k)`$. Then the intersection between $`U`$ and the level $`g`$ sheets for all $`g(<k)`$ is generically immersed $`1`$-manifold with boundary, i.e., arcs with transverse double points. When the $`\epsilon `$-neighborhood in the definition of the crossing diagram is removed from these arcs, then we obtain a classical knot diagram on $`U`$ with boundary of arcs liying on $`U`$. Give the color by $`X`$ on the regions and arcs in $`U`$ by using the rule of the shadow color depicted in Fig. 10. Then the color is well-defined on the diagram on $`U`$, proving that the color at the end point of $`\gamma _i`$, $`i=0,1`$, coincide. $`\mathrm{}`$ ###### 4.8 Lemma. Let $`K=[f:M^nN^{n+1}]`$ be an abstract $`n`$-knot diagram, and let $`𝒞`$ be a coloring by a quandle $`X`$. Suppose $`N`$ is simply connected. Then for any $`(n+1)`$-region $`R`$ and $`xX`$, there is a shadow color $`𝒞^{}`$ for $`K`$ such that $`𝒞^{}`$ restricts $`𝒞`$ and $`𝒞^{}(R)=x`$. Proof. This is proved by defining colors using paths as in Lemma 4.7, and the proof is similar. The condition $`\pi _1(N)=\{1\}`$ guarantees the existence of a homotopy between paths. In the case $`n=1`$ and $`2`$, we included end/branch points and hem $`1`$\- and $`2`$-crossings, in addition to $`k`$-crossing $`n`$-diagrams. For these diagrams, a similar argument applies, using the definition of shadow colors for these diagrams. $`\mathrm{}`$ ###### 4.9 Extended Remark. Let $`\text{SC}(1)`$ denote the collection of shadow colored abstract $`1`$-knot diagrams, and let $`C(2)`$ denote the collection of colored abstract closed $`2`$-knot diagrams. We have maps $`𝒟:C(2)\text{SC}(1)`$ and $`:\text{SC}(1)C(2)`$ that are defined as follows (see also Fig. 13 for a local description). Consider a colored closed abstract $`2`$-knot diagram, $`K^2=[f:M^2N^3]`$. The map $`𝒟`$ assigns to $`K^2`$ the following shadow colored $`1`$-knot diagram. The abstract $`1`$-knot diagram is the lower decker set in a regular neighborhood of the lower decker set in the surface $`M`$ — the lower decker points form the $`1`$-crossings and $`2`$-crossings. Here, lower decker set (see ) means the preimage of the double point set of the surface $`M`$ that are in the under sheet and represented as the broken sheet in the abstract $`2`$-knot diagram. The $`2`$-crossing points of $`𝒟(K^2)`$ correspond to the triple points of the abstract $`2`$-knot diagram, and the endpoint diagrams correspond to branch points (see for details). We color the arcs of $`𝒟(K^2)`$ with the quandle elements that appear on the regions that contain the corresponding upper sheets. At a 2-crossing of $`𝒟(K^2)`$, the over-arc is colored by the color on the upper sheet, the under arcs are colored by the colors on the two portions of the corresponding middle sheet. The 2-dimensional regions $`𝒟(K^2)`$ are colored by the quandle elements that are on the pieces of the surface that contain the arc diagram. On the other other hand, the map $``$ assigns to an abstract shadow colored 1-knot diagram, $`K^1`$, a colored abstract $`2`$-knot diagram as follows. Decompose the abstract $`1`$-knot diagram into crossing and endpoint diagrams. Construct an abstract $`2`$-knot diagram using the local correspondence between $`k`$-crossing $`1`$-diagrams and $`(k+1)`$-crossing $`2`$-diagrams, and between endpoints and branch points that is depicted in Fig. 13. Colors on the surface are determined as in the figure. In this way, a surface with double points on the boundary is constructed. The colors on the double points on the boundary just above and below the surface of the arc diagram agree. Thus these double points can be joined together as they were in Fig.7. Given $`K^1`$, the $`1`$-knot diagram $`𝒟((K^1))`$ differs from $`K^1`$ in that in addition to $`K^1`$, it contains unknotted unlinked colored components that may not be found in $`K^1`$. In particular, given a shadow colored classical knot diagram, we obtain a colored abstract $`2`$-knot diagram. If $`K`$ is a knot, the $`2`$-knot diagram consists of a sphere (in the plane of the knot diagram) and a torus of the form $`K\times S^1`$. The $`2`$-knot diagram is in $`S^2\times S^1`$. We call this the suspension of a classical knot. See Fig. 14. A similar construction applies to give maps $`𝒟:C(1)\text{SC}(0)`$ and $`:\text{SC}(0)C(1)`$ between the sets, $`C(1)`$, of colored abstract closed $`1`$-knot diagrams and $`\text{SC}(0)`$ of shadow colored $`0`$-knot diagrams. The essense of this construction is found in (see also Fig. 15). ## 5 Representing Homology Classes In Greene (see also ), rack cycles are represented by colored knot diagrams. We generalize this method to quandle cycles, using end/branch point and hem diagrams. Specifically, for a fixed quandle $`2`$\- or $`3`$-cycle we construct a colored or shadow colored diagram that represents the cycle. Moreover, if a cycle is a boundary, we construct an abstract $`2`$-knot diagram in a 3-manifold, the boundary of which is the given abstract $`1`$-knot diagram that represents the cycle. Care will be taken at endpoint diagrams to make this precise. Finally, we represent $`4`$-cycles by shadow colored abstract $`2`$-knot diagrams with hems. We do not pursue boundaries in this case. ###### 5.1 Lemma. (1) There is a one-to-one correspondence between the set of $`(n+1)`$-tuples of quandle elements and the set of quandle colored positive $`(n+1)`$-crossing $`n`$-diagrams. (2) There is a one-to-one correspondence between the set of $`(n+1)`$-tuples of quandle elements and shadow colored positive $`n`$-crossing $`n`$-diagrams. Proof. The correspondence (1) for $`n=1`$ is illustrated in Fig. 16 (A). The correspondence (2) for $`n=1`$ is illustrated in Fig. 16 (B), where the correspondence between the shadow colors and the lower decker set is illustrated. The correspondence (1) for $`n=2`$ is illustrated in the center of Fig. 17. The correspondence (2) for $`n=2`$ is illustrated in Fig 18. In general, let $`X`$ be a quandle and consider an $`(n+1)`$-tuple $`(x_1,\mathrm{},x_{n+1})X^{n+1}`$. Recall that $`E_j`$ is the coordinate $`n`$-disk in $`[1,1]^{n+1}`$ whose $`j`$th coordinate is $`0`$. Color the region of $`E_j`$ at which the $`\mathrm{}`$th coordinates are less than $`\epsilon `$ for $`\mathrm{}<j`$ with $`x_{n+2j}.`$ Color the remaining regions of $`E_j`$ in such a way that the quandle condition holds at each double point of the diagram. Such a coloring can be uniquely extended to a shadow coloring by coloring the $`(n+1)`$-region of $`E`$ at which all the coordinates are less than $`\epsilon `$ by $`x_0`$. Thus associated to an $`(n+1)`$-tuple there is a colored $`(n+1)`$-crossing diagram or a shadow colored $`n`$-diagram. By choosing the color $`x_{n+2j}`$ from the region of $`E_j`$ at which the $`\mathrm{}`$th coordinates are less than $`\epsilon `$ for $`\mathrm{}<j`$, we construct an $`(n+1)`$-tuple. A shadow colored diagram works similarly. $`\mathrm{}`$ ###### 5.2 Scholium. The boundary $`(x_1,\mathrm{},x_{n+1})`$ of a generating rack chain $`(x_1,\mathrm{},x_{n+1})`$ corresponds to the boundary of the quandle colored $`(n+1)`$-crossing diagram. Proof. The correspondence is illustrated in the Figs. 16 (A) and 17. We leave the rest of the details to the reader. $`\mathrm{}`$ ###### 5.3 Remark. Colored and shadow colored abstract knot diagrams represent chains — formal-sums of signed $`(n+1)`$-tuples (or $`(n+2)`$-tuples) of quandle elements corresponding to the $`0`$-dimensional crossings as in Lemma 5.1. Recall that the $`0`$-dimensional crossings of an abstract $`n`$-knot diagram are the $`(n+1)`$-crossing points. Each such (shadow) colored crossing point in the diagram has an associated sign. The chain determined by the diagram is the sum of these signed $`(n+1)`$-tuples (or $`(n+2)`$-tuples) taken over all the $`0`$-dimensional crossings. This is the chain represented by a (shadow) colored diagram. An abstract diagram (of a closed manifold $`M`$) that has no exceptional points has its $`1`$-dimensional crossing set closed. Therefore the sum of the representative chain is a rack cycle in this case . In low dimensions, shadow colored exceptional points or colored branch points represent degenerate chains. Specifically, a shadow colored endpoint or a colored branch point represents a chain of the form $`\pm (a,a)`$, a shadow colored branch point represents a chain of the form $`\pm (a,b,b)`$ and a shadow colored hem $`2`$-crossing point represents a chain $`\pm (a,a,b)`$. The signs are determined by the direction of the arrow along the arc that originates or terminates at the exceptional point. Thus a shadow colored abstract $`0`$-knot diagram represents a quandle $`2`$-cycle (for example, Fig. 11), as does a colored closed abstract $`1`$-knot diagram (for example, Fig. 4). A shadow colored abstract $`1`$-knot diagram represents a quandle $`3`$-cycle as does a colored closed abstract $`2`$-knot diagram (for example, Fig. 7). And a shadow colored abstract $`2`$-knot diagram represents a quandle $`4`$-cycle. For example the 2-twist-spun trefoil may be shadow colored by $`R_3`$ (see , and apply Lemma 12). This is the context of Theorem 5.5 (1a) and (1b). ###### 5.4 Examples. The $`R_3`$ $`2`$-chain $`(\alpha ,\beta )+(\gamma ,\alpha )+(\beta ,\gamma ),`$ represented by the diagram in Fig. 11 is a cycle in $`Z_2^\mathrm{Q}(R_3)`$. The $`R_3`$ $`2`$-chain represented by the colored abstract $`1`$-knot diagram in Fig. 4 is the cycle $`(\alpha ,\beta )+(\beta ,\gamma )(\beta ,\alpha )Z_2^Q(R_3;𝐙)`$. It is known that $`H_2^Q(R_3;𝐙)=0`$, so that the above examples are in fact boundaries. The $`3`$-chain represented by the colored abstract $`2`$-knot diagram in Fig. 7 is the cycle $`(\alpha ,\beta ,\gamma )+(\alpha ,\gamma ,\alpha )Z_3^Q(R_3;𝐙)`$. The corresponding (in the sense of 4.9) shadow colored abstract $`1`$-knot diagram is depicted at the right bottom of the figure, which is the same as Fig. 12. It is known that $`H_3^Q(R_3;𝐙)𝐙_3`$, and the above $`3`$-cycle is a generator. ###### 5.5 Theorem. Let $`X`$ denote a quandle. (1a) Let $`n=1,2`$. Any colored closed abstract $`n`$-knot diagram represents a quandle $`(n+1)`$-cycle in $`Z_{n+1}^\mathrm{Q}(X;𝐙)`$. (1b) Let $`n=1,2,3.`$ Any shadow colored abstract $`(n1)`$-knot diagram (possibly with exceptional points if $`n=2,3`$) represents a quandle $`(n+1)`$-cycle in $`Z_{n+2}^\mathrm{Q}(X;𝐙)`$. (2a) Let $`n=1,2`$. Let $`\eta Z_{n+1}^\mathrm{Q}(X;𝐙)`$. Then there is a colored $`n`$-knot diagram, $`𝒟_\eta ^n`$, that represents $`\eta `$. (2b) Let $`n=1,2,3.`$ Let $`\eta Z_{n+1}^\mathrm{Q}(X;𝐙)`$. Then there is a shadow colored $`(n1)`$-knot diagram (possibly with exceptional points for $`n=2,3`$), $`𝒮D_\eta ^{n1}`$, that represents $`\eta `$. Proof. Statements (1a) and (1b) (for $`n=2,3`$), follow from Remark 5.3. For $`n=1`$ statement (1b) is also easy: A shadow colored $`0`$-knot diagram represents a rack $`2`$-cycle which is also a quandle $`2`$-cycle. Consider statement (1b) for $`n=2`$, the endpoints of shadow arc diagrams represent chains of the form $`(a,a)`$ which are trivial in quandle homology. Thus shadow colored arc diagrams represent quandle $`2`$-cycles. Branch points of shadow colored abstract $`2`$-knot diagrams represent chains of the form $`(a,b,b)`$ while hem $`2`$-crossing points represent chains of the form $`(a,a,b)`$. Both such chains are trivial in quandle homology. This proves (1b) for $`n=3`$. Consider a quandle $`2`$-cycle $`\eta =_jϵ_j\stackrel{}{x}_jZ_2^Q(X;𝐙)`$ where $`\stackrel{}{x}_j=(x_1^j,x_2^j)`$ and $`ϵ_j=\pm 1`$. For each $`(x_1^j,x_2^j)`$, construct a colored $`2`$-crossing $`1`$-diagram as in Lemma 5.1 (1); the sign of the crossing is $`ϵ_j`$. The boundary of such a chain is $`(x_1^j,x_2^j)=x_1^jx_1^jx_2^j`$. Since $`\eta `$ is a cycle the sum of these boundary terms adds to $`0`$. Thus we can interconnect the under-arcs to form a collection of simple closed curves of under-arcs. The over-arcs each can be joined end-to-end locally to form the abstract diagram as depicted in Fig. 15 right top. For example, Fig. 4 illustrates an abstract diagram that represents the cycle $`(\alpha ,\beta )+(\beta ,\gamma )(\beta ,\alpha )Z^\mathrm{Q}(R_3)`$. By Remark 4.9, a desired shadow colored $`0`$-knot diagram for (2b) is obtained from the above constructed representative $`1`$-knot diagram. Thus statements (2a) and (2b) hold for $`n=1`$. Note that the way canceling pairs are joind together is not unique. Now suppose that $`\eta =_jϵ_j\stackrel{}{x}_jZ_3^Q(X;𝐙)`$ (where $`\stackrel{}{x}_j=(x_0^j,x_1^j,x_2^j)`$ and $`ϵ_j=\pm 1`$) is a $`3`$-cycle. We represent each $`ϵ_j\stackrel{}{x}_j`$ by a shadow colored $`2`$-crossing $`1`$-diagram as in Lemma 5.1 (2); the sign of the crossing is $`ϵ_j`$. The boundary of such a chain is the sum of four $`2`$-chains. Each $`2`$-chain is represented as a shadow colored $`0`$-knot diagram. Take the cartesean product of such a diagram with the unit interval to form a shadow colored $`1`$-crossing $`1`$-knot diagram. Such a diagram is a $`1`$-handle the attaching region of which is identified with the segments on the boundary of the squares that represent the chains $`\stackrel{}{x}_j`$. Thus we can attach $`1`$-handles to these squares. Since $`\eta `$ is a quandle cycle, the boundary of squares that are not attached to these $`1`$-handles have the same color assigned to the arc and the region. Cap each of these boundary segments by a colored endpoint diagram. This constructs a colored abstract $`1`$-knot diagram with endpoints representing $`\eta `$. To construct an abstract $`2`$-knot diagram that represents $`\eta `$, we follow the procedure outlined in Remark 4.9. This gives statement (2a) and (2b) for $`n=2`$. If $`\eta =_jϵ_j\stackrel{}{x}_j`$ (where $`\stackrel{}{x}_j=(x_0^j,x_1^j,x_2^j,x_3^j)`$ and $`ϵ_j=\pm 1`$) is a $`4`$-cycle, then we construct a shadow colored abstract $`2`$-knot diagram that may have hems following a procedure analogous to the ones above. Specifically, for each chain $`\stackrel{}{x}_j`$ we construct a shadow colored triple point diagram where the colors on the regions away from which normals point is $`x_i^j`$. The boundary of the chain $`\stackrel{}{x}_j`$ corresponds to the shadow colored $`2`$-crossing diagrams on the boundary of the colored cubes. We attach shadow colored $`1`$-handles between pairs of square faces to cancel like terms. Since $`\eta `$ is a quandle cycle, its boundary consists of terms of the form $`(a,a,b)`$ and $`(a,b,b)`$. There are faces that remain representing such terms. In case the term is of the form $`(a,b,b)`$, we attach a branch point as in Fig. 19. In case the term is the form $`(a,a,b)`$, then we extend the over-sheet to create a hem diagram. The double point of the hem diagram represents this chain. This completes the proof of (2b) for $`n=3`$. $`\mathrm{}`$ ###### 5.6 Remark. We included condition (1) of Definition 3.8 for expository convenience. Otherwise boundary points might be confused with hems or endpoints. In the following theorem, we use abstract diagrams with actual boundaries — diagrams for which condition (2) is satisfied, but not necessarily condition (1). The notions of colored and shadow colored diagrams extend to diagrams with actual boundary in a straight-forward way. In ordinary homology theory, boundaries of complexes correspond to boundary terms of chains. This is our motivation of extending the abstract diagrams to those with boundaries, as will be seen in the following theorem. The exceptional points that are defined above are not regarded as boundaries in quandle homology, as they were introduced to represent degenerate chains, instead of boundary terms. For example, an abstract $`1`$-knot diagram with actual boundary $`[f:MN]`$ has two types of endpoints. One is the endpoint diagram defined already. The actual boundary end point lies on the boundary $`N`$. When a diagram is shadow colored by a quandle, the former (the endpoint diagram) corresponds to a degenerate chain, and the latter corresponds to the sum of the boundary terms in $`(x_0,x_1,x_2)`$. See the right-side of Fig. 20, for example. In case a shadow colored $`2`$-knot diagram with actual boundary has a hem, the hem may have a boundary point in the form of an endpoint diagram, but the hem itself is not regarded as the actual boundary. ###### 5.7 Theorem. Let $`n=1,2`$. If the given cycle, $`\eta Z_{n+1}^\mathrm{Q}(X;𝐙)`$, is a boundary, so $`\eta =\nu `$, for $`\nu C_{n+2}^\mathrm{Q}(X;𝐙)`$, then $`𝒟_\eta ^n`$ (or $`𝒮D_\eta ^{n1}`$) is the boundary of a colored $`(n+1)`$-knot diagram, $`𝒟_\nu ^{n+1}`$, (or shadow colored $`n`$-knot diagram, $`𝒮D_\nu ^n`$,) with actual boundary (and with hems when appropriate) that represents $`\nu `$. Proof. Suppose that the quandle $`(n+1)`$-cycle $`\eta =\nu `$. In general, there is a colored abstract $`(n+1)`$-knot diagram (or shadow colored abstract $`n`$-knot diagram) with actual boundary that represents $`\nu `$. For example, take the disjoint union of colored $`(n+2)`$-crossing $`(n+1)`$-diagrams (or shadow colored $`(n+1)`$-crossing $`n`$-diagrams) that represent the chains which constitute $`\nu `$. Let $`(𝒮)𝒟_\nu `$ denote such a diagram, where the parenthesis represents that it is either a colored diagram (without $`𝒮`$) or shadow colored diagram (with $`𝒮`$). We also use the notation $`(𝒮)𝒟_\nu =[g:M_\nu N_\nu ]_{𝒞_\nu }`$ to specify the map $`g`$. Similarly, let $`(𝒮)𝒟_\eta =[f:M_\eta N_\eta ]_{𝒞_\eta }`$ denote a given colored diagram that represents $`\eta `$. The sum of the (shadow) colored crossings of $`(𝒮)𝒟_\nu `$ may differ from the sum of the (shadow) colored crossings on $`(𝒮)𝒟_\eta `$ in that either may include degenerate chains or canceling terms. That is, there can be crossings that represent chains of the form $`(a,a)`$ in case $`n=1`$, or chains of one of the forms $`(a,a,b)`$ or $`(a,b,b)`$ if $`n=2`$. Or there may be a canceling pair of (shadow) colored crossings. Take the product $`(𝒮)𝒟_\eta \times [0,1]`$ which is $`g(M_\eta )\times [0,1]N_\eta \times [0,1]`$ with the given coloring ($`\times [0,1]`$), where the original $`(𝒮)𝒟_\eta `$ is regarded as embedded in $`N_\eta \times \{0\}`$. For each pair of canceling (shadow) colored crossings, join them by a (shadow) colored $`1`$-handle (with coloring induced from the crossings), and cancel them. Such $`1`$-handles are attached on $`N_\eta \times \{1\}`$, and give a new diagram $`[g^{}:M_\eta ^{}N_\eta ^{}]`$ with a (shadow) color such that $`N_\eta ^{}`$ consists two pieces, one of which $`_1N_\eta ^{}`$ contains the original $`g(M_\eta )`$, and the other $`_2N_\eta ^{}`$ contains $`g(M_\eta )`$ with all canceling pairs eliminated. Perform the same process on $`N_\nu `$ (we continue to use the same notation $`N_\nu `$ after this process), so that we now assume that $`_2N_\eta ^{}`$ and $`N_\nu `$ do not contain canceling pairs. Next we cancel degenerate chains, case by case. Case 1. Suppose that $`\eta `$ is a $`2`$-cycle represented by a shadow colored $`0`$-knot diagram, $`𝒮D_\eta =[f:M_\eta N_\eta ]_{𝒞_\eta }`$. Then $`M_\eta `$ is a $`0`$-dimensional manifold, $`N_\eta `$ is a $`1`$-dimensional manifold. We call the components of $`M_\eta `$ oriented vertices. A $`2`$-chain $`(a,b)`$ is represented by a shadow colored $`1`$-crossing $`0`$-diagram where $`b`$ is the color of the vertex and $`a`$ is the color on the arc away from which the normal arrow to the vertex points. The description of $`𝒮D_\nu `$ is the same. For each degenerate vertex (shadow colored as $`(a,a)`$), we construct a shadow colored endpoint diagram with color $`a`$ on the $`2`$-dimensional region and color $`a`$ on the edge. The orientation of the arc is determined by the sign of the $`0`$-crossing. Each such endpoint diagram is a $`0`$-handle that is disjoint from $`N_\eta \times [0,1]`$ and $`N_\nu `$. Now attach a $`1`$-handle (in the guise of a shadow colored $`1`$-crossing $`1`$-diagram) between the degenerate crossing and the endpoint diagram. In this way, we may assume that there are no degenerate $`2`$-chains among the $`1`$-crossings representing the summands of $`\eta `$ and $`\nu `$. See Fig. 20 for an example. Case 2. Suppose that $`\eta `$ is a $`2`$-cycle that is represented by a colored $`1`$-diagram. Each degenerate colored $`2`$-crossing in $`𝒟_\eta `$ or $`𝒟_\nu `$ is a crossing at which all the colors are the same; thus a degenerate crossing represents a chain for the form $`(a,a)`$. For each such crossing we construct a colored branch point diagram that functions as a $`0`$-handle. The color on the surface is $`a`$. This handle is attached to the diagram via a $`1`$-handle in the guise of a colored $`2`$-crossing $`2`$-diagram. Therefore, as above we can assume that there are no degenerate $`2`$-chains among the crossings representing the summands $`\eta `$ and $`\nu `$. See Fig. 19. Case 3. Suppose $`\eta `$ is a $`3`$-cycle that is represented by a shadow colored $`1`$-knot diagram, $`𝒮D_\eta `$. The degenerate chains are of the form $`(a,a,b)`$ and $`(a,b,b)`$. The former are represented by colored crossings with color $`a`$ on one face, color $`a`$ on one of the lower arcs, and color $`b`$ on the upper arc. The latter are represented by shadow colored crossings where $`a`$ is the color on one face, and $`b`$ is the color on all edges. For each degenerate chain of the form $`(a,a,b)`$, we insert a shadow colored hem $`2`$-crossing diagram as a $`0`$-handle disjoint from $`N_\eta ^{}`$ and $`N_\nu `$. For each degenerate chain of the form $`(a,b,b)`$ we insert a shadow colored branch point diagram as a $`0`$-handle. Then the $`0`$-handles are attached to either $`_1N_\eta ^{}`$ or $`N_\nu `$ via $`1`$-handles that are $`2`$-crossing $`2`$-diagrams. The $`1`$-handles are attached to the $`0`$-handles at the crossing diagram on the boundary. Thus in this case, we may assume that there are no degenerate $`3`$-chains among the crossings representing the summands of $`\eta `$ and $`\nu `$. Next, we consider the cases when the colored crossings on the resulting manifolds (that we again denote by $`_1N_\eta ^{}`$ and $`N_\nu `$) are in one-to-one correspondence. In each case, we can attach $`1`$-handles between pairs of similarly (shadow) colored crossings. The $`1`$-handles are of the form: (a) shadow colored $`n`$-crossing $`n`$-diagrams in case 1 and 3; (b) colored $`(n+1)`$-crossing $`(n+1)`$-diagrams in case 2. After all of the crossings that represent summands of $`\eta `$ have been glued to the chains representing $`\nu `$, we attach (shadow) colored $`2`$-handles in the guise of (1) shadow colored $`0`$-crossing $`0`$-diagrams in case 1. (2) $`0`$-crossing $`1`$-diagrams in case 2. (3) shadow colored $`1`$-crossing $`2`$-diagrams in case 3. This completes the proof. $`\mathrm{}`$ ## 6 Equivalence of Colored and Shadow Colored Diagrams So far we have introduced abstract diagrams, colorings, and shadow colorings thereof. Here we discuss moves to such diagrams. ###### 6.1 Theorem. For $`i=0,1`$, let $`𝒮D_i`$ denote shadow colored abstract $`0`$-knot diagrams where the color set is the quandle $`X`$. Suppose that $`[𝒮D_0]=[𝒮D_1]H_2^\mathrm{Q}(X;𝐙)`$. Then $`𝒮D_0`$ can be obtained from $`𝒮D_1`$ by a finite sequence of moves taken from those depicted in Fig. 21. Proof. Let $`𝒮D_i`$ be represented by maps with colorings $`[f_i:M_iN_i]_{𝒞_i}`$ as in the proof of the preceding theorem. By Theorem 5.5, $`[f_0:M_0N_0]_{𝒞_0}`$ and $`[f_1:M_1N_1]_{𝒞_1}`$ cobound a shadow colored abstract $`1`$-knot diagram $`M`$ in a $`2`$-manifold $`N`$ with actual boundary. Let $`F:N[0,1]`$ be a smooth function such that $`F^1(i)=N_i,`$ for $`i=0,1.`$ We may assume (after a small perturbation if necessary) that $`F`$ satisfies the following conditions: (1) $`F`$ is transverse at $`0`$ and $`1`$. (2) $`F`$ is generic on $`M`$, $`N`$, $`M`$, $`N`$, and on all the self intersections and singularities of $`M`$. Thus $`F`$ has isolated Morse critical points on all the sets listed in (2), at distinct critical values. By taking the inverse images $`F^1(hϵ)`$ and $`F^1(h+ϵ)`$ at every critical value $`h`$, we obtain a sequence of moves. Hence we obtain the result by classifying these Morse critical points as follows. They are classified into two categories: Reidemeister moves for $`0`$-dimensional knot diagrams as singular sets and critical points of $`M`$ (Fig. 21 left $`3`$ figures), and Morse critical points of $`N`$ (Fig. 21 right $`2`$ figures). The left $`3`$ figures correspond to the endpoints of $`M`$, maxima/minima of $`M`$, and transverse double points of $`M`$, respectively. The right $`2`$ figures correspond to the maxima/minima of $`N`$, and saddle points of $`N`$, respectively. These exhaust generic critical points of $`f`$, and the theorem follows. $`\mathrm{}`$ ###### 6.2 Theorem. Let $`X`$ be a quandle. For $`i=0,1`$, let $`𝒮D_i=[f_i:M_iN_i]_{𝒞_i}`$ denote shadow colored $`1`$-knot diagrams that represent the same homology class in $`H_3^\mathrm{Q}(X;𝐙)`$. Then $`[f_1:M_1N_1]_{𝒞_1}`$ can be obtained from $`[f_0:M_0N_0]_{𝒞_0}`$ by a finite sequence of moves taken from those depicted in Figs. 22 and 23. Proof. By Theorem 5.5, $`[f_0:M_0N_0]_{𝒞_0}`$ and $`[f:M_1N_1]_{𝒞_1}`$ cobound a shadow colored abstract $`2`$-knot diagram $`M`$ in a $`3`$-manifold $`N`$ with actual boundary. Let $`F:N[0,1]`$ be a smooth function such that $`F^1(i)=N_i`$, for $`i=0,1`$. We may assume (after a small perturbation if necessary) that $`F`$ satisfies the following conditions. (1) $`F`$ is transverse at $`0`$ and $`1`$. (2) $`F`$ is generic on $`M`$, $`N`$, $`M`$, $`N`$, and on all the self intersections and singularities of $`M`$. Thus $`F`$ has Morse isolated critical points on all the sets listed in (2), at distinct critical values. The proof proceeds as follows in a similar way as in the preceding theorem. The singularities and critical points of $`M`$ are listed in Fig. 22 and are Reidemeister moves. From top to bottom they represent the intersection of the boundary points and interior of $`M`$, branch points, maxima/minima of double curves of $`M`$, and triple points, respectively. The Morse critical points as handle moves are listed in Fig. 23. From top to bottom, they are critical points of $`N`$, Int$`N`$, $`M`$, and Int$`M`$. The critical points of $`N`$ are maxima/minima (the left entry) or saddle points (the right). From the point of view of the boundary $`1`$-manifold, they correspond to handles of indices $`0/2`$ and $`1`$, respectively. The critical points of Int$`N`$ are similar, and depicted in the second row left and right. The critical points of $`M`$ are the creation or deletion of a pair of points. There are two types, left and right, of these $`0/1`$-handles, in relation to interior points. The bottom entry in the figure illustrates the critical points of the interior of $`M`$. The Theorem follows as these exhaust generic singularities and critical points. $`\mathrm{}`$ ###### 6.3 Scholium. Colored abstract closed 1-knot diagrams represent the same $`2`$-dimensional quandle homology class if and only if one is obtained from the other by a finite number of moves that are taken from those depicted in Figs. 22 and 23 where: (1) those moves that involve endpoints are excluded; and (2) colors on the 2-dimensional regions are ignored. Proof. Imitate the proof of Theorem 6.2 replacing shadow colored with colored throughout. Since $`M_i`$ are closed, we may assume that the cobounded surface $`M`$ has no hems. Thus the moves involving endpoints are excluded. $`\mathrm{}`$ ###### 6.4 Remark. For higher dimensions, the authors expect similar theorems. The classifications of moves for boundaries, however, become subtle. For one dimensional higher case, for example, the moves include Roseman moves (see for example ) with quandle colors, Morse critical points of the ambient space $`N`$ and $`N`$, and singularities and critical points involving hems will provide the set of moves. The moves for hems are depicted in Fig. 27. ###### 6.5 Discussion. Abstract $`1`$-knot diagrams that have no endpoints correspond to virtual knot diagrams (see ). In , a set of moves to abstract knot diagrams was given, and it was shown that up to these moves the set of abstract diagrams is equivalent to virtual knots up to virtual Reidemeister moves. Also, the correspondence was made by thickening the virtual knots and making them abstract knots. However, the moves to abstract knots are not to be confused with virtual Reidemeister moves. Virtual Reidemeister moves for thickened virtual knots are not equivalence moves for abstract knots thus obtained. For example, type II moves to thickened virtual knots do not necessarily preserve shadow coloring. We defined moves to colored and shadow colored diagrams to avoid diagrams that are type II equivalent but that have different quandles. Therefore, quandle homology provides an invariant of colored diagrams under (shadow) colored Reidemeister moves listed above. ## 7 The Quandle of an Abstract Knot and Examples The (fundamental) quandle of an abstract $`n`$-knot diagram is defined in . It is generated by the $`n`$-regions of the diagram; the relations in the quandle can be read from the $`2`$-crossings. It generalizes the knot quandle of classical knots (). See for Wirtinger presentations of knot quandles defined from knot diagrams, which are similar to Wirtinger presentations of knot groups. In this case, arcs of knot diagrams represent generators, and crossings give relations of Wirtinger form. Let $`Q(K)`$ represent such a quandle for a knot diagram $`K.`$ ###### 7.1 Example. If $`K_0`$ is the diagram illustrated in Fig. 4, then its quandle has presentation $$Q(K_0)=<x,y,z:xy=z,yz=x,yx=z>.$$ The quandle $`<x,y,z:xy=z,yz=x,yx=z>`$ is isomorphic to the 3-element dihedral quandle $`R_3`$. The abstract knot $`K_0`$ represents the $`2`$-cycle $`(x,y)+(y,z)(y,x)Z_2^\mathrm{Q}(Q(K_0))`$. ###### 7.2 Proposition. Every abstract closed $`1`$-knot diagram $`K`$ represents a cycle $`[K]Z_2^\mathrm{Q}(Q(K))`$. The abstract knot $`K_0`$ is a boundary since the dihedral quandle, $`R_3`$, has trivial $`2`$-dimensional homology. ###### 7.3 Question. Under what circumstances is the cycle $`[K]`$ a boundary? ###### 7.4 Discussion and Example. On the other hand, we may define a shadow quandle for an abstract diagram. We can show that the $`3`$-dimensional homology class that is represented by a shadow coloring of Fig. 4 is non-trivial as follows. The 5-element quandle $`QS(5)`$ (consisting of non-identity permutations in the symmetric group on 3 letters) colors the diagram non-trivially. Transpositions go on the arcs and $`3`$-cycles go on the $`2`$-dimensional regions. By a Mathematica calculation, we have determined that the $`3`$-cycle represented thereby is not a boundary. ###### 7.5 Example. The example depicted in Fig. 24 represents a $`3`$-cycle over the quandle $`QS(6)`$. The homomorphism $`p:QS(6)R_3`$ defined in 2.2 induces a map on homology. The corresponding $`R_3`$ shadow colored diagram represents a trivial $`3`$-cycle. We can show, by mean of a Mathematica calculation, that the $`QS(6)`$ colored cycle is non-trivial. The same program show that the colored diagram on the left of Fig. 25 represents a generator of $`H_3^\mathrm{Q}(QS(6);𝐙)=𝐙_{24}`$. The induced map $`f_{}`$ is illustrated to be surjective. ###### 7.6 Remark. Let $`K`$ be a knot diagram on a compact oriented surface $`F`$. Then the fundamental shadow quandle $`SQ(K)`$ is defined as follows. The generators correspond to over-arcs and connected components of $`F\text{universe of}K`$. The relations are defined for each crossing as ordinary fundamental quandles, and at each arc dividing regions. Specifically, if $`a`$ and $`b`$ are generators corresponding to adjacent regions such that the normal points from the region colored $`a`$ to that colored $`b`$, and if the arc dividing these regions is colored by $`c`$, then we have the relation $`b=ac`$. This defines a presentation of a quandle, which is called the fundamental shadow quandle of $`K`$. Two diagrams on $`F`$ that differ by Reidemeister moves on $`F`$ have isomorphic fundamental shadow quandles. The shadow colors are regarded as quandle homomorphisms from the fundamental shadow quandle to a quandle $`X`$. ## 8 Boundary Homomorphisms $`H_{n+1}^QH_n^D`$ In , from a split short exact sequence $`0C_n^\mathrm{D}(X)\stackrel{i}{}C_n^\mathrm{R}(X)\stackrel{j}{}C_n^\mathrm{Q}(X)0,`$ (7) the following homology long exact sequence $`\mathrm{}\stackrel{_{}}{}H_n^\mathrm{D}(X;G)\stackrel{i_{}}{}H_n^\mathrm{R}(X;G)\stackrel{j_{}}{}H_n^\mathrm{Q}(X;G)\stackrel{_{}}{}H_{n1}^\mathrm{D}(X;G)\mathrm{}`$ (8) was constructed. We give an application of Theorem 5.5 to boundary homomorphisms of this exact sequence. The endpoints of shadow colored arc diagrams are oriented by the orientation of the arcs. Each such endpoint represents a rack $`2`$-chain $`\pm (a,a)`$ where $`a`$ is the color on the surrounding region. The sum of these represents the image of the boundary map $`_{}:H_3^\mathrm{Q}(X)H_2^\mathrm{D}(X)`$ in the long exact homology sequence. Similarly, each oriented shadow colored branch point represents a rack $`3`$-chain $`\pm (a,b,b)`$ where $`a`$ is the color in one of the surrounding regions and $`b`$ is the color on the local surface. An oriented shadow colored hem $`2`$-crossing diagrams represent $`3`$-chains of the form $`(a,a,b)`$ where $`a`$ is the color on one side of the upper sheet, and $`b`$ is the color on the upper sheet. the image of the boundary map $`_{}:H_4^\mathrm{Q}(X)H_4^\mathrm{D}(X)`$ is represented as the sum over all branch points and hem $`2`$-crossings of the chains these represent. In Theorems 8.1 and 8.2, we will use these descriptions of the homomorphisms $`_{}`$ and geometric techniques to show that these boundary maps are trivial. ###### 8.1 Theorem . Let $`X`$ be a quandle. The boundary homomorphism $`_{}:H_3^Q(X)H_2^D(X)`$ in the long exact sequence of quandle homology is trivial. Proof. We give two proofs. (1) Let $`\eta =_{i=1}^kϵ_i(x_i,y_i,z_i)Z_3^Q(X)`$. Then $`\eta =_{j=1}^hϵ_j(w_j,w_j)`$ for some $`w_jX`$ for $`j=1,\mathrm{},h`$, from the definition of quandle cycle groups. Other terms coming out from each $`(x_i,y_i,z_i)`$ cancel out. Let $`\tau _i`$, $`i=1,\mathrm{},k`$, be the $`3`$-crossing diagrams colored with the triple $`(x_i,y_i,z_i)`$ for each $`i`$. According to the cancelation of the terms $`(x_i,y_i,z_i)`$, paste together $`\tau _i`$. There are boundary crossings (boundaries of $`\tau _i`$) that are colored by $`(w_j,w_j)`$. As the double curves form an immersed $`1`$-manifold with boundary, these boundary crossings are paired as $`_{g=1}^{h/2}[(a_g,a_g)(b_g,b_g)]`$, where the negative sign for $`b_g`$ represents that the orientation of the double curve points into the corresponding crossing. As one traces the double curve from the crossing with color $`(a_g,a_g)`$ to that with $`(b_g,b_g)`$, the curve under-goes the middle or top sheets at some of $`\tau _i`$’s. The color changes accordingly, but we have that $`a_g`$ and $`b_g`$ belong to the same orbit (as $`b_g=a_gw_g`$ for some word $`w_g`$ in $`X`$). On the other hand, one computes that $`(a,a,b)=(a,a)+(ab,ab)`$, so that $`(a_g,a_g)(b_g,b_g)B_2^D(X)`$, as desired. (2) Represent the homology class by a cycle $`\eta `$ and choose a shadow colored abstract $`1`$-knot diagram to represent $`\eta `$. Along each arc there are a collection of crossings. These represent the summands of $`\eta `$. We push the arc back along a parallel as indicated in Fig 26. When pushing backwards we make sure that the endpoint always passes under the crossings. Each $`2`$-crossing that is introduced in the process is shadow colored of the form $`(x,x,y)`$. Thus the new crossings do not affect the quandle homology class represented by the diagram. The shadow colored endpoints can be pairwise canceled. Thus $`_{}[\eta ]=0.`$ This completes the proof. $`\mathrm{}`$ ###### 8.2 Theorem. Let $`X`$ be a quandle. The boundary homomorphism $`_{}:H_4^\mathrm{Q}(X)H_3^\mathrm{D}(X)`$ in the long exact sequence of quandle homology is trivial. Proof. Consider a homology class $`[\eta ]H_4^\mathrm{Q}(X)`$ and represent it by the cycle $`\eta `$. Construct a shadow colored abstract $`2`$-knot diagram $`𝒮D_\eta =[f:MN]`$ with hems to represent $`\eta `$. We may assume (by attaching handles if necessary) that the ambient 3-manifold, $`N`$, is closed. Moreover, the boundary of the surface $`M`$ consist of simple closed curves formed by the hems in the diagram. The image $`_{}[\eta ]`$ is represented on $`𝒮D_\eta `$ as the collection of shadow colored branch point diagrams (terms of the form $`(a,b,b)`$) and hemmed $`2`$-crossing diagrams (terms of the form $`(a,a,b)`$). We first eliminate the terms of the form $`(a,a,b)`$. Consider the collection, $`B^1`$ of hems. These form a closed $`1`$-manifold in the $`3`$-manifold, $`N`$. The $`1`$-manifold $`B^1`$ is null homologous in $`N`$ since it is the boundary of the surface $`M`$ of $`𝒮D_\eta `$. Therefore $`B^1`$ bounds an embedded Seifert surface in $`N`$. The Seifert surface can be assumed to intersect the surface $`M`$ in general position. Now decompose the Seifert surface into $`1`$\- and $`2`$-handles where the $`1`$-handles are attached along a neighborhood of the hem. We will show that the hem can be eliminated by attaching the Seifert surface to the diagram $`𝒮D_\eta `$ along the hem. In the process, triple points will be introduced, but these will not affect the quandle homology class represented. The core disk of a $`1`$-handle of the Seifert surface intersects $`M`$ at a finite number of points. We push a segment of the hem along this core disk using the move depicted in (A) of Fig. 27 until segments are in a small ball neighborhood that contains no further sheet of the surface. The situation is depicted in (C) left in Fig. 27, and connect the segments as indicated in (C). Having pushed the $`B^1`$ and joined up segments, we may assume that each component of the hem is unknotted and the collection is unlinked. Then an unknotting disk for a component of $`B^1`$ intersects the double point set of $`𝒮D_\eta `$ in a finite number of points. We push the hem across these points introducing triple points, as indicated in Fig. 27 (D), that represent chains of the form $`(a,a,b,c)`$. (Recall that the hemmed sheet is always the under-sheet at each hemmed $`2`$-crossing diagram). But $`[\eta \pm (a,a,b,c)]=[\eta ]H_4^\mathrm{Q}(X)`$. The number of intersections with the resulting hem and $`M`$ is necessarily even since the hem is null homologous. We can push the components further, using the move depicted in Fig. 27 (B), until each bounds a disk that does not intersect the surface $`M`$. Now attach these disks to $`M`$ to obtain a shadow colored diagram of a closed surface that represents $`\eta `$. In this way we may assume that $`_{}[\eta ]`$ consists entirely of terms of the form $`(a,b,b)`$. The terms of the form $`(a,b,b)`$ are represented on the diagram by branch points. There are an even number of these since each arc of double points has two ends. Using the same technique as in , we can cancel these in pairs. The process of moving the surface involves quandle colored Roseman moves. The cancelation may introduce shadow colored triple points, but a consecutive pair of two of the four colors adjacent to the triple points introduced agree. Thus the triple points do not affect the quandle homology class that the diagrams represent. In this way $`\eta `$ is represented by a shadow colored diagram of a closed surface without branch points. Thus $`_{}[\eta ]=0`$. $`\mathrm{}`$ Acknowledgements. JSC is being supported by NSF grant DMS-9988107. MS is being supported by NSF grant DMS-9988101. SK is being supported by a Fellowship from the Japan Society for the Promotion of Science. We have had productive conversations with Dan Silver about this paper.
warning/0006/hep-ex0006006.html	ar5iv	text	# HEAVY FLAVOUR PHYSICS RESULTS FROM LEP 1 ## 1 $`𝒃`$-fragmentation studies Understanding the production of $`b`$-hadrons in $`Z`$ decays is important for many heavy flavour analyses. The $`b`$-quark hadronisation can be described in terms of the variable $`x_E=E_{b\text{-hadron}}/E_{\mathrm{beam}}`$, the fraction of the beam energy retained by the weakly-decaying $`b`$-hadron produced in a $`b`$ jet. Being not known accurately, the distribution of $`x_E`$ is simulated by the LEP experiments using the JETSET generator together with phenomenological models that relate the energy of the $`b`$-hadron with that of the initial $`b`$-quark. The most commonly used of these models, from Peterson et al., relies on a single parameter which has merely been tuned to reproduce the experimental spectra of high transverse momentum leptons originating mostly from the decay $`bc\mathrm{}^{}\overline{\nu }`$. This tuning corresponds to a mean $`x_E`$ of $`x_E=0.702\pm 0.008`$, which is the value recommended up to now for heavy flavour analyses at LEP. $`^\mathrm{?}`$ The ALEPH collaboration has submitted to this conference a new measurement $`^\mathrm{?}`$ of the shape of the $`x_E`$ distribution (see Fig. 2) based on approximately 3000 $`BD^{}l\nu `$ decays, where the $`B`$ meson energy has been estimated in a model-independent way from an identified lepton, a fully reconstructed $`D^{}`$ meson and missing energy information. The energy spectrum is found to be somewhat harder than assumed before, with $`x_E=0.7198\pm 0.0045_{\mathrm{stat}}\pm 0.0053_{\mathrm{syst}}`$, consistent with a previous ALEPH measurement. This confirms a recent result from SLD, $`^\mathrm{?}`$ $`x_E=0.714\pm 0.005_{\mathrm{stat}}\pm 0.007_{\mathrm{syst}}\pm 0.002_{\mathrm{model}}`$, which has small model-dependent systematics, although based on an inclusive sample of $`b`$-hadrons. These direct measurements of the shape of the $`x_E`$ distribution have now sufficient precision to envisage tests of the $`b`$-fragmentation model predictions and to discriminate amongst these models for the first time. For example, both ALEPH and SLD data favour the description of Kartvelishvili et al. over the one from Peterson et al. ## 2 Measurements of $`\mathbf{\|}𝑽_{𝒄𝒃}\mathbf{\|}`$ and $`\mathbf{\|}𝑽_{𝒖𝒃}\mathbf{\|}`$ The study of the $`B^0D^{}\mathrm{}^+\nu `$ decay kinematics allows the extraction of $`\|V_{cb}\|`$. The differential decay rate as a function of the boost $`\omega `$ of the $`D^{}`$ in the $`B^0`$ rest frame is predicted by the Heavy Quark Effective Theory to be $$\frac{d\mathrm{\Gamma }}{d\omega }=𝒦(\omega )_D^{}^2(\omega )\|V_{cb}\|^2,$$ (1) where $`𝒦(\omega )`$ is a known phase-space function and $`_D^{}(\omega )`$ a single form factor which, in the heavy quark limit, is equal to unity at zero recoil. The interesting observable is thus the decay rate at zero recoil, but since the phase space vanishes at $`\omega =1`$, the quantity $`_D^{}(1)\|V_{cb}\|`$ must be extracted from an extrapolation of the measured differential rate at $`\omega >1`$. At LEP, this extrapolation relies on a specific parametrization $`^\mathrm{?}`$ of the shape of $`_D^{}(\omega )`$ in terms of the slope $`\rho ^2`$ at $`\omega =1`$. The OPAL collaboration has recently updated their result obtained with fully reconstructed decays and performed a new analysis $`^\mathrm{?}`$ based on an inclusive $`D^{}`$ reconstruction relying on the identification of the slow pion from the $`D^{}`$ decay. Their new combined result is displayed on Fig. 2, together with earlier and similar measurements from ALEPH and DELPHI. A combination of these results, performed by the LEP $`V_{cb}`$ working group, $`^\mathrm{?}`$ takes into account all correlations and yields $$_D^{}(1)\|V_{cb}\|=(34.9\pm 0.7\pm 1.6)\times 10^3\text{and}\rho ^2=1.13\pm 0.08\pm 0.16.$$ (2) Systematic uncertainties are dominated by the limited knowledge of the $`D^{}`$ recoil spectrum for $`BD^{}\mathrm{}\nu X`$ background events. With $`_D^{}(1)=0.88\pm 0.05`$ from theoretical calculations taking into account finite quark masses and QCD corrections, $`^\mathrm{?}`$ this leads to the combined LEP estimate $`\|V_{cb}\|=(39.7\pm 0.8_{\mathrm{stat}}\pm 1.8_{\mathrm{syst}}\pm 2.2_{\mathrm{theory}})\times 10^3`$ from exclusive decays. Current theoretical calculations based on heavy quark symmetry relate $`\|V_{cb}\|`$ and $`\|V_{ub}\|`$ to the inclusive $`bc\mathrm{}^{}\overline{\nu }`$ and $`bu\mathrm{}^{}\overline{\nu }`$ decay widths, $$\|V_{cb}\|=0.0411\sqrt{\frac{\mathrm{BR}(bc\mathrm{}^{}\overline{\nu })}{0.105}}\sqrt{\frac{1.55\mathrm{ps}}{\tau _b}}\text{and}\|V_{ub}\|=0.00445\sqrt{\frac{\mathrm{BR}(bu\mathrm{}^{}\overline{\nu })}{0.002}}\sqrt{\frac{1.55\mathrm{ps}}{\tau _b}},$$ (3) with total uncertainties estimated to be $`5\%`$. $`^\mathrm{?}`$ While the measurements of the inclusive $`b\mathrm{}^{}`$ branching ratio and $`b`$-hadron lifetime are well established since several years, with current averages $`^\mathrm{?}`$ of $`\mathrm{BR}(b\mathrm{}^{})=(10.58\pm 0.07\pm 0.17)\%`$ and $`\tau _b=1.564\pm 0.014`$ ps, analyses measuring the $`bu\mathrm{}^{}\overline{\nu }`$ branching ratio are quite recent and unique to LEP. They face the difficulty of dealing with a very large $`bc\mathrm{}^{}\overline{\nu }`$ background, but have the advantage to be sensitive to the whole lepton spectrum (rather than only to the its end-point where the $`bc\mathrm{}^{}\overline{\nu }`$ decays are supressed). L3, ALEPH, and DELPHI have now all published evidence for $`bu\mathrm{}^{}\overline{\nu }`$ transitions, and their measurements average to $`^\mathrm{?}`$ $`\mathrm{BR}(bu\mathrm{}^{}\overline{\nu })=(1.67\pm 0.36\pm 0.37\pm 0.20)\times 10^3`$, where the first uncertainty summarizes statistics and experimental systematics, the second uncertainty reflects the limited knowledge of $`bc\mathrm{}^{}\overline{\nu }`$ transitions, and the third one results from the modelling of $`bu\mathrm{}^{}\overline{\nu }`$. Using Eq. 3, the LEP averages from inclusive semileptonic $`b`$ decays are $`\|V_{cb}\|=(40.8\pm 0.4_{\mathrm{exp}}\pm 2.0_{\mathrm{theory}})\times 10^3`$ and $`\|V_{ub}\|=(4.04_{0.74}^{+0.62})\times 10^3`$. The former can be combined $`^\mathrm{?}`$ with the less precise but consistent LEP estimate from exclusive decays to yield $`\|V_{cb}\|=(40.5\pm 1.8)\times 10^3`$. ## 3 Results on the $`𝑩_𝒔^\mathrm{𝟎}`$ decay width difference Information on $`\mathrm{\Delta }\mathrm{\Gamma }_s`$, the decay width difference between the two mass eigenstates of the $`B_s^0\text{}\overline{B}_s^0`$ system, can be obtained by studying the proper time distribution of untagged data samples enriched in $`B_s^0`$ mesons. In the case of an inclusive or a semileptonic $`B_s^0`$ decay selection, both the short- and long-lived components are present, and the proper time distribution is a superposition of two exponentials with decay constants $`\mathrm{\Gamma }_s\pm \mathrm{\Delta }\mathrm{\Gamma }_s/2`$. In principle, this provides sensitivity to both $`\mathrm{\Gamma }_s`$ and $`(\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s)^2`$. Ignoring $`\mathrm{\Delta }\mathrm{\Gamma }_s`$ and fitting for a single exponential leads to an estimate of $`\mathrm{\Gamma }_s`$ with a relative bias proportional to $`(\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s)^2`$. An alternative approach, which is directly sensitive to first order in $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$, is to determine the lifetime of $`B_s^0`$ candidates decaying to CP eigenstates; measurements exist for $`B{}_{s}{}^{0}J/\psi \varphi `$, $`^\mathrm{?}`$ and now also for $`B{}_{s}{}^{0}D_s^{()+}D_s^{()}`$, $`^\mathrm{?}`$ which are predicted to be mostly CP-even states. Recently, ALEPH has also obtained for the first time an estimate of $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ directly from a measurement of the $`B{}_{s}{}^{0}D_s^{()+}D_s^{()}`$ branching ratio, $`^\mathrm{?}`$ under the assumption that these decays practically account for all the CP-even final states. Figure 4 shows confidence contours in the plane $`(1/\mathrm{\Gamma }_s,\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s)`$ obtained from a combined likelihood built with all the available information from LEP and CDF, including dedicated $`\mathrm{\Delta }\mathrm{\Gamma }_s`$ studies as well as $`B_s^0`$ lifetime measurements. The corresponding results for $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ are $`^\mathrm{?}`$ $$\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s=0.24_{0.12}^{+0.16}\text{or}\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s<0.53\text{at 95\% CL}$$ (4) without external constraint, and $$\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s=0.17_{0.10}^{+0.09}\text{or}\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s<0.31\text{at 95\% CL}$$ (5) when constraining $`1/\mathrm{\Gamma }_s`$ to the current world average of the $`B^0`$ lifetime. Such a constraint is well motivated theoretically, since the $`B^0`$ and $`B_s^0`$ decay widths are predicted to differ by $`1\%`$ at most, but the current experimental check of this assumption, $`\tau _{B_s^0}/\tau _{B^0}=0.937\pm 0.040`$, $`^\mathrm{?}`$ is still of limited precision. These combined results on $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ are not yet precise enough to test the Standard Model predictions, which typically lie between 5% and 20%. ## 4 Search for $`𝑩_𝒔^\mathrm{𝟎}`$ oscillations $`B_s^0\text{}\overline{B}_s^0`$ oscillations have been the subject of many studies from ALEPH, DELPHI and OPAL, as well as SLD and CDF. No oscillation signal has been found so far. Because of the limited statistics available, the most sensitive analyses are currently the ones based on inclusive lepton samples, and on samples where a lepton and a $`D_s`$ meson have been reconstructed in the same jet. However, with larger samples, the most promising approach would be to use fully reconstructed $`B_s^0`$ mesons, which have a much better proper time resolution suitable to resolve higher oscillation frequencies. DELPHI $`^\mathrm{?}`$ have fully reconstructed 44 $`B_s^0`$ candidates in the $`\overline{D}^0K^{}\pi ^+`$, $`\overline{D}^0K^{}a_1^+`$, $`D_s^{()}\pi ^+`$ and $`D_s^{()}a_1^+`$, channels, whereas ALEPH have recently reported 50 candidates in the latter two channels. The number of signal events is estimated to be $`20`$ in each experiment, but with a proper time resolution of $`0.08`$ ps, more than two times better compared to more inclusive selections. As a result, these analyses, which have very poor sensitivity by themselves due to the lack of statistics, do nonetheless have a non-negligible impact on the average measurement of the oscillation amplitude $`𝒜`$ at high values of $`\mathrm{\Delta }m_s`$, the mass difference between the two mass eigenstates of the $`B_s^0`$ system. All results have been combined, including the latest ones from ALEPH $`^\mathrm{?}`$ released for this conference and based on $`D_s^{}\mathrm{}^+`$ correlations and fully reconstructed $`B_s^0`$ candidates, to yield the amplitudes $`𝒜`$ shown in Fig. 4 as a function of $`\mathrm{\Delta }m_s`$. In the combination procedure, $`^\mathrm{?}`$ the sensitivities of the inclusive lepton analyses, which depend directly on the assumed fraction $`f_{B_s^0}`$ of $`B_s^0`$ mesons in an unbiased sample of weakly-decaying $`b`$ hadrons, have been rescaled to a common value of $`f_{B_s^0}=0.096\pm 0.012`$. This value is obtained from direct production measurements, measurements of the time-integrated mixing probability $`\overline{\chi }`$ of $`b`$-hadrons at LEP, as well as the new world average of the $`B^0`$ oscillation frequency, $`\mathrm{\Delta }m_d=0.484\pm 0.015`$ ps. The combined sensitivity for 95% CL exclusion of $`\mathrm{\Delta }m_s`$ values is found to be 14.6 ps<sup>-1</sup>, which is also the actual limit below which all values of $`\mathrm{\Delta }m_s`$ are excluded by the data at 95% CL. No oscillation signal can be claimed based on the deviation from $`𝒜=0`$ seen in Fig. 4 around 17 ps<sup>-1</sup>: a fast Monte Carlo study $`^\mathrm{?}`$ shows indeed that statistical fluctuations can produced a more significant deviation (anywhere in the explored range in $`\mathrm{\Delta }m_s`$) in $`3\%`$ of the samples generated with a very large true value of $`\mathrm{\Delta }m_s`$. The information on $`\|V_{ts}\|`$ obtained, in the framework of the Standard Model, from the combined $`\mathrm{\Delta }m_s`$ limit is hampered by the hadronic uncertainty, as is the case when extracting $`\|V_{td}\|`$ from $`\mathrm{\Delta }m_d`$. However, many uncertainties cancel in the frequency ratio $$\frac{\mathrm{\Delta }m_s}{\mathrm{\Delta }m_d}=\frac{m_{B_s^0}}{m_{B^0}}\xi ^2\left\|\frac{V_{ts}}{V_{td}}\right\|^2,$$ (6) where $`\xi `$ is currently known to $`6\%`$ from lattice QCD. This relation can be used in fits of the CKM matrix, together with the experimental results on $`\mathrm{\Delta }m_s`$, $`\mathrm{\Delta }m_d`$, $`\|V_{ub}/V_{cb}\|`$ and $`ϵ_K`$, as well as theoretical inputs and unitarity constraints. Examples of such fits, $`^\mathrm{?}`$ shown in Fig. 5, illustrate the fact that the combined $`\mathrm{\Delta }m_s`$ results from Fig. 4 provide, together with the measured value of $`\mathrm{\Delta }m_d`$, a significant constraint on the CKM matrix, favouring positive values of the Wolfenstein parameter $`\rho `$. ## 5 CP violation in $`𝑩^\mathrm{𝟎}\mathbf{}𝑱\mathbf{/}𝝍𝑲_𝑺^\mathrm{𝟎}`$ decays ALEPH has recently released a new measurement $`^\mathrm{?}`$ of the CP asymmetry in $`B^0,\overline{B}{}_{}{}^{0}J/\psi K_S^0`$ decays, $$A(t)=\frac{N_{B^0}(t)N_{\overline{B}^0}(t)}{N_{B^0}(t)+N_{\overline{B}^0}(t)}=\mathrm{sin}(2\beta )\mathrm{sin}(\mathrm{\Delta }m_dt),$$ (7) where $`N_{B^0}(t)`$ and $`N_{\overline{B}^0}(t)`$ are the number of events produced as $`B^0`$ and $`\overline{B}^0`$ as a function of the proper time $`t`$, and $`\beta `$ is one of the angles of the CKM unitarity triangle. From a sample of 23 fully reconstructed candidates, selected with a signal efficiency of $`(28\pm 2)\%`$ and an estimated purity of $`(71\pm 12)\%`$, ALEPH measures $`\mathrm{sin}(2\beta )=0.93{}_{0.88}{}^{+0.64}\text{(stat)}{}_{0.24}{}^{+0.36}\text{(syst)}`$. The systematic uncertainty is dominated by the limited knowledge of the probability of mistagging the initial state, measured to be $`(25\pm 6)\%`$ using $`B^\pm J/\psi K^\pm `$ events. This $`\mathrm{sin}(2\beta )`$ result can be combined with previous measurements from OPAL and CDF to yield $`^\mathrm{?}`$ $`\mathrm{sin}(2\beta )=0.91\pm 0.35`$ or $`\mathrm{sin}(2\beta )>0`$ at 98.5% CL, increasing the confidence that CP violation has been observed in the $`B`$ sector. ## Acknowledgments I would like to thank D. Abbaneo and R. Forty for useful comments and careful reading of this writeup, A. Stocchi for providing Fig. 5, and the organizers of these “Rencontres” for an enjoyable conference. ## References
warning/0006/astro-ph0006246.html	ar5iv	text	# Limits on the cosmic-ray ionization rate toward massive young stars ## 1 Introduction The ionization fraction of molecular clouds is an important parameter for their dynamics through its control over the influence of any magnetic field. The ionization also has a major effect on the chemistry of molecular clouds because ion-neutral reactions are generally much faster than neutral-neutral reactions. In dense regions shielded from direct ultraviolet irradiation, the ionization is dominated by cosmic rays. However, the rate of this process $`\zeta _{\mathrm{CR}}`$ has not yet been constrained directly. The current best estimate comes from chemical models to reproduce the observed abundances of OH and HD in diffuse interstellar clouds (Hartquist et al. 1978; van Dishoeck & Black 1986; Federman et al. 1996), notably those toward Perseus OB2. These models indicate that $`\zeta _{\mathrm{CR}}`$$`=10^{16}10^{17}`$ s<sup>-1</sup> per H atom, but with a factor of 10 uncertainty because of uncertainties in temperature, radiation field and the effects of shocks. In addition, it is unknown if $`\zeta _{\mathrm{CR}}`$ varies with location in the Galaxy, since the diffuse cloud results are limited to the solar neighbourhood. Each cosmic-ray ionization of H<sub>2</sub> yields one H$`{}_{}{}^{+}{}_{3}{}^{}`$ molecule, so that H$`{}_{}{}^{+}{}_{3}{}^{}`$ has a constant concentration which depends only on $`\zeta _{\mathrm{CR}}`$ and the abundances of its main destroyers: CO, O and electrons. Hence, the recent detections of H$`{}_{}{}^{+}{}_{3}{}^{}`$ infrared absorption lines by Geballe & Oka (1996) and McCall et al. (1999) toward massive protostars provide a novel way to measure $`\zeta _{\mathrm{CR}}`$. We constrain $`\zeta _{\mathrm{CR}}`$ using models by van der Tak et al. (2000) of the envelopes of these stars, and use the derived values to model observations of HCO<sup>+</sup>, the abundance of which is also proportional to $`\zeta _{\mathrm{CR}}`$. ## 2 Models McCall et al. (1999) present observations of rovibrational lines of H$`{}_{}{}^{+}{}_{3}{}^{}`$ in absorption against seven luminous ($`10^410^5`$ L) young stars, which are still embedded in envelopes of $`100`$ M of dust and molecular gas. These same sources have been studied by Mitchell et al. (1990) in <sup>13</sup>CO infrared absorption and by van der Tak et al. (2000) in submillimeter dust continuum and CS, C<sup>34</sup>S and C<sup>17</sup>O line emission. Based on these data sets, van der Tak et al. (1999, 2000) modeled the temperature and density structure of the envelopes using a power law structure $`n=n_0(r/r_0)^\alpha `$. The radial dust temperature profile is calculated self-consistently from the luminosity and $`n_0`$ is determined from submillimeter photometry which probes the dust column density. The parameter $`\alpha `$ is constrained by modeling the relative strengths of the CS and C<sup>34</sup>S $`J=21`$ through $`109`$ lines with a non-LTE radiative transfer program based on the Monte Carlo method. The outer radii of the models are twice the half-intensity radii of the CS $`J=54`$ emission, given in Table 5 of van der Tak et al. (2000). For the sources studied here, the values are $`(311)\times 10^{17}`$ cm, which are accurate to a factor of 2. However, the dust and gas maps also reveal emission extending outside the envelopes of W 3 IRS5 and GL 490, and maybe W 33A. No “skin” appears to surround GL 2591 and GL 2136. The case of S 140 is complicated: the dust appears to be heated by multiple sources, and its emission is not well fitted by a centrally heated model. We do not have a dust map of NGC 2264, but this source is part of an extended molecular cloud complex, so that extended material may also contribute to the H$`{}_{}{}^{+}{}_{3}{}^{}`$ absorption. Evidence for extended components at lower temperature and/or column density than those probed by the dust emission comes from emission in low-$`J`$ lines of CO at $`\genfrac{}{}{0pt}{}{_>}{^{}}1^{}`$ offsets, and self-absorptions on their central line profiles. These features are present in the data for all the sources discussed in this paper. Before considering foreground contributions in § 5, we concentrate on the dense molecular envelopes. ## 3 Results Given the temperature and density profiles, we calculate the H$`{}_{}{}^{+}{}_{3}{}^{}`$ concentration at each position in the envelopes. Considering only cosmic rays as producers of H$`{}_{}{}^{+}{}_{3}{}^{}`$ and reactions with CO and O as destroyers, the concentration of H$`{}_{}{}^{+}{}_{3}{}^{}`$ is given by $`n`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$)= $`\zeta _{\mathrm{CR}}`$/\[$`x`$(CO)$`\times k_{\mathrm{CO}}`$ \+ $`x`$(O)$`\times k_\mathrm{O}`$\]. In this expression, $`x`$(CO) and $`x`$(O) are the abundances of CO and O relative to H<sub>2</sub>, and $`k_{\mathrm{CO}}`$ and $`k_\mathrm{O}`$ the rate coefficients for their respective reactions with H$`{}_{}{}^{+}{}_{3}{}^{}`$, taken from Millar et al. (1997). We neglect any dependence of $`k_{\mathrm{CO}}`$ on temperature since the dipole moment of CO is small. The models use an abundance of CO of $`2\times 10^4`$ at temperatures above $`20`$ K, and zero below due to freeze-out on dust grains. This abundance behaviour is consistent with observations of C<sup>17</sup>O emission lines by van der Tak et al. (2000). The abundance of O is assumed to be $`1.5\times 10^4`$ based on the models of Lee et al. (1996), and the temperature in our models does not drop below 14 K, where O would freeze out. The ortho/para ($`o`$/$`p`$) or $`(J,K)=(1,0)/(1,1)`$ ratio of H$`{}_{}{}^{+}{}_{3}{}^{}`$ changes with radius since the ground state of ortho-H$`{}_{}{}^{+}{}_{3}{}^{}`$ lies $`32.86`$ K above that of para-H$`{}_{}{}^{+}{}_{3}{}^{}`$ (Dinelli et al. (1997)), and reactive collisions with H<sub>2</sub> tie the $`o`$/$`p`$ ratio to the kinetic temperature. Figure 1 illustrates the results for the case of GL 2136. Integration of the concentrations of ortho- and para-H$`{}_{}{}^{+}{}_{3}{}^{}`$ over radius yields total column densities $`N`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) and mean $`o`$/$`p`$ ratios, which are compared with the data in Table 1. The calculated $`o`$/$`p`$ ratios are consistent with the data within the observational errors, but the model values are systematically lower than those observed. For CO, the models, which were constrained by emission data, typically overproduce absorption measurements of $`N`$(CO) by factors of 3, probably due to deviations from spherical geometry on small scales, consistent with several other tracers (van der Tak et al. 2000). Since the model values of $`N`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) may be less affected because H$`{}_{}{}^{+}{}_{3}{}^{}`$ is more evenly distributed than CO (Fig. 1), Table 1 presents the values of $`\zeta _{\mathrm{CR}}`$ both before (case 1) and after (case 2) scaling the model down by the ratio of observed to modeled $`N`$(CO). The uncertainty in the model is a factor of two due to the uncertain radii of the envelopes. The estimates of $`\zeta _{\mathrm{CR}}`$ in case (2) are considerably larger than those in previous work (§ 1), which together with the low $`o`$/$`p`$ ratios indicates that there may be an additional component of warm H$`{}_{}{}^{+}{}_{3}{}^{}`$ along the line of sight. We will estimate the contribution to $`N`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) by the dense envelopes by modeling emission lines of H<sup>13</sup>CO<sup>+</sup> which have critical densities of $`10^6`$ cm<sup>-3</sup>, and hence cannot arise in the foreground. ## 4 Comparison with HCO<sup>+</sup> In the dense envelopes, the main destruction route of H$`{}_{}{}^{+}{}_{3}{}^{}`$ is the reaction with CO into HCO<sup>+</sup>. The concentration of HCO<sup>+</sup> is given by $`n`$(HCO<sup>+</sup>)=$`x`$(CO)$`n`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$)$`k_{\mathrm{CO}}`$/ \[$`(x(e)k_e`$ \+ $`x`$(H<sub>2</sub>O)$`k_{\mathrm{H}_2\mathrm{O}}`$\], with $`k_e`$ the rate coefficient for dissociative recombination of HCO<sup>+</sup>. The electron fraction $`x(e)`$ has been calculated at each point in the envelopes with a small chemical network (cf. de Boisanger et al. 1996) based on the UMIST reaction rates (Millar et al. 1997). The main difference with the analysis of de Boisanger et al. (1996) is the use of a detailed physical structure to interpret the high-excitation lines. We assume that O<sub>2</sub> and H<sub>2</sub>O have negligible ($`\genfrac{}{}{0pt}{}{_<}{^{}}10^6`$) abundances in the bulk of the envelopes, but that at $`T>100`$ K, $`x`$(H<sub>2</sub>O) jumps to $`5\times 10^5`$ due to grain mantle evaporation. We neglect metals such as Mg, Fe and S as contributors to $`x(e)`$ and large molecules such as polycyclic aromatic hydrocarbons as sinks of $`x(e)`$; using the low metal abundances inferred from dark cloud chemistry models would increase $`x(e)`$ by a factor of 2–3 (Lee et al. 1996). The values of $`\zeta _{\mathrm{CR}}`$ derived above give $`x(e)10^7`$ at the outer radii and $`10^9`$ at the inner radii, as illustrated in Fig. 1. The precipitous drop of HCO<sup>+</sup> at $`100`$ K, caused by reactions with evaporated water, occurs at too small radii to affect our results. In the comparison with data, we use the $`60\times `$ less abundant isotope H<sup>13</sup>CO<sup>+</sup> to avoid optical depth effects. The maximum optical depth in the lines is $`1`$ in our models. Table 2 lists the calculated fluxes of the H<sup>13</sup>CO<sup>+</sup> $`J`$=3$``$2 and 4$``$3 lines in $`18^{\prime \prime }`$ and $`14^{\prime \prime }`$ beams. Observations are from van der Tak et al. (1999) for GL 2591 and from de Boisanger et al. (1996) for NGC 2264 and W 3 IRS5. The data for W 33A, GL 490, S 140 and GL 2136 were obtained with the James Clerk Maxwell Telescope in the way described in van der Tak et al. (2000). Using $`\zeta _{\mathrm{CR}}`$ derived from H$`{}_{}{}^{+}{}_{3}{}^{}`$, the models overproduce HCO<sup>+</sup> by factors of $`27`$. Adjusting the models to the H<sup>13</sup>CO<sup>+</sup> data yields refined estimates for $`\zeta _{\mathrm{CR}}`$ (Table 2) which pertain strictly to the dense molecular gas, unaffected by any intervening clouds along the line of sight. The data for the various sources span the range of $`\zeta _{\mathrm{CR}}`$$`=(2.6\pm 1.8)\times 10^{17}`$ s<sup>-1</sup>, in good agreement with the diffuse cloud estimates (§ 1), and also consistent with recent data from the Voyager and Pioneer spacecraft at distances up to $`60`$ AU from the Sun (Webber (1998)). Figure 2 shows that the source-to-source variation in $`\zeta _{\mathrm{CR}}`$ is not related to Galactic structure through differences in cosmic-ray flux, nor to shielding against cosmic rays at high H<sub>2</sub> column densities. The values of $`\zeta _{\mathrm{CR}}`$ are also unrelated to luminosity, which implies that local ionization such as by X-rays (Maloney et al. (1996)) is unimportant on the scales traced by our data. Variations of the cosmic-ray density by $`50`$% on scales of a few kpc are in good agreement with results from $`\gamma `$ray observations (e.g., Hunter et al. 1997). However, why does H$`{}_{}{}^{+}{}_{3}{}^{}`$ give systematically higher values of $`\zeta _{\mathrm{CR}}`$? ## 5 Contributions by foreground layers Figure 3 plots the observed $`N`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) versus heliocentric distance, including all data from McCall et al. (1999) as well as the results for the Galactic Center and the diffuse cloud in front of Cyg OB2 #12 from Geballe et al. (1999). The dense cloud data have a correlation coefficient of 93%, suggesting that absorption by intervening clouds plays an important role for the more distant sources. This section investigates the possible nature of these absorbers. First, the absorptions may occur at the edges of the dense cores studied here, where carbon is in neutral or ionized form. This “photodissociation region” occupies $``$3–4 magnitudes of visual extinction (Hollenbach & Tielens (1997)), corresponding to $`N_H\genfrac{}{}{0pt}{}{_<}{^{}}8\times 10^{21}`$ cm<sup>-2</sup>. The ionized layer is negligible because the high electron fraction ($`10^4`$) limits $`n`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) to $`10^7`$ cm<sup>-3</sup>. For the neutral component, assuming $`n10^4`$ cm<sup>-3</sup> as derived specifically for S 140 by Timmermann et al. (1996), and $`n`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) $`10^4`$ cm<sup>-3</sup>, we find $`N`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) $``$ few $`\times 10^{13}`$ cm<sup>-2</sup>, comparable to the dense envelopes. Second, the absorbers may consist of cold ($`\genfrac{}{}{0pt}{}{_<}{^{}}20`$ K) molecular gas. For $`n\genfrac{}{}{0pt}{}{_>}{^{}}10^4`$ cm<sup>-3</sup>, the CO will be frozen out on the grains. Tielens et al. (1991) observed solid CO in absorption toward all our sources and found $`N`$(CO)$`10^{17}`$ cm<sup>-2</sup>, or $`N_\mathrm{H}10^{21}`$ cm<sup>-2</sup>, assuming that most carbon is in solid CO. The implied column lengths are too short to be of importance for H$`{}_{}{}^{+}{}_{3}{}^{}`$, and the low temperatures are incompatible with the observed $`o`$/$`p`$ H$`{}_{}{}^{+}{}_{3}{}^{}`$ ratios. Third, the H$`{}_{}{}^{+}{}_{3}{}^{}`$ absorptions may arise in clouds with $`n\genfrac{}{}{0pt}{}{_<}{^{}}10^4`$ cm<sup>-3</sup>, which either surround the power-law envelopes or happen to lie along the line of sight. Such “translucent” foregrounds are visible in our data (§ 2), and can contribute $`N`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) $`10^{14}`$ cm<sup>-2</sup> each based on models by van Dishoeck & Black (1989). These tenuous clouds have long path lengths and may dominate the H$`{}_{}{}^{+}{}_{3}{}^{}`$ absorption. At low densities, HCO<sup>+</sup> may form through OH + C<sup>+</sup> $``$ H + CO<sup>+</sup> followed by H<sub>2</sub> \+ CO$`{}_{}{}^{+}`$ H + HCO<sup>+</sup>. However, for our sources, \[CII\] $`158`$ $`\mu `$m data indicate $`N`$(C<sup>+</sup>)/$`N`$(CO) $`\genfrac{}{}{0pt}{}{_<}{^{}}10^2`$. Translucent clouds are generally weak in HCO<sup>+</sup> emission (Gredel et al. (1994)). The velocities of the H$`{}_{}{}^{+}{}_{3}{}^{}`$ absorptions are consistent with those of the submillimeter emission lines of C<sup>17</sup>O and C<sup>34</sup>S, suggesting that the H$`{}_{}{}^{+}{}_{3}{}^{}`$ absorbers are in the vicinity of the infrared sources. However, the correlation of $`N`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) with distance remains after subtracting the dense core contribution (Table 2), suggesting a non-local origin. Altogether, the data indicate that the contribution of the envelopes to $`N`$(H$`{}_{}{}^{+}{}_{3}{}^{}`$) varies from $`10^{13}`$ to $`10^{14}`$ cm<sup>-2</sup>, and that any additional absorption seen in sources at $`d>2`$ kpc occurs in intervening clouds. In summary, observations of H$`{}_{}{}^{+}{}_{3}{}^{}`$ absorption and H<sup>13</sup>CO<sup>+</sup> emission lines, combined with models of the temperature and density structure of the sources, constrain the cosmic-ray ionization rate to $`\zeta _{\mathrm{CR}}`$$`=(2.6\pm 1.8)\times 10^{17}`$ s<sup>-1</sup>, with upper limits that are factors of 3–5 higher. Future tests of the results include more sensitive observations of H$`{}_{}{}^{+}{}_{3}{}^{}`$ toward W 3 IRS5 and NGC 2264, velocity-resolved observations to search for H$`{}_{}{}^{+}{}_{3}{}^{}`$ absorption at offset velocities from the dense cores, observations of more distant sources to test the correlation with distance, and observations of HCO<sup>+</sup> infrared absorption lines to directly compare with H$`{}_{}{}^{+}{}_{3}{}^{}`$ and CO infrared absorption. ###### Acknowledgements. We thank Neal Evans, Dan Jaffe, Tom Geballe and John Black for useful discussions. This research is supported by NWO grant 614-41-003.
warning/0006/quant-ph0006047.html	ar5iv	text	# Quantum Disentanglers ## I Introduction Information encoded in qubits can be used for reliable quantum communication or efficient quantum computing . This information is encoded in a quantum state $`\|\psi (\vartheta ,\phi )`$ which in the case of a qubit can be parameterized as $`\|\psi (\vartheta ,\phi )=\mathrm{cos}{\displaystyle \frac{\vartheta }{2}}\|0+\mathrm{e}^{i\phi }\mathrm{sin}{\displaystyle \frac{\vartheta }{2}}\|1;`$ (1) where $`\|0`$ and $`\|1`$ are basis vectors of the 2-dimensional space of the qubit and $`0\vartheta \pi `$; $`0\phi 2\pi `$. Qubits are very fragile, that is the state of a qubit can easily be changed by the influence of the environment or a random error. One (very inefficient) way to protect the quantum information encoded in a qubit is to measure it. With the help of an optimal measurement one can estimate the state of a qubit, with an average fidelity equal to 2/3 (see below). In this way a quantum information is transformed into a classical information which can be stored, copied, and processed according the laws of classical physics with arbitrarily high precision. However, in order to utilize the full potential of quantum information processing we have to keep the information in states of quantum systems, but then we are forced to face the problem of decoherence. Recently it has been proposed that quantum information and quantum information processing can be stabilized via symmetrization . In particular, the qubit in an unknown state is entangled with a set of $`N1`$ (ancilla) qubits in a specific reference state (let us say $`\|0`$) so the symmetric state $`\|\mathrm{\Psi }`$ of $`N`$ qubits, $`\|\mathrm{\Psi }\left(\|\psi ,0\mathrm{},0+\|0,\psi ,\mathrm{},0+\mathrm{}+\|0,0,\mathrm{}\psi \right),`$ (2) is generated. If we introduce a notation for completely symmetric states $`\|N;l`$ of $`N`$ qubits with $`l`$ of them being in the state $`\|1`$ and $`Nl`$ of them in the state $`\|0`$, then the state (2) can be expressed in the simple form $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })=\mathrm{cos}{\displaystyle \frac{\overline{\vartheta }}{2}}\|N;0+\mathrm{e}^{i\overline{\phi }}\mathrm{sin}{\displaystyle \frac{\overline{\vartheta }}{2}}\|N;1`$ (3) where the parameters $`\overline{\vartheta }`$ and $`\overline{\phi }`$ are specified by the relations $`\mathrm{cos}{\displaystyle \frac{\overline{\vartheta }}{2}}={\displaystyle \frac{\sqrt{N}\mathrm{cos}\frac{\vartheta }{2}}{\sqrt{\mathrm{sin}^2\frac{\vartheta }{2}+N\mathrm{cos}^2\frac{\vartheta }{2}}}};`$ (4) and $`\mathrm{sin}\frac{\overline{\vartheta }}{2}=\sqrt{1\mathrm{cos}^2\frac{\overline{\vartheta }}{2}}`$, while $`\overline{\phi }=\phi `$. We see that symmetric $`N`$ qubit state $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })`$ is isomorphic to a single qubit state. But in this case the information is spread among $`N`$ entangled qubits - the original quantum information is “diluted”. Each of the qubits of the $`N`$-qubit state (3) is in the state $`\rho _j=\frac{N1}{N}\|00\|+\frac{(1\sqrt{N})}{N}(\mathrm{cos}^2\frac{\overline{\vartheta }}{2}\|00\|+\mathrm{sin}^2\frac{\overline{\vartheta }}{2}\|11\|)+\frac{1}{\sqrt{N}}\|\psi (\overline{\vartheta },\overline{\phi })\psi (\overline{\vartheta },\overline{\phi })\|`$. We define the average fidelity between the single state $`\rho _j`$ and the original qubit $`\|\psi (\vartheta ;\phi )`$ as $`\overline{}={\displaystyle 𝑑\mathrm{\Omega }\psi (\vartheta ;\phi )\|\rho _j(\overline{\vartheta },\overline{\phi })\|\psi (\vartheta ;\phi )}`$ (5) where $`d\mathrm{\Omega }=\mathrm{sin}\vartheta d\vartheta d\phi /4\pi `$ is the invariant measure on the state space of the original qubit (i.e. we assume no prior knowledge about the pure state $`\|\psi (\vartheta ;\phi )`$). For this fidelity we find the expression $`\overline{}_0={\displaystyle \frac{N^212\mathrm{ln}N}{2(N1)^2}}.`$ (6) We see that for $`N=1`$ the fidelity $`\overline{}_0`$ is equal to unity (as it should, because in this case $`\|\mathrm{\Psi }=\|\psi `$) while in the limit $`N\mathrm{}`$ we find $`\overline{}=1/2`$. In fact in this limit density operators of individual qubits are approximately equal to $`\|00\|`$. In other words, individually the qubits of the symmetric state $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })`$ in the large $`N`$ limit do not carry any information about the original single-qubit state $`\|\psi `$. So how can we extract the information from the $`N`$-qubit symmetric state (3)? The ideal possibility would be to have have a perfect universal disentangler which would perform a unitary transformation $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })\|\mathrm{\Psi }_{ideal}\|N1;0\|\psi (\vartheta ,\phi ).`$ (7) But quantum mechanics does not allow this type of disentangling transformation . While the perfect transformation is impossible, there are a number of things we can do to concentrate the information from the $`N`$-qubit state $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })`$ back into a single qubit. In principle, we have the following possibilities: i) We can either optimally measure the $`N`$ qubit state and based on the information obtained prepare a single-qubit state. ii) We can design a quantum disentangler which would perform a transformation as close as possible to the ideal disentangling (7). In this quantum scenario we have several options \- the process of disentanglement can be input-state dependent. This means that states (3) for some values of the parameters $`\overline{\vartheta }`$ and $`\overline{\phi }`$ will be disentangled better than for other values of these parameters. Alternatively, we can construct a quantum device which disentangles all the state with the same fidelity. iii) Finally, we propose a probabilistic disentangler, such that when a specific projective measurement over an ancilla is performed at the output, the desired single-qubit state is generated. The probability of the outcome of the measurement in this case is state-dependent. In what follows we shall investigate all these possibilities. Before proceeding we note that a different type of disentangler has been considered by Terno and Mor \- . They considered two different operations. The first would take the state of a bipartite quantum system and transform it into a state that is just the product of the reduced density matrixes of the two subsystems. The second, which is a generalization of the first, would again start with a state of a bipartate quantum system, and map it into a separable state which has the same reduced density matrixes as the original state. They showed that while both of these processes are impossible in general, they can be realized for particular sets of input states. An approximate disentangler of the first type has been considered by Bandyopadhyay, et. al. . The disentanglers we are considering extract, to some degree of approximation, an unknown state from an entangled state formed from that state and a known state. ## II Measurement scenario Here we first describe a measurement scenario utilizing a set of specific projection operators. Then we present the optimal measurement-based approach to quantum disentanglement and we derive an upper bound on the fidelity of the measurement-based disentangler. We utilize the fact that the $`N`$ qubit system prepared in the state $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })`$ is isomorphic to a single qubit. Therefore we first consider a strategy based on a a projective measurement with two projectors $`P_j(\vartheta ^{},\phi ^{})=\|\mathrm{\Xi }_j(\vartheta ^{},\phi ^{})\mathrm{\Xi }_j(\vartheta ^{},\phi ^{})\|`$ ($`j=0,1`$) with $`\|\mathrm{\Xi }_0(\vartheta ^{},\phi ^{})`$ $`=`$ $`\mathrm{cos}{\displaystyle \frac{\vartheta ^{}}{2}}\|N;0+\mathrm{e}^{i\phi ^{}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{}}{2}}\|N;1;`$ (8) $`\|\mathrm{\Xi }_1(\vartheta ^{},\phi ^{})`$ $`=`$ $`\mathrm{e}^{i\phi ^{}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{}}{2}}\|N;0;\mathrm{cos}{\displaystyle \frac{\vartheta ^{}}{2}}\|N;1,`$ (9) such that $`\mathrm{\Xi }_j(\vartheta ^{},\phi ^{})\|\mathrm{\Xi }_k(\vartheta ^{},\phi ^{})=\delta _{j,k}`$ and $`_jP_j(\vartheta ^{},\phi ^{})=𝟙`$, where the angles $`\vartheta ^{}`$ and $`\phi ^{}`$ are chosen randomly if no prior information about the measured $`N`$-qubit state is available. We can use the result of the measurement to manufacture a a single-qubit state. Specifically, if the result of the measurement is positive for $`P_0`$ then the single qubit is prepared in the state $`\|\eta _0(\vartheta ^{},\phi ^{})=\mathrm{cos}{\displaystyle \frac{\vartheta ^{}}{2}}\|0+\mathrm{e}^{i\phi ^{}}\mathrm{sin}{\displaystyle \frac{\vartheta ^{}}{2}}\|1,`$ (10) while if the output is positive for $`P_1`$ then the single qubit is prepared in the orthogonal state $`\|\eta _1(\vartheta ^{},\phi ^{})`$ . For a particular orientation of the measurement apparatus (i.e. the angles $`\vartheta ^{},\phi ^{}`$) this measurement-based scenario gives us a single qubit prepared in the state described by the density operator $`\rho ^{(meas)}(\overline{\vartheta },\overline{\phi };\vartheta ^{},\phi ^{})={\displaystyle \underset{j=0}{\overset{1}{}}}\left\|\mathrm{\Psi }\|\mathrm{\Xi }_j\right\|^2\|\eta _j\eta _j\|`$ (11) After we average over all possible orientations of the measurement apparatus we obtain on average a single qubit prepared in the state $`\rho ^{(est)}(\overline{\vartheta },\overline{\phi })={\displaystyle \frac{1}{3}}\|\psi (\overline{\vartheta },\overline{\phi })\psi (\overline{\vartheta },\overline{\phi })\|+{\displaystyle \frac{1}{3}}𝟙.`$ (12) To find the average fidelity of this measurement-based disentangling procedure we have to evaluate the mean fidelity $`\overline{}_1`$, that is the overlap between the state (12) and the original input state $`\|\psi (\vartheta ,\phi )`$ averaged over all possible orientations of the input qubit: $`\overline{}_1={\displaystyle 𝑑\mathrm{\Omega }\psi (\vartheta ,\phi )\|\rho ^{(est)}(\overline{\vartheta },\overline{\phi })\|\psi (\vartheta ,\phi )}.`$ (13) Taking into account the relation (4) we perform the integration in Eq.(13) and we find $`\overline{}_1={\displaystyle \frac{1}{3}}(1+f_N)`$ (14) where the function $`f_N`$ reads $`f_N={\displaystyle \frac{N^2+4N^{3/2}4N^{1/2}1+2N\mathrm{ln}N}{2(N1)(N^{1/2}+1)^2}}.`$ (15) For $`N=1`$: $`\overline{}_1=2/3`$ which is the optimal fidelity of estimation of the state of a single qubit. From Fig. 1 we see that the fidelity (14) is a decreasing function of $`N`$ and in the limit $`N\mathrm{}`$ we find $`\overline{}_1=1/2`$, which is equal to the fidelity of a random guess associated with a binary system such as the two projectors under consideration. In other words, when the original qubit is diluted into an infinite qubit state of the form (3) no relevant information can be gained from the measurement. The estimated density operator (12) in this case is simply equal to $`𝟙/\mathrm{𝟚}`$, which is understandable, because as we have shown earlier in this limit the $`N`$-qubit state is approximately in the state $`\|N,0`$, so information about the original is “almost” totally lost. ### A Optimal measurement scenario We now want to find an upper bound $`\overline{}^{max}`$ for the average fidelity which can be achieved by a wide class of measurement-based disentanglement procedures. We assume that it is a priori known that our $`N`$-qubit is prepared in the symmetric state (2) with unknown parameters $`\vartheta `$ and $`\phi `$ associated with a single-qubit state (1). The integration measure on the state space of the single qubit is $`d\mathrm{\Omega }=\frac{1}{4\pi }\mathrm{sin}\vartheta d\vartheta d\phi `$ and the corresponding prior probability density distribution on this state space is constant. Our strategy is to measure the input state $`\|\mathrm{\Psi }`$ along the vector $`\|\mathrm{\Xi }_0`$ \[see Eq. (9)\], where the angles $`\vartheta ^{}`$ and $`\phi ^{}`$ are chosen according to the distribution $`q(\vartheta ^{},\phi ^{})`$, which will be left unspecified for the moment. If the answer is positive, we produce the output density matrix $`\rho _0(\vartheta ^{},\phi ^{})`$, and if it is negative we produce $`\rho _1(\vartheta ^{},\phi ^{})`$, where $$\rho _j(\vartheta ^{},\phi ^{})=𝑑\mathrm{\Omega }^{\prime \prime }p_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime }\|\vartheta ^{},\phi ^{})\|\eta (\vartheta ^{\prime \prime },\phi ^{\prime \prime })\eta (\vartheta ^{\prime \prime },\phi ^{\prime \prime })\|$$ (16) with $`j=0,1`$ and $`\|\eta `$ given by Eq. (10). We shall also leave the conditional probabilities, $`p_j`$ unspecified, as this allows us to consider a wide range of strategies. For a fixed $`\|\mathrm{\Xi }_0`$, the probability of the output being $`\rho _0(\vartheta ^{},\phi ^{})`$ is $`\|\mathrm{\Xi }_0\|\mathrm{\Psi }\|^2`$ and the probability of it being $`\rho _1(\vartheta ^{},\phi ^{})`$ is $`\|\mathrm{\Xi }_1\|\mathrm{\Psi }\|^2`$. Averaging over all vectors, $`\|\mathrm{\Xi }`$ gives us $`\rho ^{(out)}(\overline{\vartheta },\overline{\phi })`$ $`=`$ $`{\displaystyle }d\mathrm{\Omega }^{}[\|\mathrm{\Xi }_0\|\mathrm{\Psi }\|^2\rho _0(\vartheta ^{},\phi ^{})`$ (17) $`+`$ $`\|\mathrm{\Xi }_1\|\mathrm{\Psi }\|^2\rho _1(\vartheta ^{},\phi ^{})]q(\vartheta ^{},\phi ^{}).`$ (18) In order to find the average fidelity of the output produced by this procedure, we compute the fidelity for a particular input state and average over the input ensemble $$\overline{}=𝑑\mathrm{\Omega }\psi (\vartheta ,\phi )\|\rho ^{(out)}(\overline{\vartheta },\overline{\phi })\|\psi (\vartheta ,\phi ),$$ (19) where $`\overline{\vartheta }`$ is a function of $`\vartheta `$ \[see Eq.(4)\]. This can be expressed as $$\overline{}=𝑑\mathrm{\Omega }^{}𝑑\mathrm{\Omega }^{\prime \prime }\underset{j=0}{\overset{1}{}}P_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{})f_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{}),$$ (20) where $$P_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{})=p_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime }\|\vartheta ^{},\phi ^{})q(\vartheta ^{},\phi ^{}),$$ (21) is a normalized joint probability distribution, and $`f_0`$ $`=`$ $`{\displaystyle 𝑑\mathrm{\Omega }\|\mathrm{\Psi }\|\mathrm{\Xi }_0\|^2\|\psi \|\eta \|^2}`$ (22) $`f_1`$ $`=`$ $`{\displaystyle 𝑑\mathrm{\Omega }\|\mathrm{\Psi }\|\mathrm{\Xi }_1\|^2\|\psi \|\eta \|^2}.`$ (23) We first note that $$𝑑\mathrm{\Omega }^{\prime \prime }p_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{})f_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{})h_j(\vartheta ^{},\phi ^{}),$$ (24) where $$h_j(\vartheta ^{},\phi ^{})=supf_j(\vartheta ^{\prime \prime },\phi ^{\prime \prime };\vartheta ^{},\phi ^{}),$$ (25) and the supremum is taken over the variables $`\vartheta ^{\prime \prime },\phi ^{\prime \prime }`$. We then have that $$\overline{}\overline{}^{max}=sup[h_0(\vartheta ^{},\phi ^{})+h_1(\vartheta ^{},\phi ^{})],$$ (26) where the supremum is now taken over $`0\vartheta ^{}\pi `$ and $`0\phi ^{}<2\pi `$. In order to calculate this upper bound we must find explicit expressions for $`f_0`$ and $`f_1`$. After performing the necessary calculations we find for $`\overline{}^{max}`$ the expression $`\overline{}^{max}={\displaystyle \frac{1}{2}}\left[1+{\displaystyle \frac{\sqrt{N}}{(N1)^3}}(N^212N\mathrm{ln}N)\right].`$ (27) This fidelity for $`N=1`$ is equal to 2/3 while in the limit $`N\mathrm{}`$ is equal to 1/2. For any other $`N`$ is larger than the fidelity $`\overline{}_1`$ of the measurement given by Eq.(14) as discussed in our previous example. Nevertheless, as we will show later it is alway smaller than the fidelity of the universal quantum device. ## III Quantum scenario In what follows we show that a quantum disentangler which preserves quantum coherences can distill the information back to a single qubit more efficiently than can the measurement-based method. As we have already said in the introduction quantum mechanics does not allow one to construct a perfect disentangler which would perform transformation (7) for an arbitrary (unknown) state $`\|\psi (\vartheta ,\phi )`$ diluted in the $`N`$ qubit symmetric state (3). Nevertheless, we can try to design optimal disentanglers which perform best under given constraints. ### A State-independent devices So let us assume our quantum disentangler, $`D`$, is a quantum system with a $`K`$-dimensional Hilbert space spanned by basis vectors $`\|d_k`$ ($`k=1,\mathrm{},K`$). The disentangler is always initially prepared in the state $`\|d_0`$, and then it interacts with the $`N`$-qubit system in the state (3). At the output we want to disentangle the $`N1`$ ancilla qubits from the original qubit, so we expect to have $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })\|d_0\|N1;0{\displaystyle \underset{k=1}{\overset{K}{}}}{\displaystyle \underset{j=0}{\overset{1}{}}}c_j(\overline{\vartheta },\overline{\phi })\|j\|d_k.`$ (28) As seen from Eq.(28) during the disentanglement process the entanglement between the $`N1`$ ancilla qubits and the original qubit is transferred (swapped) into the entanglement between the original qubit and the disentangler itself. By tracing over the disentangler we then expect to obtain the best possible disentangled qubit in the state $`\rho ^{(out)}(\overline{\vartheta },\overline{\phi })`$. Now we impose several constraints which would specify what we mean by the optimal covariant (universal) disentangler: (1) The fidelity between the output of the disentangler and the original state $`\|\psi (\vartheta ,\phi )`$ has to be invariant with respect to rotations of the original qubit, so the fidelity has to be input-state independent. This universality of the disentangler would then guarantee that the information from the symmetric state (3) is extracted for all states equally well. (2) We are looking for the optimal disentangler which would disentangle the information with the highest fidelity. Imposing these two conditions we have found the unitary transformation which realizes the optimal covariant disentangler, i.e. which disentangle the qubit-state $`\|\psi `$ from the $`N`$-qubit state $`\|\mathrm{\Psi }`$ in the optimal and the $`\|\psi `$-state independent way (see Appendix). This disentangler is described by the transformation: $`\|N;0\|d_0`$ $``$ $`\|N1;0\left[\gamma _N\|0\|d_1+\delta _N\|1\|d_2\right];`$ (29) $`\|N;1\|d_0`$ $``$ $`\|N1;0\left[\delta _N\|0\|d_3+\gamma _N\|1\|d_1\right];`$ (30) where $`\|d_j`$ are three orthonormal basis vectors of the disentangler. The amplitudes $`\gamma _N`$ and $`\delta _N`$ given by the relation $`\gamma _N=\left({\displaystyle \frac{N+1}{2(N+1\sqrt{N})}}\right)^{1/2};\delta _N=\sqrt{1\gamma _N^2}.`$ (31) We can directly verify, that the fidelity $`_2=\psi (\vartheta ,\phi )\|\rho _d^{(out)}(\overline{\vartheta },\overline{\phi })\|\psi (\vartheta ,\phi )`$ is input-state independent and equal to $`_2=\gamma _N^2`$. Moreover, it can be shown that the transformation (31) is optimal, i.e. among all unitary transformations satisfying the given conditions the transformation (31) has the largest fidelity. We see that for $`N=1`$ the fidelity $`_2=1`$, which is obvious, because the original qubit has not been entangled with ancilla qubits. We plot $`_2`$ in Fig. 1. We see, that it is always larger than the fidelity of the disentanglement via measurement. In the limit $`N\mathrm{}`$ even the quantum disentangler gives us a totally random outcome. So in this limit, even optimal quantum entangler on which we impose the universality condition, is not able to extract information from the state (3). This is one of the main results of our paper - the optimal covariant quantum disentangler operates better than if the information is extracted (disentangled, distilled) from the symmetrized state (3) with the help the of optimal measurement. This is due to the fact that $`\overline{}^{max}_2`$. One can also ask the opposite question, how can we generate out of a qubit in an unknown state $`\|\psi `$ the symmetric state of the form (3). It can be shown that within quantum mechanics perfect universal entanglers, which would realize the inverse of the relation (7) do not exist. If one wants to create a state (3) from a qubit in an unknown state and $`N1`$ ancilla qubits in the known state $`\|0`$ again two scenarios are possible, the measurement-based and quantum scenarios. It is not surprising that the quantum scenario works better. We have found the optimal universal (covariant with respect to rotations of the input qubit) quantum entangler given by the transformations: $`\|0\|N1;0\|e_0`$ $``$ $`\left[\gamma _N\|N;0\|e_1+\delta _N\|N;1\|e_2\right];`$ (32) $`\|1\|N1;0\|e_0`$ $``$ $`\left[\delta _N\|N;0\|e_3+\gamma _N\|N;1\|e_1\right];`$ (33) where $`\|e_k`$ are three orthonormal basis states of the quantum entangler, $`\|e_0`$ is its initial state and the parameters $`\gamma _n`$ and $`\delta _N`$ are given by Eq.(31). One can check that the fidelity between the output of this entangler described by the density operator $`\rho _e^{(out)}(\vartheta ,\phi )`$ and the ideally entangled state (3) is input-state independent (i.e. does not depend on the parameters $`\vartheta ,\phi `$) and is equal to $`\gamma _N^2`$. This is the best possible universal (covariant) entangler. ### B State-dependent devices The universal disentangler gives a higher fidelity than does the best measurement-based procedure, but it is not obvious that this is the best that one can do. In the case of quantum cloning, the universal cloners are the ones which maximize the average fidelity . As we shall see, however, in the case of disentanglers this is no longer the case; there are state-dependent devices which are better. Consider the general disentangler transformation $`\|N;0\|b`$ $``$ $`\|N1;0(\|0\|D_1+\|1\|D_2)`$ (34) $`\|N;1\|b`$ $``$ $`\|N1;0(\|0\|D_3+\|1\|D_4),`$ (35) where the vectors $`\|b_{jk}`$, are states of the disentangler itself and need not be orthogonal. They must, however, satisfy the constraints imposed by the unitarity of the above transformation. The input state for the device is assumed to be $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })`$, and the ideal output state, to which the actual output should be compared, is $`\|\mathrm{\Psi }_{ideal}=\|N1;0\|\psi (\vartheta ,\phi )`$. The output state is calculated by starting with the input state, using the above transformation, and then tracing over the disentangler to obtain an output density matrix, $`\rho ^{(out)}`$. One then finds the average fidelity for this process, which we shall call $`\overline{}_3`$, from $$\overline{}_3=𝑑\mathrm{\Omega }\mathrm{\Psi }_{ideal}\|\rho ^{(out)}\|\mathrm{\Psi }_{ideal},$$ (36) Note that we are assuming a specific ensemble of input states; the probability of the one-qubit state $`\|\psi (\vartheta ,\phi )`$ is assumed to be constant on the Bloch sphere. Our result for the average fidelity for a state-dependent device depends on our choice of input ensemble, while for a state-independent device the average fidelity is independent of this ensemble. The calculation of the average fidelity is given in the Appendix, and will not be given in detail. We find that $`D_2^2=D_3^2=0`$ and $`\|D_1=\|D_4`$. This implies that the final state is just a product of the state of the $`N`$ particles and the entangler state, which means that the entangler states can be dropped from the problem. Therefore, the transformation which maximizes the average fidelity is just $`\|N;0`$ $``$ $`\|N1;0\|0`$ (37) $`\|N;1`$ $``$ $`\|N1;0\|1,`$ (38) and we have that $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })\|0\|N,0\|\psi (\overline{\vartheta },\overline{\phi }),`$ (39) which is a kind of state swapping transformation. The average fidelity itself is given by $`\overline{}_3=f_N`$, where the coefficient $`f_N`$ is given by Eq.(15). This average fidelity is larger than the fidelity of the optimal universal disentangler (see Fig. 1). In this case, the fact that the universality condition forces us to use an additional quantum device, the disentangler, with which the qubit at the output becomes partially entangled,results in a net loss of information. As a result the fidelity of the universal (covariant) entangler is smaller. Analogously, we find that quantum state-dependent entanglement can also be performed by a kind of state swapping transformation, i.e. $`\|\psi (\vartheta ,\phi )\|N;0\|0\|\mathrm{\Psi }(\vartheta ,\phi ).`$ (40) with input-state dependent fidelity $`\|\mathrm{\Psi }(\vartheta ,\phi )\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })\|^2`$. Nevertheless, when averaged over all values of $`\vartheta ,\phi `$ we find the mean fidelity of this state-dependent entangler to be equal to $`f_N`$ which on average is larger than the fidelity of the state-independent entangler. ## IV Probabilistic disentangler Let us examine a simple quantum network which takes as an input the $`N`$-qubit state (3). The network is composed of a sequence of $`N1`$ C-NOT gates $`𝒫_N=\mathrm{\Pi }_{k=1}^{N1}C_{kN}`$ where $`C_{kl}`$ is the C-NOT with $`k`$ being the control bit and $`l`$ being the target bit. This sequence of the C-NOT gates acts on the two vectors $`\|N;0`$ and $`\|N;1`$ as $`𝒫_N\|N;0`$ $``$ $`\|N1;0\|0`$ (41) $`𝒫_N\|N;1`$ $``$ $`{\displaystyle \frac{1}{\sqrt{N}}}\left(\sqrt{N1}\|N1;1+\|N1;0\right)\|1`$ (42) from which it follows that the input vector (3) is transformed as $`\|\mathrm{\Psi }(\overline{\vartheta },\overline{\phi })`$ $``$ $`{\displaystyle \frac{\sqrt{N}}{𝒩}}(\|v_+\|\psi (\vartheta ,\phi )`$ (44) $`+\sqrt{N1}\mathrm{cos}{\displaystyle \frac{\vartheta }{2}}\|v_{}\|0)`$ where $`𝒩=\sqrt{N^2\mathrm{cos}^2\frac{\vartheta }{2}+N\mathrm{sin}^2\frac{\vartheta }{2}}`$ is the normalization constant. In Eq.(44) we have introduced two orthogonal vectors of $`N1`$ qubits $`\|v_\pm `$. $`\|v_+`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N}}}\left\{\sqrt{N1}\|N1,1+\|N1,0\right\}`$ (45) $`\|v_{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N}}}\left\{\sqrt{N1}\|N1,0\|N1,1\right\}`$ (46) At the output of the network a projective measurement on the first $`N1`$ qubits is performed in order to determine whether they are in the state $`\|v_+`$ or $`\|v_{}`$. If the result $`\|v_+`$ is obtained, then the $`N`$th qubit is in the desired state $`\|\psi (\vartheta ,\phi )`$. The probability of this outcome is given by $`P_{\|v_+}={\displaystyle \frac{1}{N\mathrm{cos}^2\frac{\vartheta }{2}+\mathrm{sin}^2\frac{\vartheta }{2}}}.`$ (47) This probability is input-state-dependent, and it decreases with $`N`$. There is a difference between this probabilistic process and those considered previously, such as probabilistic cloning . Those only work for set of input states which is finite. The process considered above, however, works for a continuous, and hence infinite, set of input states. It, in fact, works for all input states of the type we are considering. Therefore, we can conclude that the range of applicability of probabilistic devices depends on the process being considered. ## V Conclusion We have considered a number of different methods of extracting an unknown state from an entangled state formed from that state and a known state. Measuring the state is, as expected, the least effective method. In the case of quantum devices, the universal device was not best one, at least if average fidelity is used as the criterion. Probabilistic quantum devices were seen to work very well for this operation in that they can be used for the entire set of input states. ###### Acknowledgements. This work was supported by the National Science Foundation under grant PHY-9970507, by the IST project EQUIP under the contract IST-1999-11053 and by the CREST, Research Team for Interacting Career Electronics. ## Appendix: Proof of optimality Let us consider the optimal quantum disentangler which acts as close as possible to the ideal transformation (7). The disentangler maps the space spanned by the vectors $`\|N;0`$ and $`\|N;1`$, into the space spanned by $`\|N1;0\|1`$ and $`\|N1,0\|1`$. This suggests that we consider a transformation of the following form $`\|N;0\|d_0`$ $``$ $`\|N1;0\left(\|0\|D_1+\|1\|D_2\right),`$ (48) $`\|N;1\|d_0`$ $``$ $`\|N1;0\left(\|0\|D_3+\|1\|D_4\right),`$ () where $`\|d_0`$ is the initial state of the disentangler which is supposed to be the same for all inputs and $`\|D_j`$ ($`j=1,\mathrm{},4`$) are some unnormalized disentangler state-vectors. Our task is to determine these vectors. Unitarity immediately implies that $`D_1^2+D_2^2`$ $`=`$ $`1`$ (49) $`D_3^2+D_4^2`$ $`=`$ $`1`$ () $`D_1\|D_3+D_2\|D_3`$ $`=`$ $`0.`$ (50) We shall now use our disentangler transformations (A.1) to calculate the fidelity of the actual output to the ideal output (7) The input of the disentangler is given by Eq. (3). If we introduce a notation $`\overline{\alpha }=\mathrm{cos}\frac{\overline{\vartheta }}{2}`$ and $`\overline{\beta }=\mathrm{e}^{i\phi }\mathrm{sin}\frac{\overline{\vartheta }}{2}`$ we can write the result of the transformation (A.1) $`\|\mathrm{\Psi }_{out}`$ $`=`$ $`\|N1;0[\overline{\alpha }(\|0\|D_1+\|1\|D_2)`$ () $`+\overline{\beta }(\|0\|D_3+\|1\|D_4)].`$ We now use this expression to find the output density matrix and trace out the disentangler itself. We define the $`N`$-qubit output density matrix to be $$\rho _{out}=\mathrm{Tr}_{\mathrm{disentangler}}(\|\mathrm{\Psi }_{out}\mathrm{\Psi }_{out}\|).$$ (52) The output fidelity is given by $$=\mathrm{\Psi }_{ideal}\|\rho _{out}\|\mathrm{\Psi }_{ideal},$$ (53) where $`\|\mathrm{\Psi }_{ideal}`$ is given by Eq. (7). If we denote $`\alpha =\mathrm{cos}\frac{\vartheta }{2}`$ and $`\beta =\mathrm{e}^{i\phi }\mathrm{sin}\frac{\vartheta }{2}`$ we can express this fidelity as $``$ $`=`$ $`{\displaystyle \frac{1}{(N\|\alpha \|^2+\|\beta \|^2)}}\{N\|\alpha \|^4D_1^2+\|\beta \|^4D_4^2`$ (54) $`+`$ $`\|\alpha \|^2\|\beta \|^2[D_3^2+ND_2^2+\sqrt{N}(D_4\|D_1+D_1\|D_4)]`$ (55) $`+`$ $`\alpha ^{}\beta \|\alpha \|^2(\sqrt{N}D_1\|D_3+ND_2\|D_1)`$ (56) $`+`$ $`\alpha \beta ^{}\|\alpha \|^2(\sqrt{N}D_3\|D_1+ND_1\|D_2)`$ (57) $`+`$ $`\alpha ^{}\beta \|\beta \|^2(\sqrt{N}D_2\|D_4+D_4\|D_3)`$ (58) $`+`$ $`\alpha \beta ^{}\|\beta \|^2(\sqrt{N}D_4\|D_2+D_3\|D_4)`$ (59) $`+`$ $`(\alpha ^{})^2\beta ^2\sqrt{N}D_2\|D_3+\alpha ^2(\beta ^{})^2\sqrt{N}D_3\|D_2\}.`$ () From this point on we will study two separate cases. Firstly, we will prove optimality of the universal disentangler and then the optimality of the state-dependent disentangler. ### A.1 Universal disentangler Demanding that the fidelity be independent of phases of $`\alpha `$ and $`\beta `$ we find that $`\sqrt{N}D_1\|D_3+ND_2\|D_1`$ $`=`$ $`0`$ (60) $`D_3\|D_2`$ $`=`$ $`0`$ () $`\sqrt{N}D_2\|D_4+ND_4\|D_3`$ $`=`$ $`0.`$ (61) Assuming these conditions to be satisfied the fidelity becomes $``$ $`=`$ $`{\displaystyle \frac{1}{(N\|\alpha \|^2+\|\beta \|^2)}}\{N\|\alpha \|^4D_1^2+\|\beta \|^4D_4^2`$ () $`+`$ $`\|\alpha \|^2\|\beta \|^2[D_3^2+ND_2^2`$ (62) $`+`$ $`\sqrt{N}(D_4\|D_1+D_1\|D_4)]\}.`$ (63) In order for this to be independent of $`\alpha `$ and $`\beta `$, the term in brackets must be proportional to $$(N\|\alpha \|^2+\|\beta \|^2)=N\|\alpha \|^4+(N+1)\|\alpha \|^2\|\beta \|^2+\|\beta \|^4.$$ (64) Comparing Eqs. (A.7) and (A.8) we find that $`D_1`$ $`=`$ $`D_4`$ (64) $`(N+1)D_4^2`$ $`=`$ $`D_3^2+ND_2^2`$ (65) $`+\sqrt{N}(D_4\|D_1+D_1\|D_4).`$ Combining these requirements with those imposed by unitarity we conclude that $$D_3^2=D_2^2=1D_4^2,$$ (66) and $`=D_4^2`$. This means that in order to maximize $``$, we must maximize $`D_4^2`$. Our first step in accomplishing this is to note that by combining the results of Eqs. (A.9) and (A.10) we have that $$(N+1)+2\sqrt{N}xD_4^2=2(N+1)D_4^2,$$ (66) where $$x=\frac{D_4\|D_1+D_1\|D_4}{2D_4^2},$$ (67) and $`1x1`$. Solving for $`D_4^2`$ we find that $$D_4^2=\frac{N+1}{2(N\sqrt{2}x)},$$ (68) which, assuming $`N2`$, is greatest when $`x=1`$. This implies that $`\|D_1=\|D_4`$ and that $`D_4^2`$ $`=`$ $`{\displaystyle \frac{N+1}{2(N+1\sqrt{N})}}`$ (69) $`D_3^2`$ $`=`$ $`D_2^2={\displaystyle \frac{N+12\sqrt{N}}{2(N+1\sqrt{N})}}.`$ () Imposing now the conditions on inner products we find that $$D_3\|D_4=D_2\|D_4=0.$$ (70) We can summarize our results in the following way. Let $`\{d_j\|j=1,2,3\}`$ be a set of three orthonormal vectors and define two parameters $`\gamma _N`$ and $`\delta _N`$ given by Eq. (31) we then have that $`\|D_4`$ $`=`$ $`\|D_1=\gamma _N\|d_1`$ (71) $`\|D_2`$ $`=`$ $`\delta _N\|d_2`$ (72) $`\|D_3`$ $`=`$ $`\delta _N\|d_3,`$ () and the universal optimal disentangler transformation is given explicitly by Eq. (30). ### A.2 Input-state dependent disentanglers In order to find the optimal input-state dependent disentangler we find the explicit form of the transformation (A.1) such that the averaged fidelity $`\overline{}=𝑑\mathrm{\Omega }`$ (with $``$ given by Eq. (A.5)) is maximized. Here, as usually, the integration measure is $`d\mathrm{\Omega }=\mathrm{sin}\vartheta d\vartheta d\phi /4\pi `$. Therefore after the integral over the phase $`\phi `$ is performed we can write the average fidelity as $`\overline{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{\xi _1ND_1^2+\xi _2D_4^2`$ () $`+`$ $`\xi _3[D_3^2+ND_2^2+\sqrt{N}(D_1\|D_4+D_4\|D_1)]\}`$ (73) with $`\xi _1`$ $`=`$ $`{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{sin}\vartheta d\vartheta }{N\mathrm{cos}^2\frac{\vartheta }{2}+\mathrm{sin}^2\frac{\vartheta }{2}}}\mathrm{cos}^4{\displaystyle \frac{\vartheta }{2}}`$ (74) $`\xi _2`$ $`=`$ $`{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{sin}\vartheta d\vartheta }{N\mathrm{cos}^2\frac{\vartheta }{2}+\mathrm{sin}^2\frac{\vartheta }{2}}}\mathrm{sin}^4{\displaystyle \frac{\vartheta }{2}}`$ () $`\xi _3`$ $`=`$ $`{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{sin}\vartheta d\vartheta }{N\mathrm{cos}^2\frac{\vartheta }{2}+\mathrm{sin}^2\frac{\vartheta }{2}}}\mathrm{sin}^2{\displaystyle \frac{\vartheta }{2}}\mathrm{cos}^2{\displaystyle \frac{\vartheta }{2}}`$ (75) After the integration over the parameter $`\vartheta `$ we find $`\xi _1`$ $`=`$ $`{\displaystyle \frac{34N+N^2+2\mathrm{ln}N}{(N1)^3}}`$ (76) $`\xi _2`$ $`=`$ $`{\displaystyle \frac{1+4N3N^2+2N^2\mathrm{ln}N}{(N1)^3}}`$ () $`\xi _3`$ $`=`$ $`{\displaystyle \frac{1+N^22N\mathrm{ln}N}{(N1)^3}}`$ (77) From the unitarity of the disentangling transformation it follows that $`D_2^2=1D_1^2`$ and $`D_3^2=1D_4^2`$. When we introduce the notation $$u=\frac{D_4\|D_1+D_1\|D_4}{2D_1D_4},$$ (78) where $`1u1`$, and $`\eta _1=D_1^2`$; $`\eta _4=D_4^2`$ we can rewrite the average fidelity (A.17) as $``$ $`=`$ $`{\displaystyle \frac{1}{2}}[\eta _1N(\xi _1\xi _3)+\eta _4(\xi _2\xi _3)`$ (78) $`+`$ $`2\sqrt{N}\xi _3u\sqrt{\eta _1\eta _4}+\xi _3(1+N)].`$ () Taking into account that $`\xi _1>\xi _3`$ and $`\xi _2>\xi _3`$ we easily find that the maximum of the mean fidelity (A.21) is achieved for $`u=1`$ and $`\eta _1=\eta _4=1`$. In this case we rewrite (A.21) as $`={\displaystyle \frac{1}{2}}[\xi _1N+\xi _2+2\sqrt{N}\xi _3].`$ () When we substitute into Eq. (A.22) the explicit expression for the parameters $`\xi _j`$ given by Eq. (A.19) we find that the mean fidelity is equal to the function $`f_N`$ given by Eq. (15). This exactly is equal to the mean fidelity of the input-state disentanglement performed via the state swapping transformation described by Eq. (38). In fact, from our conditions $`\eta _1=\eta _4=1`$ it directly follows that $`D_2^2=D_3^2=0`$ while $`D_1^2=D_4^2=1`$. In addition, from $`u=1`$ it follows that $`\|D_1=\|D_4`$, so that the optimal state-dependent disentangling transformation is indeed equal to Eq. (38), which we wanted to prove.
warning/0006/astro-ph0006225.html	ar5iv	text	# 150 keV Emission from PKS2149-306 with BeppoSAX ## 1 Introduction High redshift quasars are of interest not only for their ‘record setting’ quality, but also because they can tell us about the formation of quasars and about conditions in the universe at early times (Turner 1991, Hamann & Ferland 1994). We should also be able to learn about the quasar emission processes by comparing the continua of quasars at low z with those of their cousins at the epoch at which the characteristic quasar luminosity peaks around z$``$2 (Boyle et al. 1987). Since small changes in physical conditions can greatly alter the tails of particle energy distributions the extreme high energy continuum may be a particularly good place to look for evolutionary effects. So far though X-ray spectra of high redshift quasars from ROSAT (Elvis et al., 1994a, Fiore et al., 1998a) and ASCA (Elvis et al., 1994b, Cappi et al., 1997) have been few, and of limited signal-to-noise. Even so they have begun to show redshift dependent features that challenge models and require further XJB investigation: low energy cut-offs were discovered with ROSAT (Elvis et al., 1994a), and appear to be intrinsic (Fiore et al., 1998a, Elvis et al., 1998) and related to ultraviolet absorbing outflows; no Compton reflection component has been seen in z$`>`$1 quasars ($`\mathrm{\Omega }_d/2\pi 0.4`$ \[90%\], Elvis et al. 1994b, Nandra et al., 1995, Siebert et al., 1995, see also Williams et al (1992) for low z examples); nor is a 6.4 keV (rest frame) iron line required by the fits (EW$`<`$20 eV \[90%\] Siebert et al. 1995; EW$`<`$120 eV (90%) Elvis et al., 1994b). These properties differentiate high z quasars from typical low redshift Seyfert galaxies which have strong iron K-lines (EW=100-300 eV), and a strong Compton hump that requires half the sky (as seen from the X-ray source) to be covered with Compton-thick ‘cold’ material ($`\mathrm{\Omega }_d/2\pi `$=1, Nandra & Pounds 1994 and references therein). This suggests a different structure within high luminosity quasars. We observed PKS2149$``$306 as part of an AO1 GO program on BeppoSAX in order to examine these issues. At z=2 the Compton hump region is well within the BeppoSAX (Boella et al., 1997a) 1.3–10 keV Medium Energy Concentrator Spectrometer (MECS, Boella et al., 1997b) band, making a measurement feasible. Moreover the good sensitivity of the high energy Phoswich Detector System (PDS, FWHM=1.4 degrees, Frontera et al., 1997) on BeppoSAX allows the extension of the spectrum to energies (15–100 keV) where a Compton hump, bremsstrahlung and a power law can be more readily distinguished PKS2149-306 is a z=2.34 flat spectrum ($`\alpha _{2.7GHz}^{4.85GHz}=0.0\pm 0.14`$, where $`f_\nu \nu ^\alpha `$, NED), core-dominated <sup>1</sup><sup>1</sup>1 The 2.3 Ghz data of Morabito et al (1986) shows the source to be $`>70`$ % compact on a 3mas scale (24 pc at z=2.34, $`H_0`$=50 km s<sup>-1</sup>Mpc<sup>-1</sup>, $`\mathrm{\Omega }`$=1), using the Parkes 2.7 Ghz total flux (Quiniento & Cersosimo 1993) as a comparison. radio-loud quasar with $`m_V`$=17.9. PKS 2149–306 is extraordinarily bright in the ROSAT Sky Survey (f<sub>X</sub>=10<sup>-11</sup> $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ , Siebert et al., 1995, Brinkmann et al., 1995b, Schartel et al., 1996). Being bright, PKS2149$``$306 also promised to give the best possible limits on high energy emission from quasars. Instead, it gave us a clear detection. We present here the BeppoSAX 0.1-130 keV data for PKS2149–306, along with the broad band radio to $`\gamma `$-ray spectral energy distribution (SED) of the quasar, and we compare these with recent Blazar models. ## 2 Observations and data analysis BeppoSAX observed PKS2149$``$306 from 1997 October 31 through 1997 November 1. The two imaging detectors - the 1.3-10 keV MECS and the 0.1-10 keV Low Energy Concentrator System (LECS, Parmar et al., 1997) - detected the quasar with high signal-to-noise at the center of the their field-of-view, at 0.105 ct/s and 0.05 ct/s respectively. (With exposures of 39,439 s and 17,862 s respectively.) Unexpectedly, the co-aligned Phoswich Detector System (PDS, Frontera et al., 1997) also detected a clear signal in its 15-100 keV bandpass of 0.26$`\pm 0.05`$ ct/s (in an effective on-source exposure time of 16,676 s). Standard data reduction was performed using the SAXDAS software package (1997). In particular, data are linearized and cleaned from Earth occultation periods and unwanted periods of high particle background (due to satellite passages through the South Atlantic Anomaly, SAA). Data from MECS units 2 and 3 were combined after gain equalization (unit 1 has been inoperative since 1997 May 7). The MECS and LECS pulse height spectra of PKS2149$``$306 were extracted using 3 arcmin and 8 arcmin radius extraction regions, respectively. These radii maximize the signal-to-noise ratio below 1 keV in the LECS and above 2 keV in the MECS. Background was taken from high Galactic latitude blank fields using the same extraction regions as for the source regions. Although background is negligible for this source, we checked that the level of the background in regions surrounding the source is within 10% of that recorded in the same regions in the blank fields. For the PDS the problem of accurate background subtraction is crucial since the PDS is a non-imaging detector. The PDS uses rocking collimators, one for each of two pairs of detectors, that switch positions quickly every 92 s so that one of the pairs is always pointed at the source, while the other is monitoring the background. The spectra from the four PDS crystals were summed together. The net source spectrum was then obtained by subtracting the ‘off-source’ background from the ‘on-source’ counts, scaling by the exposure time. The background in the equatorial BeppoSAX orbit varies by only $``$10-20%, thanks to shielding by the Earth’s particle belts (this was the main reason for selecting an equatorial orbit for BeppoSAX), and has an active anti-coincidence CsI(Na) particle shield (which is also used to detect gamma-ray bursts, Frontera et al., 1997). The main background variations are recorded just after the passages through the SAA. (The instrument is switched off while in the SAA). The first five minutes after SAA passages are excluded from the data analysis to allow an automatic gain control system to return the high voltages to their normal working value. This system keeps the pulse-height to energy scale constant to 1%. Using the rocking technique it is possible to take into account accurately the residual background variations. The magnitude of residual systematic errors in PDS spectra have been evaluated by Guainazzi and Matteuzzi (1997) at $`\genfrac{}{}{0pt}{}{_<}{^{}}0.05`$ counts s<sup>-1</sup>, using deep observations of blank fields. Even if this maximum residual is subracted from the observed count rate, the signal to noise in this observation would remain $`>4\sigma `$. The source is detected at $`>3\sigma `$ up to 40-50 keV, which is about 150 keV in the quasar rest frame. The chance of finding a source in any given 2 sq.degree, the FWHM beam area of the PDS, is small. The HEAO-1 A4 all sky catalog (Levine et al., 1984) lists just 7 high Galactic latitude ($`\|b\|>20^{}`$) sources in the 13-80 keV band, which is closely comparable to the PDS band, down to a flux of 2$`\times `$10<sup>-10</sup>erg cm<sup>-2</sup> s<sup>-1</sup> (10mCrab). The PKS2149$``$306 signal is 7.5 times fainter so, assuming a logN-logS slope of $``$1.5, we expect a chance coincidence rate of 1.4%, hence we believe the PDS signal does come from PKS2149–306. Our confidence in the identification of the PDS signal with PKS2149$``$306 is increased further by the agreement to better than 10% of the normalizations of the PDS and MECS spectra and the similarity of their slopes (see next section). The EGRET source reported later is too far away to produce the PDS signal (see §4). ## 3 X-ray Spectral Fitting Spectral fits were performed using the XSPEC 10.0 software package and versions of the response matrices made public on 1997 August 31. PI channels in MECS and LECS spectra are rebinned sampling the instrument resolution with the same number of channels at all energies and ensuring at least 20 counts per bin for the other. This guarantees the applicability of the $`\chi ^2`$ method in determining the best fit parameters, since the distribution in each channel can be considered Gaussian. Constant factors have been introduced in the model fitting in order to take into account the intercalibration systematics between instruments (BeppoSAX Cookbook, Fiore et al. 1998b). If we use the two MECS as reference instruments, then the LECS and PDS normalizations lie in the ranges 0.7-1.0 and 0.80-0.95 relative to the MECS, respectively. The (observed frame) energy ranges used for the fits are: 0.1-4 keV for the LECS, 1.65-10.5 keV for the MECS, and 14-130 keV for the PDS. All the results given in table 1 are in the observed frame. Quasar frame values are given in the text where appropriate. Errors in Table 1 are 90 % confidence intervals for 1 interesting parameter. A simple power-law plus absorption model gives a reasonably good fit (Figure 1, Table 1), but with systematic deviations at low energies. PKS2149–306 appears to have faded slightly over the three years between the obsevations: the 2-10 keV BeppoSAX flux, using the simple power law fit, is ($`8.0\pm 0.2)\times 10^{12}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ , 80% of the flux recorded in the 1994 October ASCA observation ($`9.9\times 10^{12}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ , Cappi et al 1997); and the monochromatic 2 keV BeppoSAX flux is ($`2.8\pm 0.1)\times 10^{12}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ keV<sup>-1</sup>, 76% of the 1990/1991 ROSAT RASS detection at ($`3.7\pm 0.3)\times 10^{12}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ keV<sup>-1</sup> (Schartel et al. 1996). The LECS shows a clear excess of counts below about 1 keV. This soft excess seems in contradiction with the above Galactic N<sub>H</sub> reported by Yaqoob et al. (1999). While variability in the 11 months (quasar frame) between the two observations is possible, we also note that the low energy calibration of ASCA is uncertain <sup>2</sup><sup>2</sup>2see: http://heasarc.gsfc.nasa.gov/docs/asca/watchout.html http://heasarc.gsfc.nasa.gov/docs/asca/cal\_probs.html , and we have more confidence in the Beppo-SAX value (Parmar et al., 1997, Mineo et al., 2000). The presence of a low energy excess means that the X-ray spectrum cannot give a reliable Galactic absorption value, so we fix the absorbing column to the Galactic 21 cm value (2.2$`\times `$10<sup>20</sup>cm<sup>-2</sup>, Dickey & Lockman 1990). The resulting energy spectral index, $`\alpha =0.40\pm 0.04`$, is extremely flat (c.f. Sakano et al. 1998, Sambruna et al., 1999), although consistent with other high redshift radio-loud quasars measured with ASCA from 2-10 keV , including PKS2149$``$306 itself ($`\alpha `$=0.42$`{}_{0.02}{}^{}{}_{}{}^{+0.03}`$ Cappi et al. 1997, 0.54$`\pm `$0.05 Yaqoob et al. 1999). This slope is close to that of the X-ray background (Marshall et al., 1980) although this class of quasar is too rare to integrate to the X-ray background (see Elvis et al., 1994 and references. therein). A better fit can be had with a broken power-law to account for the slight upturn in the LECS spectrum below $``$1 keV (Table 1). In this case the steeper low energy slope is $`2.0_{1.0}^{+2.0}`$ allowing the absorbing column to be free to vary, or $`1.0_{0.3}^{+0.6}`$ fixing $`N_H`$ to its Galactic value. Both slopes are consistent with the ROSAT RASS value of $`1.24\pm 0.80`$ (Schartel et al., 1996). The break energy is about 0.8 keV (observed frame, 2.7 keV quasar frame) in both cases. The high energy index is indistinguishable from the single power law index. An apparent soft excess can also be generated though an absorption edge. In the quasar frame there is no sign of an ionized oxygen absorber, although at $`z=2.34`$ the energy of the deepest edges usually seen in low luminosity Seyfert galaxies (OVII and OVIII) fall at 0.22 and 0.26 keV (observer frame) respectively, close to the carbon edge (0.283 keV) produced by the LECS window, and are therefore affected by large systematic errors. A fit with an absorption edge and a single power law (Table 1) gives a good $`\chi ^2`$. The edge energy is 2.8 keV in the quasar frame, which has no ready interpretation, but the observed energy of 0.8 keV is consistent with either OVIII at $`z<0.2`$ or OVII at $`z=0`$. The optical depth is large, $`\tau 1`$, implying a total hydrogen column of a few $`10^{22}`$ cm<sup>-2</sup> assuming cosmic metal abundances. This is about one hundred times what it is seen from the Galactic 21 cm line. The high energy slope in this fit is consistent with the value from the broken power law fit. No line is detected around 2 keV, the redshifted iron K$`\alpha `$ energy. The 90 % upper limits to the emission line equivalent width are 50 eV and 25 eV (observer’s frame; 167 eV and 84 eV quasar frame) for the 6.4 and 6.7 keV (quasar frame energies) iron K$`\alpha `$ lines respectively. No line is seen at 5.1 keV (observer’s frame, 17 keV quasar frame) with a limit of 63 eV (90% confidence, observer’s frame), or 210 eV (quasar frame). This compares with the 298$`{}_{205}{}^{}{}_{}{}^{+202}`$ eV (quasar frame, 90% errors for two interesting parameters) measurement reported at thsi unusual energy by Yaqoob et al. (1999) from 1994 October ASCA data. The presence of a steeper low energy component or absorption edge has little effect on the high energy slope. Since this steeper component is only important in the LECS data we shall fit only the MECS and PDS data when investigating the high energy component, as we do in the rest of this section. This has the advantage of removing extra free parameters from the models, which usefully decreases the allowed parameter space. To determine the overall 5-150 keV rest frame spectral shape we next fitted three XSPEC models: a power-law with an exponential cutoff (cutoffpl); a Compton reflection model for neutral reflectors (pexrav, Magdziarz & Zdziarski 1995); and a thermal bremsstrahlung model. In pexrav the redshift was frozen at the quasar redshift, the disk inclination was frozen at cos $`i`$=0.45, and the cosmic abundances of Anders & Ebihara (1982) were assumed. For both models the local absorption was fixed at the Galactic value. The Compton reflection model (Table 1) gave a 90% upper limit of 0.3 to the relative normalization, $`R`$, of the reflected component <sup>3</sup><sup>3</sup>3$`R=\mathrm{\Omega }/2\pi `$, for an isotropic source above a flat infinite disk . This compares with typical values of 1.0 in Seyfert galaxies (Matt 1998). The cut-off power-law model (Table 1) requires the cut-off energy, E<sub>c</sub>, to be $`>`$36 keV (90% confidence, observed frame), or $`>`$120 keV in the quasar frame. A bremsstrahlung fit gives a rest frame temperature of 46$`{}_{16}{}^{}{}_{}{}^{+32}`$ keV (90% confidence, 1 interesting parameter), similar to that found in several z=3 radio-loud quasars in the 2-10 keV band with ASCA (Cappi et al., 1998). The Beppo-SAX result suggests that these temperatures are not merely an artifact of the upper energy limit of ASCA. As for ASCA, a bremsstrahlung model is a slightly worse fit than power-law models. ## 4 The Spectral Energy Distribution of PKS2149$``$306 The addition of the broad band BeppoSAX spectrum from 0.1-50 keV (observed) allows the construction of a well-sampled Spectral Energy Distribution (SED) since PKS2149$``$306 is already XJB well-observed at radio wavelengths, as shown by a NED <sup>4</sup><sup>4</sup>4The NASA/IPAC Extragalactic Database (NED) is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. search. Additional radio (Quiniento et al., 1993) and millimeter data (Steppe et al., 1988) fill out the long wavelength SED. Although PKS2149–306 has a flat radio spectrum, Cersosimo et al (1994) include PKS2149-306 as a candidate GigaHertz Peaked Spectrum (GPS) radio source, because of a mild curvature in its radio spectrum. PKS2149–306 also has low radio polarization ($`<1.7`$% at 1.4 MHz, Condon et al 1998). These two characteristics of GPS sources make it unlikely that such sources are dominated by beamed emission (O’Dea 1998). PKS 2149–306 though, has a Blazar-like SED and is unlikely to be a normal GPS source. The curvature in the PKS2149–306 radio spectrum may instead be due to the Compton downshifting of the originating electron spectrum, caused by a quenching UV source (see §5.1) In the 100$`\mu `$m to 0.1$`\mu `$m range there are only IRAS Faint Source Survey (FSS, Moshir et al., 1989) upper limits and an optical spectrum (Wilkes et al., 1983, Wilkes 1997, private communication). These few data points though are importantly constraining. A search of the IRAS FSS noise plates gives 3 $`\sigma `$ upper limits of 130 mJy, 96 mJy, 70 mJy and 600 mJy at 12$`\mu `$m, 25$`\mu `$m, 60$`\mu `$m and 100$`\mu `$m respectively. These limits are stable against small changes in position, so that a complex background (e.g. from Galactic cirrus) is not causing spurious values. The Compton Gamma Ray Observatory EGRET summed phase 1-4 (1991–1995) exposure (9$`\times `$10<sup>8</sup> seconds) of the field shows no evidence for emission from PKS2149–306. A point source fit at the position of the quasar gives a 3 sigma upper limit of $`6.9\times 10^8`$ photons cm<sup>-2</sup>s<sup>-1</sup> (30 MeV-3 GeV). We do recover the well known BL Lac object PKS 2155–304 (3EG J2158–3023, Hartman et al 2000) with a flux of ($`1.2\pm 0.3)\times 10^7`$ photons cm<sup>-2</sup>s<sup>-1</sup> (30 MeV-3 GeV). However, this source lies 1.7 deg from PKS2149–306, well outside the 99 percent confidence region of the EGRET source position. The PDS reaches a transmission of 50 percent at 0.7 degrees and is effectively at zero transmission by 1.7 degrees, so PKS 2155–304 cannot be the origin of the PDS signal. The complete SED of PKS2149$``$306 (Figure 2) shows how unusually hard and energetically dominant the hard X-ray, PDS, spectrum is when compared with normal radio-loud quasar SEDs (dot-dash curve, Elvis et al., 1994c). The whole Beppo-SAX spectrum lies an order of magnitude above the normal level for radio-loud quasars. The X-ray ‘soft excess’ (E$`<`$3 keV in the quasar frame) appears to be a smooth continuation of the $`\alpha 1.0`$ optical slope (with the caveat that an ionized absorber could fake an excess, see §3). In contrast the 3-150 keV slope is harder by $`\mathrm{\Delta }\alpha 0.5`$. However, unless the source varies by a factor greater than 50 at $``$1 GeV, this hard X-ray MECS/PDS slope cannot extend up to EGRET energies, and must instead turn down somewhere between 120 keV (the quasar frame lower limit on a cut-off in the PDS spectrum) and 40 MeV. The SED of PKS 2149–306 (figure 2) seems to have two peaks at about $`10^{12}`$ and $`10^{21}`$ Hz. This ‘two peaks’ SED shape is similar to those of blazars (Fossati et al. 1998), strongly suggesting that the observed SED, outside the optical band, is dominated by a blazar continuum originating in a jet. The solid and dashed lines on the SED (figure 2) are spline interpolations of the highest frequency radio point and the X-ray data above 2 keV constrained not to exceed the IRAS and EGRET upper limits. The solid line is a spline fit close to one of the synchro-Compton models from Ghisellini et al.(1998). The lines are chosen to indicate the uncertainty in determining the ratio of the two peaks. (The lines are not fits to the data with a physical model, and are only useful to guide the eye.) The two spline curves define the ranges of the two peak frequencies and amplitudes. The $`\gamma `$-ray peak lies at a few MeV, $`\nu _{peak}(IC)=`$10<sup>21.0±1.0</sup>Hz, given the EGRET upper limit. The peak of the low frequency, radio-mm, peak is in the sub-millimeter range, $`\nu _{peak}(synch)=10^{12.0\pm 0.5}`$ Hz. These are both unusually low frequencies. The ratio of the two peak luminosities, the Compton dominance, $`C_D`$ (=$`logL_{IC}/L_{synch}`$, Ghisellini et al., 1998 )$`=1.4\pm 0.4`$, is near the top of the known range. ## 5 Discussion ### 5.1 PKS 2149–306: An Extreme Compton Dominant Quasar The detection of hard X-ray emission in PKS 2149$``$306 reaches an energy some 4 times higher than detected in any other z$`>`$1 quasar, except for the EGRET blazars. Moreover the X-ray to optical ratio of PKS 2149–306 is large, a factor 10 larger than normal for radio-loud quasars. Expressed as $`\alpha _{OX}`$, the optical (2500Å) to X-ray (2 keV) slope, 0.81, is unusually small. Only two radio-loud quasars out of some 250 in the survey by Brinkmann et al. (1995a) have $`\alpha _{OX}<0.9`$. Clearly there is something unusual about this quasar. The SED of PKS 2149$``$306 is not that of a typical, unbeamed, AGN. The ‘two peaks’ SED shape has not been seen in ‘normal’ unbeamed quasars (Elvis et al., 1994, Mattox 1994), although existing limits are surprisingly weak), but is common in relativistically beamed objects, both BL Lacs and quasars (Padovani & Giommi 1996). The absence of reflection components (Compton hump and iron line) strongly limits the contribution to the observed flux of any Seyfert-like component, although these features are also weak in radio-quiet high luminosity quasars (Iwasawa & Taniguchi 1993). The ‘two peaks’ SEDs of BL Lac objects have often been considered within the Synchrotron-Inverse Compton formalism (e.g. Sambruna, Maraschi & Urry 1996). In this scheme the low frequency (radio-mm) peak is the primary synchrotron peak, while the high frequency ($`\gamma `$-ray) peak comes from photons Inverse Compton scattered off the relativistic electrons. These scattered photons may be the synchrotron photons themselves (the self-Compton case), or from an external radiation source. The frequencies at which the two components peak in $`log\nu f_\nu vs.log\nu `$ space varies by 5 decades, which changes their observational appearance greatly. Fossati et al. (1998) and Ghisellini et al. (1998) have suggested that the Blazars form a single, smooth sequence with bolometric luminosity, from the lower luminosity High energy cut-off BL Lacs (HBL) to the Low energy cut-off BL Lacs (LBL) and then to the high luminosity flat spectrum radio quasars (FSRQs, Sambruna, Maraschi & Urry 1996, Padovani, Giommi & Fiore 1997). Unlike HBLs and LBLs, FSRQs show normal quasar emission in the optical/UV, i.e. both the broad emission lines and the ‘big bump’ blue continuum (with a presumed accretion disk origin) are seen. This optical emission can occasionally dominate, as it does in PKS2149–306. We conclude that PKS2149–306 is a high luminosity FSRQ. In PKS2149–306 both SED peaks lie at the extreme lower ends of the distributions in the Fossati et al (1998) and Ghisellini et al (1998) samples. In the Fossati et al. scheme the ‘Compton dominance’ (the ratio between the luminosity in the two peaks) is inversely proportional to the frequency of the synchrotron peak, roughly to the one half power (Figure 3). PKS2149–306 fits this correlation, in that the large Compton dominance of PKS2149–306 agrees with low frequency position of the synchrotron peak of its SED. Ghisellini et al. go on to suggest that a strong external radiation source, such as an accretion disk powered ‘optical/UV big bump’, provides a bath of cool photons that can efficiently drain energy from jet electrons by Compton scattering them to high energies. The low energy photons may come directly from the disk below (Dermer & Schlickeiser 1993) or from above or from the side (Siemiginowska & Elvis 1994, Sikora et al., 1996). This cooling makes inverse Compton losses dominate over synchrotron losses making the radio spectrum weak relative to the X-ray, so creating a ‘Compton dominant’ SED. PKS2149–306 seems to have a large ‘optical/UV big bump’, as expected in this model. Such large optical/UV big bumps are unusual, but not unprecedented in FSRQ (e.g. 3C273, Elvis et al., 1994c). An illustrative accretion disk spectrum, that roughly fits the big bump, is shown in figure 2 ($`M=10^{10}M_{}`$, $`\dot{m}/\dot{m}_{Edd}`$=0.1, from the tabulation given in Siemiginowska et al. 1995). ### 5.2 Low Energy Excess or Ionized Absorber? The soft X-ray excess in PKS 2149–306 appears unusually strong (Figure 2). We noted in §3 that an absorption edge imprinted onto a single power law gives an equally good fit. Since ionized absorbers of the kind that could produce such an edge are common in AGN and blazars we must consider this interpretation on an equal footing with a ‘soft excess’ component. An absorption edge can be more tightly constrained than a soft excess. A local z$``$0 absorber is unlikely since no molecular cloud is nearby (Marscher 1988, Turner et al., 1995), nor is any low z galaxy seen on the STScI digitization of the sky surveys. An absorber associated with the quasar is more likely, X-ray absorption with a similar column density is common in radio-loud high redshift quasars (Elvis et al., 1994a, Cappi et al., 1997, Fiore et al., 1998a, Elvis et al., 1998), and there is evidence for high ionization oxygen absorbers of similar column density in some low redshift blazars (Canizares & Kruper 1984, Koenigl et al. 1995, Sambruna et al. 1997, Grandi et al. 1997, Cagnoni & Fruscione 2000), though these absorbers are not always seen (Chiappetti et al. 1999). Chandra high resolution grating spectra however do show comparable features (Fruscione et al. 2000, in preparation). To produce an apparent $`z`$=0 absorber would require a coincidence of outflow velocity and cosmological redshift. A substantial ejection velocity, $`v=0.6c`$, is implied. In PKS 2155–304 the absorber has a suggested velocity of $`\mathrm{\Delta }v0.1c`$ relative to the core (Canizares & Kruper 1984) so a high ejection velocity in PKS 2149–306 is not inconceivable. $`v=0.6c`$ is larger than the $`v=0.10.2c`$ seen in broad absorption line quasars (Turnshek 1988); it is much slower, though, than the velocities needed to explain superluminal radio sources (Zenzus 1997), although these may be pattern speeds rather than true expansion velocities. Only X-ray spectra with significantly better resolution can definitively identify absorption edges. Chandra grating observations are scheduled. ## 6 Conclusions We have detected the z=2.34 quasar PKS2149$``$306 up to 150 keV in its rest frame, using the PDS on BeppoSAX. This is some 4 times higher in energy than any other z$`>`$1 quasar, except for the EGRET blazars. PKS 2149-306 is also an order of magnitude X-ray loud relative to the optical ($`\alpha _{OX}`$=0.81). The $`>`$3 keV spectrum of PKS 2149–306, as for several other high redshift radio-loud quasars, is extremely hard ($`\alpha =0.41\pm 0.05`$, kT=46$`{}_{16}{}^{}{}_{}{}^{32}`$ keV, observed frame), comparable to the X-ray background spectrum. The PKS2149$``$306 SED seems to show the two peaks typical of BL Lacertae objects, suggesting that the X-ray 1-100 keV observed emission is dominated by an inverse Compton component. Unusually the SED is strongly dominated by the $`\gamma `$-ray peak. A strong ‘big bump’ in the optical/ultraviolet seems to be present in PKS2149–306. Possibly, as suggested by Ghisellini et al. (1998), this bump is the source of the photons that cool the jet, suppressing the radio emission and creating the Compton dominated spectrum. This makes PKS2149$``$306 an extreme example of the luminosity based one-parameter unification scheme of Ghisellini et al.. The low energy X-ray spectrum of PKS 2149–306 shows no absorption above the Galactic value, contrary to the ASCA 1994 data (Yaqoob et al. 1999), showing instead a soft excess. An absorption edge fits the data as well as a broken power-law. One interpretation of an absorption edge is an oxygen absorber with N$`{}_{H}{}^{}`$10<sup>22</sup>cm<sup>-2</sup>, blueshifted from the quasar redshift and implying an ejection velocity, $`v0.6c`$. Such an absorber would provide a link with absorbers in other high redshift quasars, and with the absorbers in some low z blazars. However, the data are ambiguous and higher spectral resolution observations are needed to check this result. The Beppo-SAX data set strong limits on the X-ray reflection features: EW(Fe-K)$`<`$15 eV (quasar frame) and $`R<0.3`$ for the Compton Hump. Hard X-ray observations, as we hoped when initiating this project, do seem to select physically extreme objects. ###### Acknowledgements. We thank Belinda Wilkes for providing a digital version of her optical spectrum of the quasar, Daryl Macomb for providing the EGRET upper limit, Seth Digel for the IRAS data, Gabriele Ghisselini for the data points in figure 3, and the referee for a careful reading that improved the paper. This work was supported in part by NASA contract NAS8-39073 (ASC), and NASA grants NAGW-2201 (LTSA) and NAGW- (ADP). This research has made use of data obtained through the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA-Goddard Space Flight Center. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, Caltech, under contract with the National Aeronautics and Space Administration Based on photographic data obtained using The UK Schmidt Telescope. The UK Schmidt Telescope was operated by the Royal Observatory Edinburgh, with funding from the UK Science and Engineering Research Council, until 1988 June, and thereafter by the Anglo-Australian Observatory. Original plate material is copyright (c) the Royal Observatory Edinburgh and the Anglo-Australian Observatory. The plates were processed into the present compressed digital form with their permission. The Digitized Sky Survey was produced at the Space Telescope Science Institute under US Government grant NAG W-2166.
warning/0006/astro-ph0006432.html	ar5iv	text	# Untitled Document Lensing Degeneracies Revisited Prasenjit Saha Astronomy Unit, School of Mathematical Sciences Queen Mary and Westfield College London E1 4NS, UK This paper shows that the mass-sheet degeneracy and other degeneracies in lensing have simple geometrical interpretations: they are mostly rescalings of the arrival-time surface. Different degeneracies appear in Local Group lensing and in cosmological lensing, because in the former the absolute magnification is measured but the image structure is not resolved, whereas in the latter the reverse usually applies. The most dangerous of these is a combination we may call the ‘mass-disk degeneracy’ in multiply-imaging galaxy lenses, which may lead to large systematic uncertainties in estimates of cosmological parameters from these systems. Subject headings: gravitational lensing To appear in AJ, Oct 2000 1. Introduction A curious feature of gravitational lensing is that most of the observables are dimensionless. This fact leads to some scaleabilities in lensing theory, which show up as parameter degeneracies when interpreting observations. These degeneracies were analyzed in detail in Gorenstein et al. (1988, hereafter G88), elaborating on Falco et al. (1985). Since that time, while lensing theory has not changed much the observational situation has changed greatly—recall that in 1988 cluster lensing was still controversial, and Milky Way microlensing was several years in the future—and with it the emphasis of theory has shifted. So it is interesting at this time to rederive the G88 degeneracies, discuss their current observational context, and try to gain some new insights into the old results. G88 discussed three basic degeneracies, which they called the similarity, prismatic and magnification transformations, and combinations of them. The most subtle of these is the magnification transformation; it seems to have been independently discovered at least two more times, and is now usually called the mass-sheet degeneracy. The present paper will be about the same transformations, but unlike G88 who started from the lens equation, we will think about transformations of the arrival-time surface. 2. The degeneracies In fact, lensing degeneracies can all be interpretated as simple transformations of arrival-time surface $$t(\theta )=t_{\mathrm{geom}}+t_{\mathrm{grav}}.$$ $`(1)`$ and we can recognize three kinds. $``$‘Similarity transformations’ scale both $`t_{\mathrm{geom}}`$ and $`t_{\mathrm{grav}}`$ by a constant factor. Such operations scale time delays between images but leave image positions and magnifications unchanged. $``$The ‘mass-sheet degeneracy’ mixes $`t_{\mathrm{geom}}`$ and $`t_{\mathrm{grav}}`$ but in such a way that $`t(\theta )`$ is scaled by a constant factor. This multiplies time delays and all magnifications by a constant factor but leaves image positions and relative magnifications unchanged. $``$Various other transformations can be written down that modify $`t_{\mathrm{grav}}`$ and possibly also $`t_{\mathrm{geom}}`$ in various ways, but leave $`t(\theta )`$ and its derivatives unchanged at all image positions. These will have no effect on observables, but they cannot be ignored because they imply uncertainties in what can be inferred about lenses from the observables. To derive these degeneracies, we start with the full expression for the arrival time $$t(\theta )=\frac{1}{2}(1+z_\mathrm{L})\frac{D_\mathrm{L}D_\mathrm{S}}{cD_{\mathrm{LS}}}(\theta \beta )^2(1+z_\mathrm{L})\frac{8\pi G}{c^3}^2\mathrm{\Sigma }(\theta ),$$ $`(2)`$ where $`\theta `$ and $`\beta `$ are angular positions on the image and source planes respectively, $`^2`$ denotes the inverse of a two-dimensional Laplacian, and the other symbols have their usual meanings. 2.1. The similarity transformations The simplest of the degeneracies appears if the distance factor is Eq. (2) is unknown (through uncertainty in one or more of $`D_\mathrm{S}`$, $`D_\mathrm{L}`$ or cosmology), which allows the transformation $$\frac{D_\mathrm{L}D_\mathrm{S}}{D_{\mathrm{LS}}}s\frac{D_\mathrm{L}D_\mathrm{S}}{D_{\mathrm{LS}}},\mathrm{\Sigma }(\theta )s\mathrm{\Sigma }(\theta ).$$ $`(3)`$ (Here and below $`s`$ is an arbitrary constant.) G88 call (3) a similarity transformation. The only effect on observables is to multiply time delays between images by $`s`$; neither image positions nor magnifications change. If the images are not resolved, another similarity transformation, $$\theta \sqrt{s}\theta ,\beta \sqrt{s}\beta ,\mathrm{\Sigma }(\theta )s\mathrm{\Sigma }(\theta ),$$ $`(4)`$ (not explicitly considered by G88) becomes possible. Here both sources and images are rescaled by $`\sqrt{s}`$, so magnifications are unaffected, while once again time delays get multiplied by $`s`$. To avoid a degeneracy of names, I suggest calling Eq. (3) a ‘distance degeneracy’ and Eq. (4) an ‘angular degeneracy’, reserving ‘similarity transformations’ for the whole category. The distance and angular degeneracies are independent, in the sense that it is possible to break one without breaking the other. Clearly, one can combine this pair to invent other pairs of independent similarity transformations. One such pair, which I suggest calling the ‘parallax’ and ‘perspective’ degeneracies, are motivated as follows. Consider the effect of parallax, i.e., moving the observer. Say the observer moves transverse to the optical axis by $`𝐫_{\mathrm{obs}}`$. For the observer, the lens will move by $`𝐫_{\mathrm{obs}}/D_\mathrm{L}`$ and the source by $`𝐫_{\mathrm{obs}}/D_\mathrm{S}`$, which amounts to keeping the $`\theta `$ fixed and moving $`\beta `$ by $`𝐫_{\mathrm{obs}}D_{\mathrm{LS}}/(D_\mathrm{L}D_\mathrm{S})`$. Applying this change to the arrival time (2) and discarding terms with no $`\theta `$-dependence gives $$t(\theta )=(1+z_\mathrm{L})\left[\frac{D_\mathrm{L}D_\mathrm{S}}{cD_{\mathrm{LS}}}\left(\frac{1}{2}\theta ^2\theta \beta \right)\frac{1}{c}𝐫_{\mathrm{obs}}\theta \frac{8\pi G}{c^3}^2\mathrm{\Sigma }(\theta )\right].$$ $`(5)`$ If $`𝐫_{\mathrm{obs}}`$ is known and non-zero, the transformations (3) and (4) are not allowed individually, but the mixture $$\theta s\theta ,\beta s\beta ,\frac{D_\mathrm{L}D_\mathrm{S}}{D_{\mathrm{LS}}}s^1\frac{D_\mathrm{L}D_\mathrm{S}}{D_{\mathrm{LS}}},\mathrm{\Sigma }(\theta )s\mathrm{\Sigma }(\theta )$$ $`(6)`$ is still possible. Here the magnifications will depend on $`𝐫_{\mathrm{obs}}`$, but Eq. (6) does not change them because it rescales $`\theta `$ and $`\beta `$ equally. I suggest calling Eq. (6) the perspective degeneracy because it preserves the product of the distance and angular scales. Meanwhile, the similarity transformation $$\theta s\theta ,\beta s\beta ,\frac{D_\mathrm{L}D_\mathrm{S}}{D_{\mathrm{LS}}}s\frac{D_\mathrm{L}D_\mathrm{S}}{D_{\mathrm{LS}}},\mathrm{\Sigma }(\theta )s^3\mathrm{\Sigma }(\theta )$$ $`(7)`$ is independent of Eq. (6) and we can think of it as the degeneracy that is broken by a parallax observation, so I suggest calling it the parallax degeneracy. One usually factors out the similarity transformation by working with a scaled arrival time surface like so $$\tau (\theta )=\frac{1}{2}(\theta \beta )^22_\theta ^2\kappa (\theta ).$$ $`(8)`$ Here the scaled arrival time $`\tau `$, the scaled surface density (or convergence) $`\kappa `$ and the operator $`_\theta ^2`$ are all dimensionless. The physical arrival time and density are $$t(\theta )=(1+z_\mathrm{L})\frac{D_\mathrm{L}D_\mathrm{S}}{cD_{\mathrm{LS}}}\tau (\theta ),\mathrm{\Sigma }(\theta )=\frac{c^2}{4\pi G}\frac{D_\mathrm{S}}{D_{\mathrm{LS}}D_\mathrm{L}}\kappa (\theta ).$$ $`(9)`$ The usual lensing potential is $`\psi =2_\theta ^2\kappa `$ and the bending angle is $`\alpha =_\theta \psi `$. 2.2. The mass-sheet degeneracy We now rewrite (8) by discarding a $`\frac{1}{2}\beta ^2`$ term, since it is constant over the arrival-time surface, and using $`_\theta ^2\theta ^2=4`$, to get $$\tau (\theta )=2_\theta ^2(1\kappa )\theta \beta .$$ $`(10)`$ The transformation $$1\kappa s(1\kappa ),\beta s\beta .$$ $`(11)`$ clearly just rescales time delays while keeping the image structure the same; but since the source plane is rescaled by $`s`$ all magnifications are scaled by $`1/s`$, leaving relative magnifications unchanged. The effect on the lens is to rescale the lensing mass and then add or subtract a constant mass sheet. G88 call (11) a magnification transformation, but ‘mass-sheet degeneracy’ is its usual name nowadays. For a circular lens, the mass-sheet degeneracy preserves the total mass inside an Einstein radius $`\theta _\mathrm{E}`$. We can see this by invoking the two-dimensional analog of Gauss’s flux law in electrostatics, which in lens notation becomes $$\alpha \times 𝑑𝐥=2\kappa d^2\theta ,$$ $`(12)`$ or that the normal component of $`\alpha `$, integrated along any closed loop, is proportional to the enclosed mass. Along an Einstein ring, $`\alpha `$ is always radial and hence normal to the ring; also, its magnitude always equals $`\theta _\mathrm{E}`$ (since a source at the centre is imaged onto the ring). Hence, the left hand integral in Eq. (12) depends only on $`\theta _\mathrm{E}`$. Meanwhile the right hand integral gives twice the enclosed mass. Thus, fixing the Einstein radius fixes the enclosed mass. The mass-sheet degeneracy is broken if there are sources at more than one redshift. The reason is that we can no longer factor out the source-redshift dependence as we did in Eqs. (8) and (9). Instead, we can replace (8) and (9) with $$\tau (\theta )=\frac{1}{2}(\theta \beta )^22\frac{D_{\mathrm{LS}}}{D_\mathrm{S}}_\theta ^2\kappa (\theta ),t(\theta )=(1+z_\mathrm{L})\frac{D_\mathrm{L}}{c}\tau (\theta ),\mathrm{\Sigma }(\theta )=\frac{c^2}{4\pi G}\frac{1}{D_\mathrm{L}}\kappa (\theta ),$$ $`(13)`$ and replace Eq. (10) with $$\tau (\theta )=2_\theta ^2\left(1\frac{D_{\mathrm{LS}}}{D_\mathrm{S}}\kappa \right)\theta \beta ,$$ $`(14)`$ Sources at different redshifts imply simultaneous equations of the type (14) but with different factors of $`D_{\mathrm{LS}}/D_\mathrm{S}`$, which prevents a transformation like (11). 2.3. Other degeneracies G88 discuss one other transformation, which they call ‘prismatic’, consisting of adding the same constant to both the source position and the bending angle. Physically, this amounts to adding a very massive lens at very large transverse distance while pushing the source in the opposite direction. So it is not as important as the similarity and magnification transformations. Figure 1. Illustration of the mass-disk degeneracy, showing the surface density (lower panel) and the arrival time (upper panel) for three circular lenses. The units, except for $`\kappa `$, are arbitrary. The arrival time indicates a saddle point (looking like a local minimum in this cut), a maximum, and a minimum. The dashed curves correspond to a non-singular isothermal lens. Stretching the time scale amounts to making lens profile steeper (dotted curves) and shrinking the time scale amounts to making the lens profile shallower (solid curves). Note that there is a limit to stretching, because otherwise $`\kappa `$ will become negative somewhere in the region with images—negative $`\kappa `$ outside that region can always be avoided by adding an external monopole. But there is no limit to shrinking. Clearly, one can concoct any number of localized transformations that leave $`t(\theta )`$ and its derivatives unchanged at all image positions and do not make $`\kappa `$ negative anywhere. An obvious one is what we may call a ‘monopole’ transformation: any circularly symmetric redistribution of mass inwards of all observed images, and any circularly symmetric change in mass outside all observed images will have no effect on observables. A more subtle example, which causes an ambiguity between close and wide binary lenses in Local Group lensing, is discussed by Dominik (1999). The monopole transformation has an important indirect effect: it changes the mass-sheet degeneracy into a ‘mass-disk degeneracy’—as long as the disk is larger than the region of images, a circular disk and an infinite sheet are equivalent in lensing—and a much more dangerous effect, since it cannot be eliminated by the requirement that $`\kappa `$ goes to 0 at large $`\theta `$. Figure 1 illustrates. 3. Digression: velocity dispersions Though not a lensing observable, velocity dispersion is often measured in connection with lensing, and is worth discussing here. In any lens having approximately critical density, a typical internal velocity $`v`$ satisfies $$vc\theta _\mathrm{E}^{\frac{1}{2}}.$$ $`(15)`$ To derive this, we write $`R`$ for the lens’s size and $`M`$ for its mass, and recall that $`\theta _\mathrm{E}R/D_\mathrm{L}`$ for critical density, $`v^2GM/R`$ from the virial theorem, and that $`\theta _\mathrm{E}GM/(c^2D_\mathrm{L})`$. The most familiar example of the scaling (15) is for an isothermal lens. As is well known, this lens derives from a stellar dynamical sphere with constant velocity dispersion $`\sigma `$: the density is $`\rho =\sigma ^2/(2\pi Gr^2)`$, which amounts to a projected density of $`\mathrm{\Sigma }=\sigma ^2/(2GD_\mathrm{L}^2\theta ^2)`$, leading to $$\theta _\mathrm{E}=4\pi \frac{\sigma ^2}{c^2}\frac{D_\mathrm{L}}{D_\mathrm{S}}.$$ $`(16)`$ One can use Eq. (16) to define a formal $`\sigma `$ for any approximately circular lens. This formal $`\sigma `$ can usefully serve as a surrogate for $`\theta _\mathrm{E}`$. Moreover, because of the relation (15) the formal $`\sigma `$ will be of order the internal velocities in the lens, but in general it will be different from the actual velocity dispersion. To elaborate, let us consider the relation between observed velocity dispersions and mass distribution. For a stellar system with no rotation or other streaming motions, the virial theorem states that $`v^2=𝐫\mathrm{\Phi }`$, where $`v`$ is the stellar velocity, the averages $`..`$ are over the stellar distribution function, and $`\mathrm{\Phi }`$ is the total gravitational potential. Thus far there are no symmetry assumptions. If, however, spherical symmetry does apply then the line-of-sight direction must contribute the same as the orthogonal directions, and hence $$v_{\mathrm{los}}^2=\frac{1}{3}𝐫\mathrm{\Phi },$$ $`(17)`$ $`v_{\mathrm{los}}`$ being the line-of-sight stellar velocity. If $`\mathrm{\Phi }`$ is due to an isothermal sphere with dispersion $`\sigma `$ then $`𝐫\mathrm{\Phi }=2\sigma ^2`$ everywhere, which gives $$v_{\mathrm{los}}^2=\frac{2}{3}\sigma ^2$$ $`(18)`$ (cf. equation 4.6 in Kochanek 1993). For other spherical lenses one may compare $`v_{\mathrm{los}}^2`$ with the formal $`\sigma ^2`$ derived from Eq. (16). For example, consider a homogeneous sphere of stars of radius $`R`$, with no non-stellar matter. Using (17), we have $$v_{\mathrm{los}}^2=\frac{GM}{5R}$$ $`(19)`$ where $`M`$ is the total mass. If this sphere is barely compact, $`M`$ and $`R`$ satisfy $`4GMD_{\mathrm{LS}}/(c^2D_\mathrm{L}D_\mathrm{S})=\theta _\mathrm{E}^2`$ and $`R/D_\mathrm{L}=\theta _\mathrm{E}`$. Inserting these values into (19) and then eliminating $`\theta _\mathrm{E}`$ using (16) gives $$v_{\mathrm{los}}^2=\frac{\pi }{5}\sigma ^2,$$ $`(20)`$ only 6% different from (18). The above relations imply two things: (i) Eq. (17) indicates that there is a considerable range allowed in $`v_{\mathrm{los}}^2`$ for lenses with given formal $`\sigma `$, and (ii) Eq. (20) shows that even if $`v_{\mathrm{los}}^2`$ is observed to have the expected isothermal value, it does not follow that the lens is isothermal. And all this uncertainty is present without even considering ellipticity and velocity anisotropy. In summary, for galaxy and cluster lenses an order-of-magnitude relation of the type $$\theta _\mathrm{E}2^{\prime \prime }\times \frac{v_{\mathrm{los}}^2}{(300\mathrm{km}\mathrm{s}^1)^2}$$ $`(21)`$ is useful, but velocity dispersion is not a precise constraint unless the mass distribution is already known. Perhaps lenses become much better constrained if there is much more detailed velocity information; the answer seems unknown, but see Dejonghe & Merritt (1992) and Romanowsky & Kochanek (1999). 4. Appearances of the degeneracies in Local Group Lensing In Local Group lensing, the similarity transformations are relevant. The mass-disk degeneracy does not apply because absolute magnifications are always measured, and anyway there are no disk-like lens components involved. In most Local Group microlensing events only the magnification as a function of time is measured; the distances are unknown and the image structure is unresolved, so both distance and angular degeneracies (alternatively, both parallax and perpective degeneracies) apply. Note that the events being time-dependent and hence furnishing a whole sequence of arrival-time surfaces does not prevent the similarity transformations—for each event one can scale the whole sequence of arrival-time surfaces by the same factor. In a few cases ($`15`$ out of $`500`$ events observed so far) one degeneracy has been broken through additional observational information, and there are prospects for breaking the degeneracies completely in future with the help of observations from satellites. The requirements for degeneracy-breaking are well known, but it is interesting to interpret them in terms of the similarity transformations. 4.1. Proper motions Proper motion measurements break the angular degeneracy (4), leaving the spatial degeneracy (3). The significance of proper motion measurements was already appreciated by Refsdal (1966b), though the configuration then envisaged (lensing of visible one star by another, with separate proper-motion measurements for both) is not now considered realistic. A more realistic situation, independently pointed out by Gould (1994a), Nemiroff & Wickramasinghe (1994), and Witt & Mao (1994), is of a lens transiting the source star, in which case the finite size of the source will flatten the peak of the light curve; if the angular size of the source star can be estimated then $`d\beta /dt`$, and hence $`\theta _\mathrm{E}`$, can be inferred. Alcock et al. (1997) observed such an event. Another situation which enables $`d\beta /dt`$ to be measured is when not the lens itself but a caustic of a binary lens crosses the source. Albrow et al. (1999, 2000a,b) Afonso et al. (2000) and Alcock et al. (2000) have made such measurements. 4.2. Parallaxes For a parallax observation, one needs to introduce the effect of a suitable $`𝐫_{\mathrm{obs}}`$ in Eq. (5), and this can be brought about in two ways. One way, suggested by Refsdal (1966b) and Gould (1992, 1994b, 1995), is to have separate observers, using one or more satellites. The other way, suggested by Gould (1992), is to exploit the Earth’s acceleration. Now, a constant $`d𝐫_{\mathrm{obs}}/dt`$ is irrelevant in Eq. (5) because it can be absorbed inside $`d\beta /dt`$. But a known $`d^2𝐫_{\mathrm{obs}}/dt^2`$ modifies both the magnification (photometric parallax) and the proper motion of the image centroid (astrometric parallax). Photometric parallax events have been observed by Alcock et al. (1995), Bennett et al. (1997) and Mao (1999). Parallax observations leave the perspective degeneracy (6), which holds the combination $$\frac{D_\mathrm{L}D_\mathrm{S}}{D_{\mathrm{LS}}}\mathrm{\Sigma }(\theta )$$ $`(22)`$ constant. For a single mass, (22) is $`\stackrel{~}{r}_\mathrm{E}^2`$ ($`\stackrel{~}{r}_\mathrm{E}`$ being the Einstein radius projected onto the observer plane). Thus we recover the well-known result that parallax measurements determine $`\stackrel{~}{r}_\mathrm{E}`$. 4.3. Prospects for combining proper motion and parallax As will be clear from the above (and elsewhere—Refsdal 1966b already made this point) combining proper motion and parallax measurements will lift all the degeneracies, and enable the lens mass to be solved for completely. Prospects for combined measurements have been discussed in several papers. Paczyński (1998) and Boden et al. (1998) suggest using interferometry to measure both proper motion and astrometric parallax, while Miyamoto & Yoshii (1995), Berlinski & Saha (1998) and Gould & Salim (1999) advocate using interferometry for proper motions and photometry for parallaxes. 5. Appearances of the Degeneracies in Cosmological Lensing The degeneracies in cosmological lensing are complementary to those in Local Group lensing. The angular degeneracy does not appear because there is always some resolved image structure. The distance degeneracy appears, but in a very simple way—with redshifts usually measurable, the distance factor in Eq. (2) is $`H_0^1`$ times a weak and readily-quantifiable dependence on cosmology. The main thing to worry about is the mass-disk degeneracy. The following very briefly discusses the various contexts. 5.1. In quasar microlensing In quasar microlensing the mass-disk degeneracy is actually useful! Modeling microlensing of lensed quasars involves computing lightcurves in a potential of stars plus smooth matter. In such computations (e.g., Refsdal & Stabell 1997) a standard trick uses the mass-disk degeneracy to transform away the effect of the smooth matter by rescaling the stellar masses appropriately—see Eq. (24) of Paczyński (1986), which appears to be an independent discovery of the degeneracy. The angular degeneracy is also present, because while one obviously cannot change the angular scale of the macro-images, it is not forbidden to rescale the micro-image system within each macro-image, along with the source’s proper motion. 5.2. In cluster lensing Another independent discovery of the the mass-disk degeneracy, this time in the context of cluster lensing, was by Schneider & Seitz (1994). Kaiser’s (1995) formula expressing $`\mathrm{ln}(1\kappa )`$ in terms of observable ellipticities makes the degeneracy particularly explicit: multiplying $`(1\kappa )`$ by a constant will not change the ellipticities. These and later papers have drawn considerable attention to the need to break the degeneracy, and research towards this end is active. Most of the effort is directed towards using the magnification information from number counts of background galaxies (see e.g., Taylor et al. 1998), but AbdelSalam et al. (1998) use the information that comes from having a range of source redshifts. 5.3. In quasar macrolensing Although lensing degeneracies were originally discovered in the context of quasar macrolensing, recent literature in this area (e.g., the article on ‘Modeling Galaxy Lenses’ by Blandford et al., 2000) usually does not discuss degeneracies. The reason, perhaps, is that the popular parametrized lens models have focused attention on their respective parameters and away from the global transformations that produce degeneracies. Considerable work has been done on fitting parametric models to the detailed image structure (e.g., Kochanek 1995, Kochanek et al., 2000). Such work often produces precisely constrained values for the radial density gradient and the core radius. But—and this is very important—those values are conditional upon particular parametrized lens models, because (a) the mass-disk degeneracy allows one to change the radial density gradient drastically without changing the image structure at all, and (b) the monopole degeneracy makes core radii if anything more free. Nor, as we saw in the previous section, do velocity dispersion measurements provide strong independent constraints unless the mass profile is assumed already known. Thus, the mass-disk and monopole degeneracies point to some significant uncertainties in our current knowledge of galaxy lens profiles, and hence to uncertainties in estimates of cosmological parameters from quasar lensing. The effect of the mass-disk degeneracy on estimates of $`h`$ was already fully appreciated in Falco et al. (1985). Given a lens model that reproduces all observations of a lensed quasar and its host galaxy, one is still free to stretch or shrink the scale of the arrival-time surface (i.e., $`h^1`$) using the mass-disk degeneracy—see Figure 1. There is a limit to stretching, because eventually $`\kappa `$ somewhere will reach zero; this means that there is an upper limit on the inferred $`h`$. There is no limit to shrinking the arrival-time surface: the lens can get arbitrarily close to a disk with $`\kappa =1`$ and the inferred $`h`$ will get closer and closer to zero! To prevent this happening one must incorporate some assumptions about the steepness of the mass profile. Model-builders are familiar with such behavior (see e.g., Wambsganss & Paczyński 1994, Williams & Saha 2000). Degeneracies are even more dangerous for inferences of $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ from lensing, because individual lenses contain no information on these parameters, only the ensemble of lenses does.<sup>1</sup><sup>1</sup>In a little known companion paper to the famous Refsdal (1964) on time delays and $`h`$, Refsdal (1966a) suggested that time-delays for systems at different redshifts could be put on a sort of Hubble diagram to determine the other cosmological parameters. But at present, researchers prefer to fit the redshift-dependencies of the density of multiple-image systems and the distribution of image separations, which also depend on cosmology; these two quantities are much easier to observe than time delays, but more awkward to interpret because magnification bias enters. Several researchers (Maoz & Rix 1993, Kochanek 1996, Park & Gott 1997, Chiba & Yoshii 1999) have attempted to constrain $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ from the redshift-dependence of the source density or the image separations, or both. The results are conditional upon different assumptions made by the various authors, and in particular upon very specific lens profiles. Williams (1997) studies the dependence of the image-separation statistics on lens profiles, and concludes that it is much larger than the dependence on cosmology. For example, by making the lens profiles less steep in the inner regions she can make small-separation systems more magnified and hence (because of magnification bias) more abundant at high redshifts, thus completely drowning out the effect of cosmology.<sup>2</sup><sup>2</sup>The image-separation statistics are complicated by the existence of a number of wide-separation quasar pairs at low redshifts, with no visible lens. These are currently thought to be binary quasars with no lensing involved (Kochanek et al. 1999), though spectral similarities in some cases cast doubts upon that interpretation (Small et al. 1997, Peng et al. 1999). Park & Gott (1997) find that the image-separation statistics with the wide-separation pairs included as lenses is not reproducible using power-law lenses. Williams (1997) finds that the same statistics can be reproduced if the lensing galaxies have changing logarithmic density profiles and follow the scaling laws characteristic of spirals rather than ellipticals, and concludes that the resolution of the nature of the wide-separation pairs as lenses or binaries will lead to constraints on the lensing population. Kochanek et al. (1999) say that Williams’s examples are “inconsistent with the known properties of galaxies and lenses” but do not explain which known properties. 6. Summary Lensing degeneracies can be simply understood as rescalings or other transformations of the arrival-time surface that leave various image properties unaffected. The most important of these are as follows. $``$‘Similarity transformations’ typically arise in microlensing. There are two independent ones which, depending on context, are usefully taken as: 1. a ‘distance’ degeneracy where the distance scale varies while the angular scale stays fixed; and 2. an ‘angular’ degeneracy where the angular scale varies while the distance scale stays fixed; or as 1. a ‘perspective’ degeneracy where the product of the distance and angular scale varies while the ratio stays fixed; and 2. a ‘parallax’ degeneracy where the ratio stays fixed while the product varies. $``$The ‘mass-disk degeneracy’ is typical of cosmological lensing. It rescales $$1\frac{\text{density}}{\text{critical density}}$$ within a finite disk larger than observed region, in the process rescaling the total magnification and the time delays, but otherwise leaving images unaffected. $``$‘Localized’ degeneracies do not change the arrival-time surface at image positions, but change it elsewhere; no lensing observable is altered, but other properties such as core radii and velocity dispersions may be. The most insidious of these is the mass-disk degeneracy when it appears in multiply-imaging galaxy lenses, where it translates into a serious source of uncertainty in estimates of $`h`$ from time delays, and even worse uncertainties in estimates of $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ from image statistics. I am grateful to the referee for a number of detailed comments and suggestions. References AbdelSalam, H.M., Saha, P., & Williams, L.L.R. 1998, AJ, 116, 1541 Afonso, C. et al. 2000, ApJ, 532, 340 Albrow, M. et al. 1999, ApJ 522, 1011 Albrow, M. et al. 2000a, ApJ 534, 894 Albrow, M. et al. 2000b, astro-ph/0004243 Alcock, C., et al. 1995, Apj 454, L125 Alcock, C., et al. 1997 ApJ 491, 436 Alcock, C., et al. 2000 ApJ astro-ph/9907369 Bennett, D.B. 1997, BAAS 191, 8303 Berlinski, P. & Saha, P. 1998, New Ast Rev, 42, 111 Blandford, R., Surpi, G., & Kundić, T. 2000, in Brainerd, T.G. & Kochanek, C.S. 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warning/0006/math0006058.html	ar5iv	text	# On the Existence and Temperedness of Cusp Forms for 𝑆⁢𝐿₃⁢(ℤ) ## 1 Introduction In the 1950s A. Selberg (\[Sel1\]) developed his trace formula to prove the existence of non-holomorphic, everywhere-unramified, cuspidal “Maass” forms. These are real-valued functions on the upper-half plane $`=\{x+iyy>0\}`$ which are invariant under the action of $`SL_2()`$ by fractional linear transformations. Unlike the holomorphic cusp forms, which can all be explicitly described, no Maass form for $`SL_2()`$ has ever been constructed and they are believed to be intrinsically transcendental. The non-constant Laplace eigenfunctions in $`L^2(SL_2()\backslash )`$ are all Maass forms. Since $`SL_2()\backslash `$ is noncompact, their existence is no triviality; they very likely do not exist on the generic finite-volume quotient of $``$ (see \[Sarnak\]). However, Selberg showed that for any congruence subgroup $`\mathrm{\Gamma }SL_2()`$, the discrete spectrum is as large as one can expect – namely, it obeys the same asymptotics as the spectrum of a compact surface of the same size: ###### Theorem 1.1. (Selberg) Let $`0=\lambda _0<\lambda _1\lambda _2\mathrm{}`$ be the discrete eigenvalues (with multiplicity) of the non-Euclidean laplacian $`\mathrm{\Delta }`$ on $`\mathrm{\Gamma }\backslash `$. If $`\mathrm{\Gamma }SL_2()`$ is a congruence subgroup, then $$N(T)=\mathrm{\#}\{\lambda _jT\}\frac{\text{area}(\mathrm{\Gamma }\backslash )}{4\pi }T$$ (1.1) as $`T\mathrm{}`$. Similar asymptotics hold for any compact manifold by theorems of H. Weyl and others; we will refer to such an asymptotic for $`N(T)`$ as the Weyl law for the space in question. For noncompact manifolds one can count the discrete and continuous spectra together, but it is extremely difficult to decouple the two. The main contribution of this paper is a technical novelty for separating them and proving the Weyl law for certain noncompact quotients. One can generalize Maass forms to groups other than $`SL_2()`$ in a variety of ways. We are most interested in the linear group and so will focus our attention there, to $`G=SL_3()`$ in particular. Let $`K=SO_3()`$ be a maximal compact subgroup of $`G`$, $`=G/K,\mathrm{\Gamma }=SL_3(),`$ and $`X=\mathrm{\Gamma }\backslash `$. The ring $`𝒟`$ of ($`G`$-)invariant differential operators on the symmetric space $``$ is a polynomial ring in two generators; it is explicitly described in \[Bump\], p. 32 for example. We will concentrate first on a particular element of $`𝒟`$, the laplacian $`\mathrm{\Delta }`$, which is normalized so that its continuous spectrum on $`L^2()`$ is the interval $`[1,\mathrm{})`$. Here again, the non-constant discrete eigenfunctions of $`\mathrm{\Delta }`$ in $`L^2(X)`$ are cusp forms (see Section 3 for definitions), which are the appropriate generalizations of Maass forms and our basic objects of study. By symmetry considerations, it is not difficult to prove the existence of odd cusp forms on $`SL_2()\backslash `$ and $`SL_3()\backslash `$. These comprise only half of the expected spectrum; the deeper issue is the existence of even cusp forms. The only other known cusp forms on $`SL_3()\backslash `$ are Gelbart-Jacquet lifts (\[GelJac\]) of the forms on $`SL_2()\backslash `$ that Selberg discovered (many of which have been numerically identified–see \[Hejhal\]). However, no version of the trace formula has been used for the higher rank $`SL_3()\backslash `$ to prove the Weyl law for cusp forms, ala Selberg’s Theorem 1.1.<sup>1</sup><sup>1</sup>1In \[StaWal\] a proof of Theorem 1.2 is claimed. This proof appears to be incorrect because it is based on taking a specifically-chosen sequence of functions in the trace formula of \[Wallace\]. Unfortunately that trace formula is incomplete, because in its estimation of the continuous spectrum, it ignores the poles of the intertwining operators which arise from contour shifts. This seems to be a complicated problem to fix, though doing so would lead to much more specific information about the discrete spectrum (e.g. perhaps an error term). For example, Arthur’s trace formula (\[Art2\]) computes the traces of certain integral operators over the discrete spectrum, but in it the discrete spectrum is paired with parabolic orbital integrals. Though it appears difficult, it would be very interesting to separate the two and give estimates on the size of the spectrum – especially because Arthur’s formula has been developed for general quotients. We will shortcut the trace formula to prove the Weyl law for $`SL_3()\backslash `$ as well as some qualitative results about the spectrum. ###### Theorem 1.2. Let $`0=\lambda _0<\lambda _1\lambda _2\mathrm{}`$ denote the eigenvalues of $`\mathrm{\Delta }`$ on $`SL_3()\backslash `$ with multiplicity. Then $$N(T):=\mathrm{\#}\{\lambda _jT\}\frac{vol(SL_3()\backslash )}{\mathrm{\Gamma }(7/2)}\left(\frac{T}{4\pi }\right)^{5/2}$$ (1.2) as $`T\mathrm{}`$. These are the same asymptotics that the spectra of closed five-dimensional manifolds obey. This result confirms a conjecture of Sarnak (\[Sarnak\]), who asserted that the Weyl law holds for the cuspidal spectrum of the laplacian on any congruence quotient of $`SL_n()/SO_n()`$. Theorem 1.2 is actually a special case of the more-general Theorem 5.3, a kind of equidistribution theorem which counts the two-dimensional joint spectrum of $`𝒟`$ lying in various sets. This generalization is used to establish many of the other findings below, so first we shall briefly describe the joint spectrum before stating the applications. The ring of invariant differential operators $`𝒟`$ on $``$ is commutative and we may take an orthonormal set of common eigenfunctions $$\frac{1}{\sqrt{\text{area}(SL_3()\backslash )}}=\varphi _0,\varphi _1,\varphi _2,\mathrm{}$$ (1.3) such that $$\mathrm{\Delta }\varphi _j=\lambda _j\varphi _j.$$ Each $`\varphi _j`$ in turn induces a homomorphism $$\lambda :𝒟,\text{where}D\varphi =\lambda (D)\varphi ,D𝒟.$$ This map is best described by “principal series” or “Langlands” parameters $$\{\mathrm{}_1,\mathrm{}_2,\mathrm{}_3\},\mathrm{}_1+\mathrm{}_2+\mathrm{}_3=0$$ (see Section 4 for more details). It is possible to compute $`\lambda (D)`$ in terms of these spectral parameters, e.g. $$\lambda (\mathrm{\Delta })=1\frac{\mathrm{}_1^2+\mathrm{}_2^2+\mathrm{}_3^2}{2}.$$ (1.4) The real parts of the $`\mathrm{}_i`$ are bounded, for example by a trivial bound from representation theory: $$\|\text{Re }\mathrm{}_i\|<\frac{1}{2}(\text{[JacSha]}).$$ (1.5) Hence the Weyl law essentially counts the number of spectral parameters $`(\mathrm{}_1,\mathrm{}_2,\mathrm{}_3)`$ lying in a large ball. In Theorem 5.3 we compute the asymptotics of how many points of the joint spectrum lie in certain sets of various shapes and sizes. Along with some representation theory, we can count how many cusp forms obey given properties. The statement of Theorem 5.3 is technical so we will only describe its corollaries here. A cusp form is called tempered if each $`\text{Re }\mathrm{}_i=0`$. The archimedean Ramanujan-Selberg conjecture asserts that all cusp forms on $`SL_3()\backslash `$ are tempered. The only nontrivial improvement of (1.5) on individual forms is the result of \[LRS\]. ###### Theorem 1.3. Almost all of the cusp forms on $`SL_3()\backslash SL_3()/SO_3()`$ obey the archimedean Ramanujan-Selberg conjectures in the sense that $$\underset{T\mathrm{}}{lim}\frac{\mathrm{\#}\{\lambda _jT\mathrm{\Delta }\varphi _j=\lambda _j\varphi _j,\varphi _j\text{ tempered}\}}{\mathrm{\#}\{\lambda _jT\}}=1.$$ For $`SL_2()\backslash `$ the nontempered cusp forms are those with Laplace eigenvalue less than $`1/4`$, and the Weyl law shows there are only a finite number of these. In higher rank, though, there may be nontempered forms with large Laplace eigenvalues. The archimedean Ramanujan-Selberg conjecture is an analog of Selberg’s famous “1/4 conjecture,” as it automatically implies that every cuspidal Laplace eigenvalue satisfies $$\lambda ^{cusp}(SL_3()\backslash )\lambda _1()=1.$$ In fact, this has already been proved in \[Miller\], but without establishing temperedness itself. The technique here is different and further-reaching in that it works whenever the Weyl law can be proved by our method. ###### Corollary 1.4. If $$N^{cusp}(T):=\mathrm{\#}\{\lambda _jT\mathrm{\Delta }\varphi _j=\lambda _j\varphi _j\text{ for a nonzero cusp form }\varphi _j\},$$ then $$N^{cusp}(T)N(T)\frac{vol(SL_3()\backslash )}{\mathrm{\Gamma }(7/2)}\left(\frac{T}{4\pi }\right)^{5/2}$$ as $`T\mathrm{}`$. Proof: The only other eigenfunctions are residues of Eisenstein series, which are not tempered, and Theorem 1.3 shows these have measure zero. Another and more-precise reason uses the classification of the discrete spectrum established by Mœglin and Waldspurger (\[MœWal\], confirming the conjecture of \[Jacquet\]). A consequence is that the only non-cuspidal discrete eigenfunctions of $`\mathrm{\Delta }`$ on $`SL_p()\backslash SL_p()/SO_p()`$ are constants when $`p`$ is prime. Thus for $`T0`$, $`N(T)=N^{cusp}(T)+1`$, as alluded to earlier. $`\mathrm{}`$ ###### Corollary 1.5. There are “native” cusp forms on $`SL_3()\backslash `$ which are not Gelbart-Jacquet lifts of Maass forms from $`SL_2()\backslash `$. Moreover, these native forms comprise 100% of the spectrum in the sense that $$\underset{T\mathrm{}}{lim}\frac{\mathrm{\#}\{\lambda _jT\mathrm{\Delta }\varphi _j=\lambda _j\varphi _j,\varphi _j\text{ a Gelbart-Jacquet lift}\}}{\mathrm{\#}\{\lambda _jT\}}=0.$$ A cusp form is called self-dual if $`\varphi (g)=\varphi ((g^t)^1)`$ for all $`gG`$. Alternatively, its joint spectral parameters are $$\{\mathrm{}_1,\mathrm{}_2,\mathrm{}_3\}=\{\mu ,0,\mu \}\text{ for some }\mu .$$ ###### Theorem 1.6. There exists a non-self-dual cusp form on $`SL_3()\backslash SL_3()/SO_3()`$, and in fact these are also of full measure: $$\underset{T\mathrm{}}{lim}\frac{\mathrm{\#}\{\lambda _jT\mathrm{\Delta }\varphi _j=\lambda _j\varphi _j,\varphi _j\text{ self-dual}\}}{\mathrm{\#}\{\lambda _jT\}}=0.$$ All Gelbart-Jacquet lifts are self-dual and it is conjectured that the converse is true. In the next section we present a self-contained account of the main ideas of this paper, but specialized towards giving a new proof of Selberg’s Theorem 1.1 for $`\mathrm{\Gamma }=SL_2()`$. The main improvement is a circumvention of the usual derivation of the trace formula. Usually one integrates the automorphic kernel only over a truncated fundamental domain, in order to prevent divergence from the parabolic orbital integrals and the Eisenstein series. After matching the growth rates of these two terms, a trace formula is obtained in the limit as the truncation parameter moves to infinity through the cusp. Here we merely truncate, and by positivity considerations arrive at an inequality giving a lower bound on the size of the spectrum. This gives the Weyl lower bound, but with the wrong constant – it involves the area of the truncated fundamental domain instead. By pushing the truncation parameter to infinity only at the end, we recover the correct lower bound; the upper bound has already been established by \[Donnelly\]. ### Acknowledgements I wish to thank Peter Sarnak, under whose direction some of this material first appeared in my Princeton University dissertation. Also I am indebted to Donna Belli, Don Blasius, Sol Friedberg, Serge Lang, Alex Lubotzky, Jonathan Rogawski, Ze’ev Rudnick, Nolan Wallach, Gregg Zuckerman, and the referee for helpful comments. This work was supported by NSF graduate and post-doctoral fellowships, the Yale Hellman fund, and NSA grant MDA904-99-1-0046. ## 2 The Weyl Law for $`SL_2()\backslash `$ This is a self-contained section illustrating our method on $`SL_2()\backslash `$. We will reprove the Weyl Law (Theorem 1.1) using the same “partial trace” technique we use in higher rank for Theorem 1.2. We stress that the reason the partial trace is successful has little to do with the complications of higher rank per se, but rather because it serves as a substitute for the trace formula in counting the spectrum. As of yet, it appears difficult to give a direct proof of the Weyl law in higher rank using a trace formula, though this is how and why Selberg proved it for $`SL_2()\backslash `$ to begin with. Our technique is also applicable to congruence covers of $`SL_2()\backslash `$, but we will focus on the full-level situation here in order to demonstrate how the main difficulty – non-compactness – is addressed. Basic references for this section are \[Sel1\], \[Sel2\], \[Hejhal\], and \[Terras\]. ### 2.1 Definitions The hyperbolic plane can be modeled by the complex upper-half plane $$=\{u+ivv>0\}$$ with area element $$dA=\frac{dudv}{v^2},$$ line element $$d\sigma ^2=\frac{du^2+dv^2}{v^2},$$ and laplacian $$\mathrm{\Delta }=v^2\left(\frac{^2}{u^2}+\frac{^2}{v^2}\right).$$ The group $`G=SL_2()`$ acts by isometries on $``$ by fractional linear transformations $$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right):z\frac{az+b}{cz+d}.$$ The stabilizer of any point in $``$ is a maximal compact subgroup isomorphic to $`K=SO_2()=\text{Stab}(i)`$, and so $``$ is isomorphic to the symmetric space $`G/K`$. Let $`\mathrm{\Gamma }=SL_2(),\overline{\mathrm{\Gamma }}=PSL_2()`$ and $`X=\mathrm{\Gamma }\backslash =\overline{\mathrm{\Gamma }}\backslash `$. The space $`X`$ is noncompact but has finite volume, $`\frac{\pi }{3}`$. Thus the laplacian $`\mathrm{\Delta }`$ has a continuous spectrum on $`L^2(X)`$ furnished by Eisenstein series $$E_s(z)=\frac{1}{2}\underset{\{\left(\begin{array}{cc}1& n\\ 0& 1\end{array}\right)n\}\backslash \mathrm{\Gamma }}{}Im(\gamma z)^s.$$ (This infinite series is only defined for $`Re(s)>1`$, where it converges absolutely, but it has a meromorphic continuation to $`s`$.) The continuous spectrum comes from $$E_{\frac{1}{2}+it}(z),t.$$ Also, since $`\mathrm{\Delta }`$ annihilates constant functions, it has a discrete spectrum. Our aim here is to count it. Let $$0=\lambda _0<\lambda _1\lambda _2\mathrm{}$$ be the discrete spectrum of $`\mathrm{\Delta }`$ on $`L^2(X)`$ and $$\sqrt{\frac{3}{\pi }}=\varphi _0,\varphi _1,\varphi _2,\mathrm{}$$ be an orthonormal set of eigenfunctions $$\varphi _jL^2(X),\mathrm{\Delta }\varphi _j=\lambda _j\varphi _j.$$ We will prove ###### Theorem 2.1. (=Theorem 1.1)(Selberg) As $`T\mathrm{}`$ $$N(T)=\mathrm{\#}\{\lambda _jT\}\frac{T}{12}.$$ ### 2.2 Convolution Operators Every $`gC_c^{\mathrm{}}(K\backslash )=C_c^{\mathrm{}}(K\backslash G/K)`$ acts on $`fL^2(G/K)`$ by convolution: $$(L_gf)(x)=(fg)(x)=_{G/K}f(y)g(y^1x)𝑑y.$$ (2.1) Bi-$`K`$-invariant functions are functions of the distance of a point in $``$ to $`i`$. Their convolution operators form a commutative algebra which also commutes with the Laplace operator. These operators share common eigenfunctions as well, so Laplace eigenfunctions on $`=G/K`$ play a special role. The simplest examples are furnished by power functions $$\mathrm{\Delta }v^s=s(1s)v^s$$ (2.2) and by the bi-$`K`$-invariant “spherical” functions $$\stackrel{~}{\varphi }_s(z)=_K\text{Im }(kz)^s𝑑k,_K𝑑k=1,$$ $$\mathrm{\Delta }\stackrel{~}{\varphi }_s(z)=s(1s)\stackrel{~}{\varphi }_s.$$ (2.3) The integral operators $`L_g`$ also act on $`fL^2(\mathrm{\Gamma }\backslash G/K)`$: $`(L_gf)(x)=`$ $`{\displaystyle _{G/K}}f(y)g(y^1x)𝑑y`$ $`=`$ $`{\displaystyle \underset{\gamma \overline{\mathrm{\Gamma }}}{}}{\displaystyle _{\mathrm{\Gamma }\backslash G/K}}f(\gamma ^1y)g(y^1\gamma x)𝑑y`$ $`=`$ $`{\displaystyle _{\mathrm{\Gamma }\backslash G/K}}f(y)K(x,y)𝑑y,`$ (2.4) where $$K(x,y)=\underset{\gamma \overline{\mathrm{\Gamma }}}{}g(x^1\gamma y)$$ (2.5) is called the automorphic kernel. Again, the $`L_g`$ commute with each other and with $`\mathrm{\Delta }`$, so we may assume that the $`\varphi _j`$ are taken to be an orthonormal set of eigenfunctions of all $`L_g,gC_c^{\mathrm{}}(K\backslash G/K)`$ as well: $$(L_g\varphi )(x)=_{G/K}\varphi (y)g(y^1x)𝑑y=\widehat{g}(\varphi )\varphi (x).$$ (2.6) Selberg showed that $$_K\varphi _j(kg)𝑑k=\varphi _j(I)\stackrel{~}{\varphi }_{s_j}(g),\mathrm{\Delta }\varphi _j=s_j(1s_j)\varphi _j$$ and hence by averaging (2.6) over $`K`$ he found a formula for $`\widehat{g}(\varphi )`$: ###### Proposition 2.2. (Selberg’s Uniqueness Principle) The $`L_g`$-eigenvalue $`\widehat{g}(\varphi _j)`$ of any $`\varphi _j`$ as above depends only on its Laplace eigenvalue. Namely, if $$\mathrm{\Delta }\varphi _j=(\frac{1}{4}\nu _j^2)\varphi _j,$$ (2.7) then $$L_g\varphi _j=\widehat{g}(\nu _j)\varphi _j,$$ where the Selberg transform $$\widehat{g}(\nu )=_{}g(u+iv)v^{1/2+\nu }\frac{dudv}{v^2}.$$ (2.8) ###### Remark 2.3. Formula (2.8) is in fact (2.6) applied to the eigenfunction $`v^{1/2+\nu }`$, and evaluated at the identity. The kernel $`K(x,y)`$ thus has two expansions: ###### Proposition 2.4. (Spectral Expansion) For $`gC_c^{\mathrm{}}(K\backslash G/K)`$, $$K(x,y)=\underset{\gamma \overline{\mathrm{\Gamma }}}{}g(y^1\gamma x)=$$ $$\underset{\stackrel{j0}{\text{disc. }L^2\text{ spec.}}}{}\widehat{g}(\nu _j)\varphi _j(x)\varphi _j(y)+\frac{1}{4\pi }_{}\widehat{g}(it)E_{1/2+it}(x)\overline{E_{1/2+it}(y)}𝑑t.$$ (2.9) We will now recall some analytic properties of the Selberg transform that will be needed in constructing our choices of $`g`$ later. First is the inversion formula $$g(x)=_{}\widehat{g}(it)\stackrel{~}{\varphi }_{1/2it}(x)\frac{t\mathrm{tanh}\pi t}{4\pi }𝑑t.$$ (2.10) Secondly, we need to clarify the relationship between the Selberg transform and the Mellin transform $``$. If $$\overline{g}(v)=_{}g(u+iv)𝑑u$$ (2.11) denotes the “Harish transform,” then (2.8) expresses $$\widehat{g}(\nu )=\left(\frac{\overline{g}(x)}{\sqrt{x}}\right)(\nu ).$$ (2.12) Through the Mellin inversion formula, formula (2.10) essentially amounts to a statement about the Harish transform. A more-precise fact was proven by Ehrenpreis and Mautner (see also \[Gangolli\]). ###### Theorem 2.5. (\[EhrMau\]). The Harish transform $`g\overline{g}`$ is a bijection between $$C_\sigma ^{\mathrm{}}(K\backslash )=\{gC_c^{\mathrm{}}(K\backslash )supp(g)Ki(\sigma ^1,\sigma )\}$$ and $$C_\sigma ^{\mathrm{}}(0,\mathrm{})=\{hC_c^{\mathrm{}}(0,\mathrm{})h(v)=h(v^1),supp(h)(\sigma ^1,\sigma )\}.$$ ### 2.3 Construction of Functions Recall that there exist smooth, non-negative functions on $``$ with arbitrarily-small support whose Fourier transforms are also non-negative. (The easiest way to make one is to convolve a non-negative, compactly supported function with itself and rescale.) Using (2.12) and Theorem 2.5, we may fix some $`gC_\sigma ^{\mathrm{}}(K\backslash )`$ such that $`g0`$, $`_{}\widehat{g}(it)𝑑t=1`$, and $`\widehat{g}(it)0`$ for $`t`$. We would like to scale $`\widehat{g}`$ so that it resembles the characteristic function of a large interval. However, in this non-euclidean setting it is difficult to control $`g`$ and $`\widehat{g}`$ simultaneously. Instead, we will convolve $`\widehat{g}(it)`$ with $`\chi _{[T,T]}`$: $$\widehat{g}_T(it)=_T^T\widehat{g}(it+ir)𝑑r0,$$ (2.13) which is the Mellin transform of $$\frac{\overline{g}(x)}{\sqrt{x}}_T^Tx^{ir}𝑑rC_\sigma ^{\mathrm{}}(0,\mathrm{}).$$ Hence Theorem 2.5 also guarantees the existence of some function $`g_TC_\sigma ^{\mathrm{}}(K\backslash )`$ whose Selberg transform $$\widehat{g_T}=\widehat{g}_T.$$ We will use the $`g_T`$’s in the automorphic kernel. Our construction allows us to conclude two important analytic properties: ###### Proposition 2.6. $$\underset{x}{\mathrm{max}}\|g_T(x)\|=g_T(i).$$ Proof: This follows from the inversion formula (2.10), the positivity of $`\widehat{g}_T`$ (2.13), and the inequality $$\|\stackrel{~}{\varphi }_{1/2+it}(x)\|1=\stackrel{~}{\varphi }_{1/2+it}(i)$$ (see \[DKV\], (2.13)). $`\mathrm{}`$ ###### Proposition 2.7. For any $`m0`$, $$\widehat{g}_T(it)=\{\begin{array}{cc}1+O_m\left((\|T\|\|t\|+1)^m\right)\hfill & t<T,\hfill \\ O_m\left((\|t\|\|T\|+1)^m\right)\hfill & tT.\hfill \end{array}$$ Proof: Since $`g`$ is smooth, $`\widehat{g}(it)=O\left((1+\|t\|)^m\right)`$ for any $`m0`$. Hence $$_T^{\mathrm{}}\widehat{g}(it)𝑑t=O\left((1+T)^m\right),T>0$$ and the proposition follows because of our normalization $`_{}\widehat{g}(it)𝑑t=1.`$ $`\mathrm{}`$ ### 2.4 Partial Trace Let $`C1`$ be a fixed truncation parameter. We will adjust it only in the final step of our analysis. Denote the usual fundamental domain for $`X`$ as $$=\{z\|z+\overline{z}\|1\|z\|\}$$ and its truncation as $$_C=\{zIm(z)C\}.$$ Of course, $`_C`$ is compact and has area $`\frac{\pi }{3}\frac{1}{C}`$; as $`C\mathrm{}`$, it exhausts $``$. If $`X`$ were compact, we would take the trace of the integral operator $`L_g`$ by integrating the two expressions for $`K(x,x)`$ in (2.9) over $``$. However, these integrals diverge and we instead take a partial trace over $`_C`$: $$_{}K(x,x)𝑑x=\underset{\gamma \overline{\mathrm{\Gamma }}}{}__Cg(x^1\gamma x)𝑑x$$ $$=\underset{j0}{}\widehat{g}(\nu _j)__C\|\varphi _j(x)\|^2𝑑x+\frac{1}{4\pi }_{}\widehat{g}(it)__C\|E_{1/2+it}(x)\|^2𝑑x𝑑t.$$ (2.14) Up until now we have essentially followed Selberg’s derivation of the trace formula. His next step is the understand the divergence of each side of (2.14) as a function of $`C`$. He is able to cancel it from each side and derive a trace formula for $`\widehat{g}(\nu _j)`$ by letting $`C\mathrm{}`$. As a final step, he uses a family of functions such as the $`g_T`$ to deduce the spectral asymptotics. This cancellation can be found in great generality (e.g. \[Art2\]), but the resulting formula is quite complicated. We will now deviate from this path by keeping $`C`$fixed for now, and performing an analysis of (2.14) with the $`g_T`$. By noting $$__C\|\varphi _j(x)\|^2𝑑x_{}\|\varphi _j(x)\|^2𝑑x=1,$$ (2.15) (2.14) becomes an inequality, giving a lower bound on $`\widehat{g}(\nu _j)`$ (which, as we will make precise later, roughly equals $`N(T^2)`$). Since only a finite number of $`\mathrm{\Gamma }`$-translates of $`_C`$ neighbor it, we may take $`\sigma `$ to be small enough so that $$g(x^1\gamma x)0,x_C$$ only for the finitely-many $`\gamma \overline{\mathrm{\Gamma }}`$ which have a fixed-point on the boundary of $`_C`$. ###### Lemma 2.8. For $`\sigma `$ small, $$\left\|\underset{\stackrel{\gamma \overline{\mathrm{\Gamma }}}{\gamma I}}{}__Cg(x^1\gamma x)𝑑x\right\|=O\left(\sigma ^2\mathrm{max}g\right).$$ Proof: By the above remark, it suffices to show $$_{}\|g(x^1\gamma x)\|𝑑x=O\left(\sigma ^2\mathrm{max}g\right)$$ for $`\gamma `$ a rotation. Since rotations sweep distant points a proportional amount, only $`x`$ of distance $`O(\sigma )`$ from the fixed point can have $`g(x^1\gamma x)0`$. $`\mathrm{}`$ Applying (2.14), (2.15), and Proposition 2.6 we have that $`{\displaystyle __C}g_T(x^1Ix)`$ $`{\displaystyle \underset{j0}{}}\widehat{g}_T(\nu _j)`$ (2.16) $`+{\displaystyle \frac{1}{4\pi }}{\displaystyle _{}}\widehat{g}_T(it){\displaystyle __C}\|E_{1/2+it}(x)\|^2𝑑x𝑑t`$ $`+O(\sigma ^2g_T(I)).`$ The next subsection is devoted to analyzing the Eisenstein series term. In Lemma 2.10 we show that it is $`O(T\mathrm{log}T)`$. Thus for some $`c_0>0`$ $$\left(\text{area}(_C)c_0\sigma ^2\right)g_T(I)=\left(\text{area}(_C)c_0\sigma ^2\right)\frac{1}{4\pi }_{}\widehat{g}_T(it)t\mathrm{tanh}(\pi t)𝑑t$$ $$\underset{j0}{}\widehat{g}_T(\nu _j)+O(T\mathrm{log}T).$$ Proof of Theorem 2.1: We need to show $$\underset{T\mathrm{}}{lim}\frac{N(T)}{T}=\frac{1}{12}.$$ The general theorem of \[Donnelly\], for example, shows that $$\underset{T\mathrm{}}{lim\; sup}\frac{N(T)}{T}\frac{1}{12}.$$ (2.17) By Proposition 2.7 $`{\displaystyle _{}}\widehat{g}_T(it)t\mathrm{tanh}(\pi t)𝑑t=`$ $`{\displaystyle _T^T}t\mathrm{tanh}(\pi t)𝑑t+O(T)`$ $`=`$ $`T^2+o(T^2).`$ (2.18) Again using (2.17), $$\underset{j0}{}\widehat{g}_T(\nu _j)N((T+\sqrt{T})^2)+O(1).$$ This proves $$\underset{T\mathrm{}}{lim\; inf}\frac{N(T)}{T}\frac{\frac{\pi }{3}\frac{1}{C}c_0\sigma ^2}{4\pi }.$$ Taking $`\sigma 0`$ and $`C\mathrm{}`$, we conclude the Weyl law. $`\mathrm{}`$ ### 2.5 The Eisenstein Series Bounds Recall that the Eisenstein series $`E_s(z)`$ has the constant term $$c(v,s)=_0^1E_s(u+iv)𝑑u=v^s+\varphi (s)v^{1s},$$ where $$Z(s)=\pi ^{s/2}\mathrm{\Gamma }(\frac{s}{2})\zeta (s),$$ and $$\varphi (s)=\frac{Z(2s1)}{Z(2s)}=\sqrt{\pi }\frac{\mathrm{\Gamma }(s\frac{1}{2})\zeta (2s1)}{\mathrm{\Gamma }(s)\zeta (2s)}.$$ Define the truncated Eisenstein series for $`u+iv`$ as $$\mathrm{\Lambda }^CE_s(u+iv)=\{\begin{array}{cc}E_s(u+iv)c(v,s),\hfill & v>C\hfill \\ E_s(u+iv),\hfill & vC.\hfill \end{array}$$ (2.19) The truncated $`\mathrm{\Lambda }^CE_s(z)`$ can be extended from $``$ to $`\mathrm{\Gamma }\backslash `$; it decays rapidly as $`v\mathrm{}`$, and so is in $`L^2(\mathrm{\Gamma }\backslash )`$. Thus for any $`s_1,s_2`$ which are not poles of the Eisenstein series $`E_s(z)`$, $$_{\mathrm{\Gamma }\backslash }\mathrm{\Lambda }^CE_{s_1}(z)\mathrm{\Lambda }^CE_{s_2}(z)\frac{dudv}{v^2}<\mathrm{}.$$ In fact the same is true if we only truncate $`E_{s_1}(z)`$: ###### Lemma 2.9. With $`s_1`$ and $`s_2`$ as above, $$_{\mathrm{\Gamma }\backslash }\mathrm{\Lambda }^CE_{s_1}(z)\mathrm{\Lambda }^CE_{s_2}(z)\frac{dudv}{v^2}=_{\mathrm{\Gamma }\backslash }\mathrm{\Lambda }^CE_{s_1}(z)E_{s_2}(z)\frac{dudv}{v^2}.$$ Proof: Unfolding, their difference $$_{\mathrm{\Gamma }\backslash }\mathrm{\Lambda }^CE_{s_1}(z)\left(E_{s_2}(z)\mathrm{\Lambda }^CE_{s_2}(z)\right)\frac{dudv}{v^2}=_C^{\mathrm{}}\frac{dv}{v^2}c(y,s_2)_0^1\left(E_{s_1}(z)c(v,s_1)\right)𝑑u$$ $$=_C^{\mathrm{}}\frac{dv}{v^2}c(v,s_2)\left(c(v,s_1)c(v,s_1)\right)=0.$$ $`\mathrm{}`$ It will simplify notation to introduce $`\delta _C(z)`$, the characteristic function of $`\{z\text{Im }z>C\}`$. Then for $`z`$ $$\mathrm{\Lambda }^CE_s(z)=E_s(z)\delta _C(z)c(v,s)=\underset{\gamma \mathrm{\Gamma }_{\mathrm{}}\backslash \mathrm{\Gamma }}{}\left[\text{Im }(\gamma z)^{s_1}\delta _C(\gamma z)c(\text{Im }(\gamma z),s_1)\right].$$ Because the rightmost expression is automorphic, it agrees with $`\mathrm{\Lambda }^CE_s(z)`$ for all $`z`$; at most one of the terms indicated by the $`\delta _C`$ functions is actually subtracted. We can thus unfold this series representation of $`\mathrm{\Lambda }^CE_{s_1}(z)`$ in the integral: $$_{\mathrm{\Gamma }\backslash }\mathrm{\Lambda }^CE_{s_1}(z)E_{s_2}(z)\frac{dudv}{v^2}=_{\mathrm{\Gamma }_{\mathrm{}}\backslash }E_{s_2}(z)\left(v^{s_1}\delta _C(z)c(v,s_1)\right)\frac{dudv}{v^2}$$ $$=_0^{\mathrm{}}\frac{dv}{v^2}\left(v^{s_1}\delta _C(z)c(v,s_1)\right)_0^1𝑑uE_{s_2}(z)$$ $$=_0^{\mathrm{}}\left(v^{s_2}+\varphi (s_2)v^{1s_2}\right)\left(v^{s_1}\delta _C(v)\left(v^{s_1}+\varphi (s_1)v^{1s_1}\right)\right)\frac{dv}{v^2}$$ $$=_0^Cv^{s_2+s_12}𝑑v+\varphi (s_2)_0^Cv^{s_1s_21}𝑑v$$ $$\varphi (s_1)_C^{\mathrm{}}v^{s_2s_1+1}𝑑v\varphi (s_1)\varphi (s_2)_C^{\mathrm{}}v^{s_1s_2}𝑑v$$ $$=\frac{C^{s_2+s_11}}{s_2+s_11}+\frac{\varphi (s_2)C^{s_1s_2}}{s_1s_2}+\frac{\varphi (s_1)C^{s_2s_1}}{s_2s_1}+\varphi (s_1)\varphi (s_2)\frac{C^{1s_2s_1}}{1s_2s_1},$$ (2.20) provided $`\text{Re }s_1>1+\text{Re }s_2`$. Each side has a meromorphic continuation in $`s_1,s_2`$, and together are the Maass-Selberg relations. Because the Eisenstein series and hence their truncations are holomorphic for $`Re(s_1)=Re(s_2)=\frac{1}{2}`$, the Maass-Selberg relations must have a removable singularity there. The limiting value is $$_{\mathrm{\Gamma }\backslash }\|\mathrm{\Lambda }^CE_{1/2+it}(z)\|^2\frac{dudv}{v^2}=2\mathrm{log}C\frac{\varphi ^{}}{\varphi }(\frac{1}{2}+it)+\frac{C^{2it}\varphi (\frac{1}{2}it)C^{2it}\varphi (\frac{1}{2}+it)}{2it}.$$ (2.21) The functional equations $$Z(s)=Z(1s),Z(\overline{s})=\overline{Z(s)}$$ show that $$\|\varphi (\frac{1}{2}+it)\|=\left\|\frac{Z(it)}{Z(it+1)}\right\|=1.$$ (2.22) Also, $$\frac{\varphi ^{}}{\varphi }(\frac{1}{2}+it)=2\left(\frac{Z^{}}{Z}(2it)+\frac{Z^{}}{Z}(2it)\right)$$ (2.23) and $$\frac{Z^{}}{Z}(s)=\underset{\rho Z(\rho )=0}{}\frac{1}{s\rho }\frac{1}{s}\frac{1}{s1},$$ (2.24) a sum over the poles and critical zeroes of $`\zeta (s)`$. ###### Lemma 2.10. $$_{}\widehat{g}_T(it)__C\|E_{1/2+it}(x)\|^2𝑑x𝑑t=O(T\mathrm{log}T).$$ Proof: By the definition of the truncation, $$__C\|E_{1/2+it}(x)\|^2𝑑x_{}\|\mathrm{\Lambda }^CE_{1/2+it}(x)\|^2𝑑x,$$ which is calculated in (2.21). In view of (2.22) and (2.23), it suffices to show $$\text{Re }_{}\widehat{g}_T(it)\frac{Z^{}}{Z}(2it)𝑑t=O(T\mathrm{log}T).$$ This follows from (2.24) and the zero counting estimate (see Proposition 7.1 or \[Titchmarsh\]) $$\text{Re }_{i(H1)}^{i(H+1)}\underset{\rho Z(\rho )=0}{}\frac{ds}{s\rho }=O(\mathrm{\#}\{\rho H1\text{Im }(\rho )H+1\})=O(\mathrm{log}H).$$ $`\mathrm{}`$ ## 3 Coordinates on $`SL_3()`$ and $`=SL_3()/SO_3()`$ The Iwasawa decomposition $`G=SL_3()=NAK`$ states that every element $`gG`$ can be uniquely factored as $`g=nak`$, where $$nN_0=\left\{\left(\begin{array}{ccc}1& & \\ 0& 1& \\ 0& 0& 1\end{array}\right)G\right\},$$ $$aA_0=\left\{\left(\begin{array}{ccc}a_1& 0& 0\\ 0& a_2& 0\\ 0& 0& a_3\end{array}\right)G,a_i>0\right\},$$ and $$kK=SO_3().$$ The minimal standard parabolic is $$P_0=N_0A_0=A_0N_0=\left\{\left(\begin{array}{ccc}& & \\ 0& & \\ 0& 0& \end{array}\right)G\right\}.$$ There are two associate maximal standard parabolics $$P_1=\left\{\left(\begin{array}{ccc}& & \\ & & \\ 0& 0& \end{array}\right)G\right\},P_2=\left\{\left(\begin{array}{ccc}& & \\ 0& & \\ 0& & \end{array}\right)G\right\}.$$ Each can be decomposed uniquely as $`P_i=M_i^{}N_iA_i=N_iM_i^{}A_i`$, where $$N_1=\left\{\left(\begin{array}{ccc}1& 0& \\ 0& 1& \\ 0& 0& 1\end{array}\right)G\right\},$$ $$M_1^{}=\left\{\left(\begin{array}{ccc}& & 0\\ & & 0\\ 0& 0& \pm 1\end{array}\right)G\right\},$$ $$A_1=\left\{\left(\begin{array}{ccc}a& 0& 0\\ 0& a& 0\\ 0& 0& a^2\end{array}\right)G,a>0\right\},$$ $$N_2=\left\{\left(\begin{array}{ccc}1& & \\ 0& 1& 0\\ 0& 0& 1\end{array}\right)G\right\},$$ $$M_2^{}=\left\{\left(\begin{array}{ccc}\pm 1& 0& 0\\ 0& & \\ 0& & \end{array}\right)G\right\},$$ and $$A_2=\left\{\left(\begin{array}{ccc}a^2& 0& 0\\ 0& a^1& 0\\ 0& 0& a^1\end{array}\right)G,a>0\right\}.$$ The period of an automorphic function $`\psi L^2(\mathrm{\Gamma }\backslash G/K)`$ along the parabolic $`P`$ is defined as the integral $$\psi _P(g)=_{\mathrm{\Gamma }N\backslash N}\psi (ng)𝑑n.$$ (3.1) If all periods along each of the parabolics $`P_0,P_1,`$ and $`P_2`$ vanish, then $`\psi `$ is called a cusp form. We shall use the following sets of coordinates for the Lie algebras. Let $$𝐚_0=\{\left(\begin{array}{ccc}h_1& 0& 0\\ 0& h_2& 0\\ 0& 0& h_3\end{array}\right)h_1,h_2,h_3,h_1+h_2+h_3=0\}$$ $$\{H=(h_1,h_2,h_3)^3h_1+h_2+h_3=0\}.$$ It has simple roots $`\alpha _1,\alpha _2𝐚_0^{}`$, which are the following linear functions on $`𝐚_0`$: $$\alpha _1(H)=h_1h_2,\alpha _2(H)=h_2h_3.$$ There is also a third root $`\alpha _3=\alpha _1+\alpha _2`$ which acts by $`\alpha _3(H)=h_1h_3`$. We can similarly coordinatize $$𝐚_0^{}=\{\lambda =(\mathrm{}_1,\mathrm{}_2,\mathrm{}_3)^3\mathrm{}_1+\mathrm{}_2+\mathrm{}_3=0\}$$ and its complexification $$𝐚_0^{}=\{\lambda =(\mathrm{}_1,\mathrm{}_2,\mathrm{}_3)^3\mathrm{}_1+\mathrm{}_2+\mathrm{}_3=0\}.$$ Since these are subsets of $`^3`$, we will use the standard norm there to define norms in these spaces. The simple roots are bases of $`𝐚_0^{}`$ and $`𝐚_0^{}`$ as vector spaces over $``$ and $``$, respectively. Also, they have corresponding co-roots $$\alpha _1^{}=(1,1,0)𝐚_0,\alpha _2^{}=(0,1,1)𝐚_0.$$ The Weyl groups $`\mathrm{\Omega }(𝐚_0)`$ and $`\mathrm{\Omega }(𝐚_0^{})`$ are isomorphic to the symmetric group $`S_3`$, and act in a compatible way by permuting the standard basis vectors in $`^3`$ or $`^3`$ (see the chart in the appendix for more details). The Cartan subgroups $`A_1,A_2`$ of the maximal parabolics have one-dimensional Lie algebras, coordinatized by $$𝐚_\mathrm{𝟏}=\{H=(h,h,2h)h\}\{h\}$$ and $$𝐚_\mathrm{𝟐}=\{H=(2h,h,h)h\}\{h\}.$$ These each have one-dimensional dual spaces, $`𝐚_{\mathrm{𝟏}}^{}{}_{}{}^{}`$ and $`𝐚_{\mathrm{𝟐}}^{}{}_{}{}^{}`$ respectively. Also, each dual has a special vector $`\rho `$, half the sum of the positive roots: $$\rho _0=\frac{\alpha _1+\alpha _2+\alpha _3}{2}:(h_1,h_2,h_3)h_1h_3,$$ $$\rho _1:(h,h,2h)3h,$$ and $$\rho _2:(2h,h,h)3h.$$ Under these coordinates $`\rho _1`$ and $`\rho _2`$ are naturally identified with $`\alpha _2`$ and $`\alpha _1`$, respectively. Each parabolic $`P_i`$ has the Langlands decomposition $$P_i=N_iA_iM_i^{}$$ and logarithm maps $`H_i:G𝐚_𝐢`$ such that $$a=e^{H_i(a)},aA_i$$ (3.2) and $$gN_ie^{H_i(g)}M_i^{}K.$$ These maps are well-defined despite the fact that the decomposition $`G=P_iK`$ is in general not unique. Similarly for the maximal parabolics, there are maps $`m_1`$ and $`m_2`$ mapping $`G`$ onto $`M_1^{}/(KM_1^{})`$ and $`M_2^{}/(KM_2^{})`$, respectively. ### Truncation Let $`\mathrm{\Delta }_P`$ denote the set of simple roots which do not vanish identically on $`P`$: $$\mathrm{\Delta }_{P_0}=\{\alpha _1,\alpha _2\},\mathrm{\Delta }_{P_1}=\{\alpha _2\},\mathrm{\Delta }_{P_2}=\{\alpha _1\}.$$ The group $`G`$ can also be viewed as a parabolic, with $`\mathrm{\Delta }_G=\{\}.`$ We can now define Langlands’ truncation (\[Lan1\], see \[Art1\]). For any parabolic $`P=P_0,P_1,P_2`$, or $`G`$, define $`\widehat{\tau }_P(x)`$ to be the characteristic function of $$\{x=c_1\alpha _1^{}+c_2\alpha _2^{}𝐚_0c_i>0\alpha _i\mathrm{\Delta }_P\}.$$ Thus, $`\widehat{\tau }_{P_0}`$ is the characteristic function of $$\{x=c_1\alpha _1^{}+c_2\alpha _2^{}𝐚_0c_1,c_2>0\},$$ $$\widehat{\tau }_{P_1}\text{ of }\{x=c_1\alpha _1^{}+c_2\alpha _2^{}𝐚_0c_2>0\},$$ $$\widehat{\tau }_{P_2}\text{ of }\{x=c_1\alpha _1^{}+c_2\alpha _2^{}𝐚_0c_1>0\},$$ and $$\widehat{\tau }_G\text{ of }𝐚_0.$$ Let $`C𝐚_0`$ be a fixed parameter. The truncation of an automorphic form $`\psi `$ is a sum over all standard parabolic subgroups $$(\mathrm{\Lambda }^C\psi )(x):=\underset{P}{}(1)^{dimA}\underset{\gamma \mathrm{\Gamma }P\backslash \mathrm{\Gamma }}{}\widehat{\tau }_P(H_0(\gamma x)C)_{\mathrm{\Gamma }N\backslash N}\psi (n\gamma x)𝑑n$$ $$=\underset{P}{}(1)^{dimA}\underset{\gamma \mathrm{\Gamma }P\backslash \mathrm{\Gamma }}{}\widehat{\tau }_P(H_0(\gamma x)C)\psi _P(\gamma x),$$ (3.3) which itself is clearly an automorphic form. It can be proven that it also decays rapidly in the cusp because of the way its constant terms have been removed (see \[Art1\]). Note that if $`\psi `$ is a cusp form to begin with, by definition all its constant terms $`\psi _P`$ vanish identically in proper parabolics, and thus $`\mathrm{\Lambda }^C\psi =\psi `$. ### Spectral Density We will later need to use the spectral density function $`\beta (\lambda )`$ on $`i𝐚_{\mathrm{𝟎}}^{}{}_{}{}^{}`$, which is a constant multiple of $$\left\|\frac{\mathrm{\Gamma }(\frac{1}{2}+\lambda (\alpha _1^{}))}{\mathrm{\Gamma }(\lambda (\alpha _1^{}))}\frac{\mathrm{\Gamma }(\frac{1}{2}+\lambda (\alpha _2^{}))}{\mathrm{\Gamma }(\lambda (\alpha _2^{}))}\frac{\mathrm{\Gamma }(\frac{1}{2}+\lambda (\alpha _3^{}))}{\mathrm{\Gamma }(\lambda (\alpha _3^{}))}\right\|^2.$$ (3.4) This function will appear below in (4.2); see \[DKV\] for more information. For $`\lambda `$ such that $`\lambda (\alpha _1^{}),\lambda (\alpha _2^{})`$, and $`\lambda (\alpha _3^{})`$ are all large, Stirling’s formula shows that $`\beta (\lambda )`$ behaves as a constant times $`\|\lambda (\alpha _1^{})\lambda (\alpha _2^{})\lambda (\alpha _3^{})\|`$. The spectral density function is normalized so that $$_{\lambda T}\beta (\lambda )𝑑\lambda \frac{T^{5/2}}{\mathrm{\Gamma }(7/2)(4\pi )^{5/2}}.$$ One can also view this as normalizing the measure $`d\lambda `$. ## 4 Convolution operators Every $`gC_c^{\mathrm{}}(K\backslash G/K)`$ acts by convolution on $`L^2(G/K)`$: $$(L_gf)(x)=(fg)(x)=_Gf(y)g(y^1x)𝑑y=_{G/K}f(y)g(y^1x)𝑑y.$$ The convolution operator $`L_g`$ also acts on $`fL^2(\mathrm{\Gamma }\backslash G/K)`$ by $$(L_gf)(x)=_{G/K}f(y)g(y^1x)𝑑y$$ $$=\underset{\gamma \mathrm{\Gamma }}{}_{\mathrm{\Gamma }\backslash G/K}f(\gamma ^1y)g(y^1\gamma x)𝑑y=_{\mathrm{\Gamma }\backslash G/K}f(y)K(x,y)𝑑y,$$ where $$K(x,y)=\underset{\gamma \mathrm{\Gamma }}{}g(y^1\gamma x)$$ is the automorphic kernel. Suppose furthermore that $`g(x)`$ is real; then $`L_g`$ is an operator on $`L^2(\mathrm{\Gamma }\backslash G/K)`$ which commutes with the laplacian $`\mathrm{\Delta }`$ and all other invariant differential operators. We may thus choose an orthonormal set $$\frac{1}{\sqrt{vol(\mathrm{\Gamma }\backslash G/K)}}=\varphi _0,\varphi _1,\varphi _2,\mathrm{}$$ of common eigenfunctions of $`L_g`$ and the ring of invariant differential operators $`𝒟`$. Continuing, we may resolve any multiplicities that remain by using the Hecke operators, which commute with these other operators. ###### Proposition 4.1. (Selberg’s Uniqueness Principle)(\[Sel1\]) If $`\varphi `$ is a common eigenfunction of the ring of invariant differential operators, then $$(L_g\varphi )(x)=\widehat{g}(\varphi )\varphi (x),$$ where $`\widehat{g}(\varphi )`$ only depends on $`\varphi `$’s eigenvalues under $`𝒟`$. In fact, one can always find some $`\lambda 𝐚_{}^{}`$ such that the function $`\varphi _\lambda =e^{(\lambda +\rho )(H_0(x))}`$ on $`G/K`$ has the same eigenvalues as $`\varphi `$ under any $`D𝒟`$. This provides the formula $$\widehat{g}(\varphi )=\widehat{g}(\lambda )=(L_g\varphi _\lambda )(I)=_{G/K}g(x)e^{(\lambda +\rho )(H_0(x))}𝑑x.$$ (4.1) If $`s`$ is any permutation in the Weyl group $`\mathrm{\Omega }(𝐚_0^{})`$, then in fact $`\varphi _\lambda (x)=e^{(\lambda +\rho )(H_0(x))}`$ has the same eigenvalues as $`\varphi _{s\lambda }`$. The uniqueness principle thus implies that $`\widehat{g}`$ is invariant under $`\mathrm{\Omega }(𝐚_0^{})`$. Formula (4.1) shows that the transform $`\widehat{g}(\lambda )`$ is the composition of an average over $`N_0`$, and a Fourier transform on $`A_0`$. More precisely, if $`g`$ is a function on $`G/K`$, denote $$\overline{g}(a)=_{N_0}g(na)𝑑n.$$ For a function $`f`$ on $`A_0`$, define the Fourier transform as $$(f)(\lambda )=_{A_0}f(a)e^{\lambda (H_0(a))}𝑑a,\lambda i𝐚_\mathrm{𝟎}^{}.$$ Then $`\widehat{g}=\left(\overline{g}()e^{\rho (H_0())}\right)`$. We will make use of the following theorem in constructing our choices of functions $`g`$. Before stating it, let us introduce the notation $$A(\sigma )=\{aA_0H_0(a)\sigma \}.$$ ###### Theorem 4.2. (\[Gangolli\]) The map $`g\overline{g}`$ provides a bijection between the sets $$\{gC^{\mathrm{}}(K\backslash G/K)suppgKA(\sigma )K\}$$ and $$\{hC^{\mathrm{}}(A)supphA(\sigma )\text{ and }h(sa)=h(a)\text{ for all }s\mathrm{\Omega }(𝐚_\mathrm{𝟎})\}.$$ ### Construction of functions Let $`\sigma >0`$ be fixed for the remainder of this section. It is possible to find a non-negative function $`gC^{\mathrm{}}(K\backslash G/K)`$ which is supported in $`KA(\sigma )K`$, and whose transform $`\widehat{g}`$ is non-negative on the joint spectrum. Having such a function without the positivity requirement on $`\widehat{g}`$ is straightforward; one then rescales and convolves it with itself to achieve positivity (see lemma 6.2 of \[DKV\] for example). We shall again normalize $`g`$ so that $$_{i𝐚_\mathrm{𝟎}^{}}\widehat{g}(\lambda )𝑑\lambda =1.$$ Let us state a property of $`g`$ which is implied by the positivity of $`\widehat{g}`$ on $`i𝐚_{\mathrm{𝟎}}^{}{}_{}{}^{}`$: ###### Proposition 4.3. If $`gC_C^{\mathrm{}}(K\backslash G/K)`$ is such that $`\widehat{g}(\lambda )0`$ for all $`\lambda i𝐚_{\mathrm{𝟎}}^{}{}_{}{}^{}`$, then $$\underset{xG}{\mathrm{max}}\|g(x)\|=g(I).$$ ###### Remark 4.4. As the proof will demonstrate, the analogous fact is true for the usual Euclidean Fourier transform. Proof of Proposition 4.3: The proof uses the inversion formula $$g(x)=_{i𝐚_{\mathrm{𝟎}}^{}{}_{}{}^{}}\widehat{g}(\lambda )\overline{\stackrel{~}{\varphi }_\lambda (x)}\beta (\lambda )𝑑\lambda ,$$ (4.2) where $$\|\stackrel{~}{\varphi }_\lambda (x)\|\stackrel{~}{\varphi }_\lambda (I)=1$$ is a spherical function (e.g. (2.3)) and $`\beta (\lambda )`$ is the spectral density. Trivially $$\|g(x)\|_{i𝐚_{\mathrm{𝟎}}^{}{}_{}{}^{}}\widehat{g}(\lambda )\beta (\lambda )𝑑\lambda =g(I).$$ $`\mathrm{}`$ Now if $`\mathrm{\Sigma }`$ is a measurable, bounded subset of $`i𝐚_{\mathrm{𝟎}}^{}{}_{}{}^{}`$ which is invariant under the Weyl group, define $$\widehat{g}_\mathrm{\Sigma }(\lambda )=\widehat{g}\chi _\mathrm{\Sigma }(\lambda )=_\mathrm{\Sigma }\widehat{g}(\lambda \mu )𝑑\mu .$$ (4.3) The function $`\widehat{g}_\mathrm{\Sigma }(\lambda )`$ is roughly concentrated on $`\mathrm{\Sigma }`$, especially for large $`\mathrm{\Sigma }`$; we shall use it to estimate the spectrum in $`\mathrm{\Sigma }`$. Also $`\widehat{g}_\mathrm{\Sigma }(\lambda )`$ decays rapidly, since $`g_\mathrm{\Sigma }`$ is a smooth function of compact support. If $`\mathrm{\Sigma }`$ is open and its boundary has a finite Hausdorff length, then $$\widehat{g}_{t\mathrm{\Sigma }}(\lambda )(1+dist(\lambda ,t\mathrm{\Sigma }))^m$$ (4.4) (cf. p. 85 of \[DKV\]). Of course $`\widehat{g}_\mathrm{\Sigma }(\lambda )`$ is the Fourier transform of $$\overline{g}(a)_\mathrm{\Sigma }e^{(\mu \rho )(H_0(a))}𝑑\mu ,$$ which is a smooth function on $`A_0`$ whose support is contained in $`A(\sigma )`$. It is furthermore invariant under the action of the Weyl group $`\mathrm{\Omega }(𝐚_\mathrm{𝟎})`$ because $`\mathrm{\Sigma }`$ is. Thus by Theorem 4.2 there exists a smooth, bi-$`K`$-invariant function $`g_\mathrm{\Sigma }`$ supported in $`KA(\sigma )K`$ such that $$\overline{g_\mathrm{\Sigma }}(a)=\overline{g}(a)_\mathrm{\Sigma }e^{\mu (H_0(a))}𝑑\mu $$ (4.5) and $$\widehat{g_\mathrm{\Sigma }}=\widehat{g}_\mathrm{\Sigma }.$$ In summary, while not changing the support of our function $`g`$, we can still smear its transform $`\widehat{g}`$ over $`\mathrm{\Sigma }.`$ Of course $`g_\mathrm{\Sigma }`$ becomes more concentrated near the identity as $`\mathrm{\Sigma }`$ gets larger; in the classical Euclidean case this is analogous to multiplying a function by Fejer’s kernel. In (5.5) we will use a result analogous to Proposition 2.7 comparing $`\widehat{g}_\mathrm{\Sigma }`$ to $`\chi _\mathrm{\Sigma }`$. ## 5 The Partial Trace In addition to the discrete spectrum there is also a continuous spectrum, furnished by Eisenstein series. Because it is complicated we will just refer to its terms in the spectral expansion as “$`Eis_g(x,y)`$” until we need to be more explicit: $$K(x,y)=\underset{j0}{}\widehat{g}(\lambda _j)\varphi _j(x)\varphi _j(y)+Eis_g(x,y)=\underset{\gamma \mathrm{\Gamma }}{}g(x^1\gamma y).$$ (5.1) Let $``$ be a fundamental domain for $`\mathrm{\Gamma }\backslash G/K`$ and $`_C`$ be a compact subset of $``$. Then $`{\displaystyle __C}K(x,x)𝑑x=`$ $`{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}{\displaystyle __C}g(x^1\gamma x)𝑑x`$ $`=`$ $`{\displaystyle \underset{j0}{}}\widehat{g}(\lambda _j){\displaystyle __C}\varphi _j(x)^2𝑑x+{\displaystyle __C}Eis_g(x,x)𝑑x`$ (5.2) Now, $$__C\varphi _j(x)^2𝑑x_{}\varphi _j(x)^2𝑑x=1$$ so $$\underset{\gamma \mathrm{\Gamma }}{}__Cg(x^1\gamma x)𝑑x\underset{j0}{}\widehat{g}(\lambda _j)+__CEis_g(x,x)𝑑x.$$ (5.3) Now we will analyze the integrals $$__Cg(x^1\gamma x)𝑑x$$ more systematically. Note that we do not group these into conjugacy classes as is done in the trace formula. ###### Proposition 5.1. For $`\sigma `$ sufficiently small and $`suppgA(\sigma )`$, the integral $$__Cg(x^1\gamma x)𝑑x=0$$ for all but finitely many $`\gamma `$ – those which have fixed points in the closure of $`_C`$. ###### Proposition 5.2. If $`\gamma I`$ has a fixed-point in the closure of $`_C`$ and $`suppgA(\sigma )`$, then $$__Cg(x^1\gamma x)𝑑x\|\mathrm{max}g\|vol(A(\sigma )).$$ Although the implied constant above depends on $`_C`$, it may be taken to be independent of $`\gamma SL_3()`$ since only finitely many elements produce a non-zero integral. Clearly $$__Cg(x^1Ix)𝑑x=vol(_C)g(I).$$ Now take $`g=g_\mathrm{\Sigma }`$ as in (4.5). Using Propositions 4.3 and 5.2, the inequality (5.3) becomes $$vol(_C)g_\mathrm{\Sigma }(I)+O\left(g_\mathrm{\Sigma }(I)vol(A(\sigma ))\right)\underset{j=0}{\overset{\mathrm{}}{}}\widehat{g}_\mathrm{\Sigma }(\lambda _j)+__CEis_{g_\mathrm{\Sigma }}(x,x)𝑑x.$$ (5.4) The proof of Theorem 8.5 in \[DKV\] shows that both $$\left\|\underset{j=0}{\overset{\mathrm{}}{}}\widehat{g}_\mathrm{\Sigma }(\lambda _j)\mathrm{\#}\{\text{Im }\lambda _j\mathrm{\Sigma }\}\right\|_{B\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda $$ and $$\left\|g_\mathrm{\Sigma }(I)_\mathrm{\Sigma }\beta (\lambda )𝑑\lambda \right\|_{B\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda ,$$ $$B\mathrm{\Sigma }=\{\lambda i𝐚_\mathrm{𝟎}^{}\text{dist}(\lambda ,\mathrm{\Sigma })1\},$$ (5.5) assuming $`\mathrm{\Sigma }`$ is a bounded, open, and Weyl-group invariant subset of $`i𝐚_\mathrm{𝟎}^{}`$. Now to control the error term we shall assume $`\mathrm{\Sigma }`$’s boundary is piece-wise smooth (cf. Lemma 8.7 of \[DKV\]). Then (3.4) and (5.5) show that for large $`t`$ $$g_{t\mathrm{\Sigma }}(I)_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda ,$$ and by using (5.4), $$\left(_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda \right)\left(vol(_C)+O(vol(A(\sigma )))\right)(1+o(1))$$ $$\mathrm{\#}\{\text{Im }\lambda _jt\mathrm{\Sigma }\}+__CEis_{g_{t\mathrm{\Sigma }}}(x,x)𝑑x.$$ We will see at the end of Section 7 that $$__CEis_{g_{t\mathrm{\Sigma }}}(x,x)𝑑x=o\left(_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda \right).$$ (5.6) Proof of Theorem 1.2: Taking $`t\mathrm{}`$ $$\underset{t\mathrm{}}{lim\; inf}\frac{\mathrm{\#}\{\text{Im }\lambda _jt\mathrm{\Sigma }\}}{_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda }vol(_C)+O(vol(A(\sigma ))).$$ We had insisted that $`suppgA(\sigma )`$, so taking $`\sigma 0`$ and exhausting $``$ through compact sets $`_C`$ we conclude that $$vol()\underset{t\mathrm{}}{lim\; inf}\frac{\mathrm{\#}\{\text{Im }\lambda _jt\mathrm{\Sigma }\}}{_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda }.$$ (5.7) If $$\mathrm{\Sigma }=\{\lambda =(\mathrm{}_1,\mathrm{}_2,\mathrm{}_3)i𝐚_\mathrm{𝟎}^{}\lambda ^21\},$$ then this indicates $$vol()\underset{T\mathrm{}}{lim\; inf}\frac{N(T)}{\left(\frac{T}{4\pi }\right)^{5/2}\frac{1}{\mathrm{\Gamma }(7/2)}}$$ because the Laplace eigenvalue of $`e^{(\lambda +\rho )(H(g))}`$ is $`1\frac{\mathrm{}_1^2+\mathrm{}_2^2+\mathrm{}_3^2}{2}`$.<sup>2</sup><sup>2</sup>2A check of this normalization is provided by the compact case, where there are no Eisenstein series and the Weyl law is already known. (Note that $`\|\text{Re }\mathrm{}_j\|<\frac{1}{2}`$ by unitary – \[JacSha\].) The upper bound $$\underset{T\mathrm{}}{lim\; sup}\frac{N(T)}{\left(\frac{T}{4\pi }\right)^{5/2}\frac{1}{\mathrm{\Gamma }(7/2)}}vol()$$ (5.8) due to \[Donnelly\] shows in fact that $$N(T)\left(\frac{T}{4\pi }\right)^{5/2}\frac{vol()}{\mathrm{\Gamma }(7/2)}.$$ (5.9) $`\mathrm{}`$ ###### Theorem 5.3. (Spectral Equidistribution) If $`\mathrm{\Sigma }`$ is a bounded, open, Weyl-group invariant subset of $`i𝐚_\mathrm{𝟎}^{}`$ with a piece-wise smooth boundary, then $$\underset{t\mathrm{}}{lim}\frac{\mathrm{\#}\{\text{Im }\lambda _jt\mathrm{\Sigma }\}}{_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda }=vol().$$ (5.10) Proof: The lower bound is in (5.7). To prove the upper bound, we will use this lower bound along with the asymptotics of the Weyl law for counting in balls. Let $`r`$ be the length of the longest vector in $`\mathrm{\Sigma }`$ (which is bounded by assumption). By defining $`\mathrm{\Sigma }_c`$ as the complement of $`\mathrm{\Sigma }`$ inside the $`r`$-ball $`B_r`$, we have that $$\mathrm{\#}\{\text{Im }\lambda _jt\mathrm{\Sigma }\}+\mathrm{\#}\{\text{Im }\lambda _jt\mathrm{\Sigma }_c\}=\mathrm{\#}\{\text{Im }\lambda _jt(\mathrm{\Sigma }\mathrm{\Sigma }_c)\}.$$ Given any $`ϵ>0`$ we can find $`t`$ large enough so that $$\mathrm{\#}\{\text{Im }\lambda _jt\mathrm{\Sigma }\}(1+ϵ)\left(_{t(\mathrm{\Sigma }\mathrm{\Sigma }_c)}\beta (\lambda )𝑑\lambda _{t\mathrm{\Sigma }_c}\beta (\lambda )𝑑\lambda \right),$$ (5.11) which implies the upper bound when $`ϵ0`$. $`\mathrm{}`$ Proofs of Theorems 1.6 and 1.3: If the cusp form $`\varphi _j`$ is not tempered, then its spectral parameter $`\lambda =(\mathrm{}_1,\mathrm{}_2,\mathrm{}_3)`$ is not a purely-imaginary vector. By the classification of the unitary dual, we have equality of the sets $$\{\mathrm{}_1,\mathrm{}_2,\mathrm{}_3\}=\{\overline{\mathrm{}}_1,\overline{\mathrm{}}_2,\overline{\mathrm{}}_3\}.$$ We know that $`\|Re(\mathrm{}_i)\|<\frac{1}{2}`$ by unitary (\[JacSha\]), so the vectors $`\lambda `$ all lie near the hyperplanes defined by $$\lambda (\alpha _1^{})=0,\lambda (\alpha _2^{})=0,\lambda (\alpha _1^{}+\alpha _2^{})=0.$$ Similarly, if $`\varphi `$ is self dual then $$\{\mathrm{}_1,\mathrm{}_2,\mathrm{}_3\}=\{\mu ,0,\mu \}\text{ for some }\mu .$$ This again constrains $`\lambda `$ to lie along a hyperplane. However, by taking the shape $`\mathrm{\Sigma }`$ in Theorem 5.3 to be further and further away from any fixed hyperplane, we conclude that no hyperplane has a positive percentage of the spectrum near it. Thus each of the exceptional sets we are considering is of measure zero compared to the rest of the spectrum. $`\mathrm{}`$ A related argument was used in \[DKV\] for cocompact subgroups $`\mathrm{\Gamma }`$. Proof of Theorem 1.5: All Gelbart-Jacquet lifts are self-dual forms on $`SL_3`$. Another proof would be simply by counting: the lift quadruples the Laplace eigenvalue of a cusp form on $`SL_2()\backslash `$, and the number of these with $`SL_3()`$-Laplace eigenvalue $`T`$ is $`O(T)`$ by Selberg’s theorem 1.1. Though forms $`\psi `$ on congruence covers of $`SL_2()\backslash `$ may also lift to $`SL_3()\backslash SL_3()/SO_3()`$, such $`\psi `$ are actually twists of forms on $`SL_2()\backslash `$, so one need only count the lifts from $`SL_2()\backslash `$ itself, and not from other congruence covers. $`\mathrm{}`$ ## 6 Eisenstein Series These periodized functions on $`\mathrm{\Gamma }\backslash G/K`$ are constructed using the Langlands decompositions of $`G`$’s parabolics. Recall that any element $`pP=NAM^{}`$ factors uniquely into a product $`p=nam`$ of elements from their respective subgroups. Writing the diagonal matrices $`a=e^{H_P(p)}`$, we get a map which extends to $`G`$: $$gH_P(g),gNe^{H_P(g)}M^{}K.$$ For the parabolics $`P_1`$ and $`P_2`$, where $`M^{}GL_2()`$, there are corresponding maps $`m:GM^{}/(KM^{})`$ (see (3.2)). If $`\lambda 𝐚_0^{}`$ and $`gG`$, the minimal parabolic Eisenstein series is defined as $$E(P_0,g,\lambda )=E(P_0,g,\lambda ,1)=\underset{\mathrm{\Gamma }P_0\backslash \mathrm{\Gamma }}{}e^{(\lambda +\rho _0)(H_0(\gamma g))}.$$ (6.1) This sum only converges when $`\alpha _1(\lambda )`$ and $`\alpha _2(\lambda )`$ have large real parts, but it has a meromorphic continuation to all of $`𝐚_0^{}`$. Since $$\mathrm{\Gamma }M_1^{}\backslash M_1^{}/(KM_1^{})GL_2()\backslash ,$$ the discrete eigenfunctions on the former are just the even<sup>3</sup><sup>3</sup>3i.e. $`f(x+iy)=f(x+iy)`$. discrete eigenfunctions for $`SL_2()`$. Take such a cusp form and $`\lambda 𝐚_1^{}`$, and define the maximal parabolic Eisenstein series $$E(P_1,g,\lambda ,\varphi )=\underset{\mathrm{\Gamma }P_1\backslash \mathrm{\Gamma }}{}e^{(\lambda +\rho _1)(H_1(\gamma g))}\varphi (m_1(\gamma g)).$$ (6.2) There is a similar Eisenstein series $`E(P_2,g,\lambda _1,\varphi )`$ for the other maximal parabolic $`P_2`$, related through a functional equation. Each series again is initially only defined for certain values of $`\lambda `$ but extends via a meromorphic continuation to $`𝐚_{}^{}`$ (\[Lan1\],\[Lan3\]). We are now in a position to define what “$`Eis_g(x,y)`$” is. The spectral expansion of the automorphic kernel is $$K(x,y)=\underset{\stackrel{\varphi _j\text{ an }L^2\text{ discrete eigenfunction}}{\text{on }SL_3()\backslash SL_3()/SO_3()}}{}\widehat{g}(\lambda _j)\varphi _j(x)\varphi _j(y)+Eis_g(x,y),$$ and up to a normalizing constant for the measure, $`Eis_g(x,y)`$ is (\[Art2\]) $$\frac{1}{3(2\pi i)^2}_{i𝐚_\mathrm{𝟎}^{}}\widehat{g}(\lambda )E(P_0,x,\lambda ,1)\overline{E(P_0,y,\lambda ,1)}𝑑\lambda +$$ $$\frac{1}{2\pi i}\underset{\stackrel{\varphi _j\text{ an even }L^2\text{ discrete eigenfunction}}{\text{on }SL_2()\backslash ,\mathrm{\Delta }\varphi _j=(\frac{1}{4}\nu _j^2)\varphi _j}}{}_{i𝐚_1^{}}\widehat{g}(\lambda +(\nu _j,\nu _j,0))E(P_1,x,\lambda ,\varphi _j)\overline{E(P_1,y,\lambda ,\varphi _j)}𝑑\lambda .$$ (6.3) We also may assume that each $`\varphi _j`$ on $``$ is a Hecke eigenform. There is a beautiful formula for the inner-products of truncated Eisenstein series, due to Langlands. It generalizes the Maass-Selberg formula (2.20) for $`SL_2()`$. See \[Art1\] for details. ###### Theorem 6.1. Langlands’ inner product formula (\[Lan1\], \[Art1\]) $$_{\mathrm{\Gamma }\backslash G/K}\left[(\mathrm{\Lambda }^CE)(P,x,\varphi ,\lambda _1)\right]\left[(\mathrm{\Lambda }^CE)(P,x,\varphi ,\lambda _2)\right]𝑑x$$ $$=\underset{PP^{}\text{ associate}}{}vol(𝐚^{}/<\alpha ^{}\alpha \mathrm{\Delta }_P^{}>)\times $$ $$\underset{s_1,s_2\mathrm{\Omega }(𝐚,𝐚^{})}{}\frac{e^{(s_1\lambda _1+s_2\lambda _2)(C)}}{_{\alpha \mathrm{\Delta }_P^{}}(s_1\lambda _1+s_2\lambda _2)(\alpha ^{})}M(s_1,\lambda _1)\varphi ,M(s_2,\lambda _2)\varphi .$$ Here $`M(s,\lambda )`$ is an intertwining operator, which sends $`\varphi `$ to a cusp form on the potentially-different parabolic $`P^{}`$. We will discuss it in the instances it arises for us, where it essentially acts as scalar multiplication. The last expression $`\psi ,\psi ^{}`$ is an inner product over $`\mathrm{\Gamma }M^{}\backslash M^{}`$, and the Weyl group $`\mathrm{\Omega }(𝐚,𝐚^{})`$ is the set of isomorphisms of $`𝐚`$ to $`𝐚^{}`$ coming from restrictions of elements in $`\mathrm{\Omega }(𝐚_0)`$. ### Interlude: $`SL_2`$ Write $`C=(c,c),c>0`$, and let $`\lambda _1=(it+ϵ,itϵ),\lambda _2=(it,it)`$ so that $$_{}\mathrm{\Lambda }^CE(P,g,\lambda _1,\varphi )E(P,g,\lambda _2,\varphi )𝑑g$$ is a constant times $$\frac{e^{(ϵ,ϵ)C}}{2ϵ}+\frac{e^{(2it+ϵ,2itϵ)C}}{4it+2ϵ}R(2it)+\frac{e^{(2itϵ,2it+ϵ)C}}{4it2ϵ}R(2it+2ϵ)+\frac{e^{(ϵ,ϵ)C}}{2ϵ}R(2it)R(2it+2ϵ).$$ Here the intertwining operator is $$R(s)=\sqrt{\pi }\frac{\mathrm{\Gamma }(\frac{s}{2})}{\mathrm{\Gamma }(\frac{s+1}{2})}\frac{\zeta (s)}{\zeta (s+1)}=\frac{Z(s)}{Z(s+1)}$$ (6.4) with $$Z(s)=\pi ^{s/2}\mathrm{\Gamma }(\frac{s}{2})\zeta (s)=Z(1s).$$ Take $`ϵ0`$ so that $`\lambda _1\overline{\lambda }_2`$. Then the last expression approaches $$_{}\|\mathrm{\Lambda }^CE(P,g,\lambda )\|^2𝑑g=(const)\left[4c\frac{R^{}}{R}(2it)+\frac{e^{4itc}}{4it}R(2it)\frac{e^{4itc}}{4it}R(2it)\right].$$ In (2.20) we derived this from a direct calculation. It is a key step in Selberg’s trace formula for $`SL_2()`$. ## 7 Bounding the Eisenstein Contribution In this final section we will complete the proofs by establishing the estimate in (5.6). From (6.3) it is sufficient to show that both $$__C_{i𝐚_\mathrm{𝟎}^{}}\widehat{g}_{t\mathrm{\Sigma }}(\lambda )\|E(P_0,x,\lambda ,1)\|^2𝑑\lambda 𝑑x=o\left(_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda \right),$$ (7.1) and $$__C_{i𝐚_\mathrm{𝟏}^{}}\underset{\stackrel{\stackrel{\varphi _j\text{ an even }L^2\text{ discrete}}{\text{ eigenfunction}}}{\stackrel{\text{on }SL_2\left(\right)\backslash }{\mathrm{\Delta }\varphi _j=(\frac{1}{4}\nu _j^2)\varphi _j}}}{}\widehat{g}_{t\mathrm{\Sigma }}(\lambda +(\nu _j,\nu _j,0))\|E(P_1,x,\lambda ,\varphi _j)\|^2d\lambda dx=o\left(_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda \right).$$ (7.2) To do this we shall use the rapid decay of $`\widehat{g}_{t\mathrm{\Sigma }}`$ in (4.4). Among other benefits, this allows us to interchange the order of integration. Keep in mind that $`\widehat{g}_{t\mathrm{\Sigma }}`$ is roughly the characteristic function of $`t\mathrm{\Sigma }`$. ### Some background on L-functions We will require some information about the density of zeroes of certain L-functions: the Riemann $`\zeta `$ function and the standard L-function of an even cusp form $`\varphi `$ on $`SL_2()\backslash `$, $`L(s,\varphi )`$. Each is defined as an Euler product<sup>4</sup><sup>4</sup>4Recall that $`\varphi `$ is assumed to be both a Hecke and Laplace eigenfunction. over the primes $$\zeta (s)=\underset{p}{}(1p^s)^1,L(s,\varphi )=\underset{p}{}(1\alpha _pp^s)^1(1\alpha _p^1p^s)^1,$$ (7.3) where the $`\alpha _p`$ satisfy the bound $`\|\alpha _p\|p^{5/28}`$ (\[BDHI\]). Each L-function can be completed: $$Z(s)=\mathrm{\Gamma }_{}(s)\zeta (s),\mathrm{\Lambda }(s,\varphi )=\mathrm{\Gamma }_{}(s+\nu )\mathrm{\Gamma }_{}(s\nu )L(s,\varphi ),$$ where $$\mathrm{\Gamma }_{}(s)=\pi ^{s/2}\mathrm{\Gamma }(s/2)$$ and $`\nu `$ is related to $`\varphi `$’s Laplace eigenvalue by $$\mathrm{\Delta }\varphi =(\frac{1}{4}\nu ^2)\varphi .$$ With this convention $`Z(s)=Z(1s),\mathrm{\Lambda }(s,\varphi )=\mathrm{\Lambda }(1s,\varphi )`$, and each is entire except for the simple poles of $`Z(s)`$ at $`s=0,1`$. The following estimate will be used to bound Eisenstein series integrals later. The analogous statement for $`Z(s)`$ is classical (e.g. see \[Titchmarsh\]) and the proof for $`\mathrm{\Lambda }(s,\varphi )`$ is essentially identical. However we do not know of a reference in the literature and include it for completeness. ###### Proposition 7.1. For $`T`$ $$\text{Re }_{T1}^{T+1}\frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}(1+it,\varphi )𝑑t\mathrm{log}(\|T\|+\|\nu \|).$$ (7.4) Proof: Using entirety of $`\mathrm{\Lambda }(s,\varphi )`$ and its Mittag-Leffler expansion, $$\frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}(s)=\frac{\mathrm{\Gamma }_{}^{}}{\mathrm{\Gamma }_{}}(s+\nu )+\frac{\mathrm{\Gamma }_{}^{}}{\mathrm{\Gamma }_{}}(s\nu )+\frac{L^{}}{L}(s,\varphi )=\underset{\{\rho \mathrm{\Lambda }(\rho )=0\}}{}\frac{1}{s\rho }$$ (the sum of the zeroes is actually only conditionally convergent, so the term with $`\rho `$ should always be summed with the term containing the zero at $`1\rho `$). From the Euler product, $$\frac{L^{}}{L}(s,\varphi )=\underset{p}{}\underset{n=1}{\overset{\mathrm{}}{}}(\alpha _p^n+\alpha _p^n)p^{ns}\mathrm{log}p,$$ so the bound $`\|\alpha _p\|p^{5/28}`$ implies that $`\|\frac{L^{}}{L}(2+it,\varphi )\|`$ is uniformly bounded in both $`\nu `$ and $`t`$. By Stirling’s formula, $$\frac{\mathrm{\Gamma }_{}^{}}{\mathrm{\Gamma }_{}}(s)=\frac{1}{2}\mathrm{log}\pi +\frac{1}{2}\mathrm{log}s+O(1/\|s\|)$$ and $$\frac{\mathrm{\Gamma }_{}^{}}{\mathrm{\Gamma }_{}}(2+it+\nu )+\frac{\mathrm{\Gamma }_{}^{}}{\mathrm{\Gamma }_{}}(2+it\nu )\mathrm{log}(\|t\|+\|\nu \|)+O(1).$$ Thus, $$\underset{\rho }{}\frac{1}{1+\|t\gamma \|^2}\underset{\rho }{}\frac{1}{2+it\rho }\mathrm{log}(\|t\|+\|\nu \|)+O(1).$$ (7.5) It follows that there are no more than $`O(\mathrm{log}(\|T\|+\|\nu \|))`$ zeroes between with imaginary part between $`T1`$ and $`T+1`$. By writing $$\text{Re }_{T1}^{T+1}\frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}(1+it,\varphi )𝑑t=\text{Re }_{T1}^{T+1}\underset{\rho }{}\frac{1}{1+it\rho }dt$$ $$=\text{Re }_{T1}^{T+1}\underset{\|\rho iT\|2}{}\frac{1}{1+it\rho }dt+\text{Re }_{T1}^{T+1}\underset{\|\rho iT\|>2}{}\frac{1}{1+it\rho }dt,$$ invoking (7.5), and using the fact that $$\text{Re }_{1+i(T1)}^{1+i(T+1)}\frac{ds}{s\rho }=O(1),$$ we bound each of the terms above by $`O(\mathrm{log}(\|T\|+\|\nu \|))`$. $`\mathrm{}`$ ### Minimal Parabolic Eisenstein Series We have appended a table of the 36 terms in the Langlands inner-product formula (also referred to as the Maass-Selberg relations) for the minimal parabolic. The calculation is aided by the identity $$M(s,(\mathrm{}_1,\mathrm{}_2,\mathrm{}_3))=\underset{\stackrel{1i<j3}{s(i)>s(j)}}{}R(\mathrm{}_i\mathrm{}_j).$$ (7.6) This is a special case of a general result of Langlands (\[Lan2\],pp. 36-47;\[Lan4\], p. 134;\[Art2\],p. 854); the function $`R`$ here is the same as the one used above in (6.4). Using a limiting procedure as in the interlude we can compute: ###### Proposition 7.2. (Diagonal terms) Let $`C=(c,0,c),\lambda _1=(it_1+ϵ_1,it_2+ϵ_2,it_3+ϵ_3),\lambda _2=(it_1,it_2,it_3)`$. Then $$\underset{ϵ_1,ϵ_2,ϵ_30}{lim}\underset{s\mathrm{\Omega }(𝐚_0)}{}\frac{e^{s(\lambda _1+\overline{\lambda }_2)(C)}}{[s(\lambda _1+\overline{\lambda }_2)(\alpha _1^{})][s(\lambda _1+\overline{\lambda }_2)(\alpha _2^{})]}M(s,\lambda _1),M(s,\lambda _2)$$ $$=3c^22c\frac{R^{}}{R}(it_1it_2)2c\frac{R^{}}{R}(it_2it_3)2c\frac{R^{}}{R}(it_1it_3)$$ $$+\frac{R^{}}{R}(it_1it_2)\frac{R^{}}{R}(it_2it_3)$$ $$+\frac{R^{}}{R}(it_1it_3)\frac{R^{}}{R}(it_2it_3)+\frac{R^{}}{R}(it_1it_2)\frac{R^{}}{R}(it_1it_3).$$ ###### Proposition 7.3. For any $`ϵ>0`$ and $`\lambda i𝐚_{\mathrm{𝟎}}^{}{}_{}{}^{}`$ with $`\lambda `$ large, we have that $$_{}\|\mathrm{\Lambda }^CE(P_0,x,\lambda )\|^2𝑑x=O_ϵ(\lambda ^ϵ).$$ Proof: We start with some estimates on $`\zeta (s)`$. The following may be found in Chapter 3 of \[Titchmarsh\]: There is an absolute constant $`\kappa >0`$ such that $$\frac{1}{\mathrm{log}(\|t\|+2)}\|\zeta (\sigma +it)\|\mathrm{log}(\|t\|+2)$$ (7.7) and $$\left\|\frac{\zeta ^{}}{\zeta }(\sigma +it)\right\|\mathrm{log}(\|t\|+2)$$ in the region $`\sigma 1\frac{\kappa }{\mathrm{log}(\|t\|+2)}`$. Thus $$R(s)=\frac{Z(1s)}{Z(1+s)}=\frac{\pi ^{(1s)/2}\mathrm{\Gamma }\left(\frac{1s}{2}\right)\zeta (1s)}{\pi ^{(1+s)/2}\mathrm{\Gamma }\left(\frac{1+s}{2}\right)\zeta (1+s)}$$ $$(1+\|\text{Im }s\|)^{\text{Re }s}\mathrm{log}(\|\text{Im }s\|+2)^2$$ and $$\frac{R^{}}{R}(s)\mathrm{log}(\|\text{Im }s\|+2)$$ for $$\|\text{Re }s\|\frac{\kappa }{\mathrm{log}(\|\text{Im }s\|+2)}.$$ Of course both $`R(s)`$ and $`\frac{R^{}}{R}(s)`$ are analytic in this region because of the nonvanishing. Fix $`\lambda =(it_1,it_2,it_3)`$. Then each of the six terms from Proposition 7.2 is trivially $`O_ϵ(\lambda ^ϵ)`$. Of the remaining thirty terms (see the chart in the appendix), some may have singularities when either of the denominators $$s(\lambda _1+\overline{\lambda }_2)(\alpha _1^{})\text{or}s(\lambda _1+\overline{\lambda }_2)(\alpha _2^{})$$ vanish. Nevertheless, the sum of all of the terms represents a holomorphic function for all $`\lambda i𝐚_{\mathrm{𝟎}}^{}{}_{}{}^{}`$, so the poles cancel with other terms. If $`\lambda `$ is such that each $$\|s(\lambda _1+\overline{\lambda }_2)(\alpha _j^{})\|\frac{\kappa }{2\mathrm{log}(\lambda +2)},$$ then each term is trivially $`O_ϵ(\lambda ^ϵ)`$ as well (the numerators have modulus one when $`\text{Re }\lambda =0`$). Otherwise, if some denominator is small, we will take estimates further away and appeal to the maximum principle. Take a small neighborhood in $$s_1=s(\lambda _1+\overline{\lambda }_2)(\alpha _1^{}),s_2=s(\lambda _1+\overline{\lambda }_2)(\alpha _2^{}).$$ Now, suppose that $$\|s_1\|\frac{\kappa }{2\mathrm{log}(\lambda +2)}\|s_2\|$$ (the other cases have almost-identical proofs). Then for $`\|s_1\|=\frac{\kappa }{2\mathrm{log}(\lambda +2)}`$, (7.7) shows that each term is $`O_ϵ(\lambda ^ϵ)`$. The maximum modulus principle shows that this bound holds uniformly for $`\|s_1\|\mathrm{log}\lambda `$. $`\mathrm{}`$ ### Maximal Parabolic $`P_1`$ Now we turn to the sum in (7.2). We will first consider the simpler case that $`\varphi `$ is a cusp form. If $`s`$ is the lone permutation in $`\mathrm{\Omega }(𝐚_1,𝐚_2)`$, then the intertwining operator acts as $$M(s,\lambda )\varphi =R(\lambda (\alpha ^{}),\varphi )\varphi ^{},$$ where $`\varphi `$ and $`\varphi ^{}`$ are cusp forms on $`M_1`$ and $`M_2`$ coming from the same even cusp form on $`SL_2()\backslash `$, and $$R(s,\varphi )=\frac{\mathrm{\Lambda }(s,\varphi )}{\mathrm{\Lambda }(s+1,\varphi )}$$ is a ratio formed from $`\varphi `$’s completed standard L-function (7.3). As before with (7.6), since $`\varphi `$ is everywhere unramified, this can be derived from Langlands’ formula (\[Lan2\],\[Lan4\]). Then the Maass-Selberg relations take the following form: ###### Proposition 7.4. The inner product $$_{}\left\|\mathrm{\Lambda }^CE(P,g,\lambda ,\varphi )\right\|^2𝑑g$$ is a constant multiple of $`3/2c\frac{R^{}}{R}(\lambda (\alpha ^{}),\varphi )`$. Proof: The inner product formula is a constant times $$\underset{ϵ0}{lim}\frac{e^{(2it+2ϵ)c}}{2it+2ϵ}\frac{e^{(2it+2ϵ)c}}{2it+2ϵ}\frac{R(2it+2ϵ)}{R(2it)}.$$ $`\mathrm{}`$ ###### Proposition 7.5. The integral $$_{t=T1}^{T+1}_{}\left\|\mathrm{\Lambda }^CE(P,g,\lambda ,\varphi _j)\right\|^2𝑑g𝑑t=O(log(T+\lambda _j)).$$ Proof: Using the previous proposition, this requires only the estimate on $`_{T1}^{T+1}\frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}(1+it,\varphi )𝑑t`$ in Proposition 7.1 (cf. (2.23)). $`\mathrm{}`$ ###### Remark 7.6. (On summing over the eigenvalues) The formula (6.3) includes a sum over the $`SL_2()\backslash `$ spectrum as well. Through a more-precise statement of Selberg’s Theorem 1.1 such as $$N(T)=\frac{1}{12}T^2+O(T\mathrm{log}T),$$ we can bound the spectral points $`\nu _j`$ in the interval $`[\nu 1,\nu +1]`$ by $`O(\sqrt{\|\nu \|}\mathrm{log}\|\nu \|)`$. ### The Constant Function on $`SL_2`$ There is only one Eisenstein series left to estimate, $$E(P_1,g,\lambda ,1),$$ the Eisenstein series induced from the constant function on the maximal parabolic $`P_1`$. The constant functions on $`SL_2()\backslash `$ may be viewed as multiples of $$Res_{s=1}E_s(z)=\frac{1}{2}Res_{s=1}\underset{(c,d)=1}{}\frac{y^s}{\|cz+d\|^{2s}}.$$ We may ignore the actual value of the constant since we are only trying to get an order-of-magnitude estimate. Thus, $`E(P_1,g,\lambda ,1)`$is a constant multiple of the residue $$Res_{\delta =0}E(P_0,g,(\frac{1}{2}+it+i\delta ,\frac{1}{2}+iti\delta ,2it),1).$$ (7.8) Taking the residue of the inner product $$_{}\|\mathrm{\Lambda }^CE(P_1,g,\lambda ,1)\|^2𝑑g$$ is more complicated, though a limiting value must exist since the Eisenstein series is meromorphic there. We will explicitly see this cancellation occurring. Slicker arguments are possible, but we will present a detailed proof for the sake of clarity. Let $$\lambda _1=(\frac{1}{2}+it+i\delta _1+ϵ,\frac{1}{2}+iti\delta _1+ϵ,2it2ϵ),$$ and $$\lambda _2=(\frac{1}{2}it+i\delta _2,\frac{1}{2}iti\delta _2,2it).$$ Then we are interested in $$\underset{\delta _1,\delta _20}{lim}\delta _1\delta _2\underset{s_1,s_2\mathrm{\Omega }(𝐚_0)}{}\frac{e^{(s_1\lambda _1+s_2\lambda _2)(C)}}{(s_1\lambda _1+s_2\lambda _2)(\alpha _1^{})(s_1\lambda _1+s_2\lambda _2)(\alpha _2^{})}M(s_1,\lambda _1)M(s_2,\lambda _2).$$ This expression is holomorphic in $`\delta _1`$ and $`\delta _2`$ near zero, so it does not matter how we take the limits $`\delta _1,\delta _20`$. We will first take the residue in $`\delta _1`$. Of the 36 terms, we will of course ignore those which do not have a pole at $`\delta _1=0`$. These can occur only in $`M(s_1,\lambda _1)`$, for the denominators do not yet vanish (note $`ϵ,\delta _20`$ at this stage). Since $$M(s_1,\lambda _1)=\underset{\stackrel{i<j}{s_1(i)>s_1(j)}}{}R(\lambda _{1_i}\lambda _{1_j}),$$ these permutations $`s_1`$ must interchange 1 and 2, which forces $$s_1\{(12),(13),(321)\}.$$ Next, when the residue at $`\delta _2=0`$ is taken, poles can occur in two different ways: $`M(s_2,\lambda _2)`$ might have a pole, or a denominator might vanish. The former occurs for $`s_2\{(12),(13),(321)\}`$ as well. ###### Proposition 7.7. Let $`s_1\{(12),(13),(321)\}`$. If one of the denominator terms $$(s_1\lambda _1+s_2\lambda _2)(\alpha _1^{})=0$$ or $$(s_1\lambda _1+s_2\lambda _2)(\alpha _2^{})=0,$$ then $$s_2\{e,(23),(123)\}=S_3\{(12),(13),(321)\},$$ i.e. $`M(s_2,\lambda _2)`$ has no pole at $`\delta _2=0`$. Proof: Now that $$\lambda _1=(\frac{1}{2}+it+ϵ,\frac{1}{2}+it+ϵ,2it2ϵ)$$ and $$\lambda _2=(\frac{1}{2}it,\frac{1}{2}it,2it),$$ the only possible way to get consecutive entries of $`s_1\lambda _1+s_2\lambda _2`$ to equal is if $$s_1=s_2(12).$$ $`\mathrm{}`$ We see that only the following terms have poles: $$\text{two from }M{}_{}{}^{}s:s_1,s_2\{(12),(13),(321)\}$$ or $$\text{one from }M,\text{ one from a denominator}:s_1=(12),s_2=e,s_1=(13),s_2=(123).$$ Note that $`(321)\times (23)`$ fails to have poles in the denominators. The 11 terms (see the chart in the appendix) are: $$(12)\times (12):\frac{e^{(1+ϵ,1+ϵ,2ϵ)(c,0,c)}}{(2)(1+3ϵ)}$$ $$(12)\times (13):\frac{e^{(\frac{1}{2}+3it+ϵ,ϵ,\frac{1}{2}3it2ϵ)(c,0,c)}}{(\frac{1}{2}+3it)(\frac{1}{2}+3it+3ϵ)}R(\frac{1}{2}3it)R(\frac{1}{2}3it)$$ $$(12)\times (321):\frac{e^{(1+ϵ,\frac{1}{2}+3it+ϵ,\frac{1}{2}3it2ϵ)(c,0,c)}}{(\frac{3}{2}3it)(6it+3ϵ)}R(\frac{1}{2}3it)$$ $$(13)\times (12):\frac{e^{(\frac{1}{2}3it2ϵ,ϵ,\frac{1}{2}+3it+ϵ)(c,0,c)}}{(\frac{1}{2}3it3ϵ)(\frac{1}{2}3it)}R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}+3it+3ϵ)$$ $$(13)\times (13):\frac{e^{(2ϵ,1+ϵ,1+ϵ)(c,0,c)}}{(13ϵ)(2)}R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}3it)R(\frac{1}{2}3it)$$ $$(13)\times (321):\frac{e^{(\frac{1}{2}3it2ϵ,\frac{1}{2}+3it+ϵ,1+ϵ)(c,0,c)}}{(6it3ϵ)(\frac{3}{2}+3it)}R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}3it)$$ $$(321)\times (12):\frac{e^{(1+ϵ,\frac{1}{2}3it2ϵ,\frac{1}{2}+3it+ϵ)(c,0,c)}}{(\frac{3}{2}+3it+3ϵ)(6it3ϵ)}R(\frac{1}{2}+3it+3ϵ)$$ $$(321)\times (13):\frac{e^{(\frac{1}{2}+3it+ϵ,\frac{1}{2}3it2ϵ,1+ϵ)(c,0,c)}}{(6it+3ϵ)(\frac{3}{2}3it3ϵ)}R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}3it)R(\frac{1}{2}3it)$$ $$(321)\times (321):\frac{e^{(1+ϵ,2ϵ,1+ϵ)(c,0,c)}}{(1+3ϵ)(13ϵ)}R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}3it).$$ The next two terms had limits taken in $`\delta _1,\delta _20`$ and, up to constants from the residues, are $$(12)\times e:\frac{e^{(ϵ,ϵ,2ϵ)(c,0,c)}}{3ϵ}$$ $$(13)\times (123):\frac{e^{(2ϵ,ϵ,ϵ)(c,0,c)}}{3ϵ}R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}+3it+3ϵ)R(\frac{1}{2}3it)R(\frac{1}{2}3it).$$ Note that $$R(\frac{1}{2}+x)R(\frac{1}{2}+x)=\frac{Z(\frac{1}{2}+x)}{Z(\frac{3}{2}+x)}=\frac{Z(\frac{3}{2}x)}{Z(\frac{3}{2}+x)}$$ by the functional equation. Accordingly, we may set $`ϵ=0`$ in the first 9 terms; the last two involve a limit at $`ϵ=0`$ producing derivatives. ###### Proposition 7.8. $$(i)R(\frac{1}{2}+it)1$$ (7.9) $$(ii)\left\|R(\frac{1}{2}+it)R(\frac{1}{2}+it)\right\|=1$$ (7.10) $$(iii)\frac{d}{dϵ}\frac{Z(\frac{3}{2}+it+ϵ)}{Z(\frac{3}{2}+itϵ)}_{ϵ=0}=O(\mathrm{log}t).$$ (7.11) Proof: (ii) follows from the facts that $`Z(\overline{s})=\overline{Z(s)}`$ and $`R(s)R(s)=1`$. For (i), recall $$R(s)=\sqrt{\pi }\frac{\mathrm{\Gamma }(\frac{s}{2})}{\mathrm{\Gamma }(\frac{s+1}{2})}\frac{\zeta (s)}{\zeta (s+1)}.$$ By Stirling’s formula $$\left\|\frac{\mathrm{\Gamma }(\frac{1}{4}+\frac{it}{2})}{\mathrm{\Gamma }(\frac{3}{4}+\frac{it}{2})}\right\|\sqrt{\frac{2}{t}}\text{ as }t\mathrm{},$$ and by the convexity bound, $$\zeta (\frac{1}{2}+it)=O_ϵ(t^{1/4+ϵ})$$ (7.12) (better bounds can be obtained – see \[Titchmarsh\]). Also, taking the logarithm of the Euler product of $`\zeta (s)`$ we find $$\mathrm{log}\zeta (\frac{3}{2}+it)=\underset{p\text{ prime}}{}\mathrm{log}(1p^{3/2it})=O(1),$$ which proves (i). Differentiating, $$\frac{\zeta ^{}}{\zeta }(\frac{3}{2}+it)=\underset{p}{}\frac{p^{3/2it}\mathrm{log}p}{1p^{3/2it}}=O(1)$$ also. So $$\frac{Z^{}}{Z}(\frac{3}{2}+it)=\frac{1}{2}\mathrm{log}\pi +\frac{1}{2}\frac{\mathrm{\Gamma }^{}}{\mathrm{\Gamma }}(\frac{3}{4}+\frac{it}{2})+\frac{\zeta ^{}}{\zeta }(\frac{3}{2}+it)=O(\mathrm{log}t),$$ proving (iii). $`\mathrm{}`$ Summarizing these pointwise bounds, ###### Proposition 7.9. For $`T`$ large, $$_{T1}^{T+1}_{}\|\mathrm{\Lambda }^CE(P_1,g,(it,it,2it),1)\|^2𝑑g𝑑t=O_ϵ(T^ϵ).$$ ### Assembling the Bounds Proof of (5.6): Choose $`c`$ large enough so that $`_C`$ is contained in $`\{x\widehat{\tau }_{P_1}(H_0(xC)),\widehat{\tau }_{P_2}(H_0(xC))0\}`$. Then the truncation does not affect $`_C`$ and $$__C\|E(g,\lambda )\|^2𝑑g_{}\|\mathrm{\Lambda }^CE(g,\lambda )\|^2𝑑g$$ for each Eisenstein series $`E(g,\lambda )`$ on $``$. In the propositions we bounded local integrals of all of the Eisenstein series as growing slower than any polynomial, with the exception of the maximal parabolic Eisenstein series (Remark 7.6). For this we must sum over the $`SL_2()\backslash `$ spectrum as well, so the contribution of the left-hand side of (7.2) near the point $`\lambda `$ is bounded by $`O_ϵ(\lambda ^{1+ϵ})`$. We note in comparison with Proposition 2.7 that (4.4) shows $$\widehat{g}_{t\mathrm{\Sigma }}(\lambda )=\{\begin{array}{cc}1+O_m(\text{dist}(\lambda ,(t\mathrm{\Sigma }))^m),\hfill & \lambda t\mathrm{\Sigma },\hfill \\ O_m(\text{dist}(\lambda ,(t\mathrm{\Sigma }))^m),\hfill & \lambda t\mathrm{\Sigma }\hfill \end{array},m0.$$ Since we have shown that all the Eisenstein series grow polynomially, we may switch the order of integration, ignore the tail, and conclude (see (5.5)) $$__CEis_{t\mathrm{\Sigma }}(x,x)𝑑x_{t\mathrm{\Sigma }}(1+\lambda )^{1+ϵ}𝑑\lambda .$$ Since $`\beta (\lambda )`$ grows at the rate of $`\lambda ^3`$ in almost all directions, we have completed the proof that $$__CEis_{g_{t\mathrm{\Sigma }}}(x,x)𝑑x=o\left(_{t\mathrm{\Sigma }}\beta (\lambda )𝑑\lambda \right).$$ $`\mathrm{}`$
warning/0006/cond-mat0006341.html	ar5iv	text	# Phase Transitions Driven by Vortices in 2D Superfluids and Superconductors: From Kosterlitz-Thouless to 1st Order ## Abstract The Landau-Ginzburg-Wilson hamiltonian is studied for different values of the parameter $`\lambda `$ which multiplies the quartic term (it turns out that this is equivalent to consider different values of the coherence length $`\xi `$ in units of the lattice spacing $`a`$). It is observed that amplitude fluctuations can change dramatically the nature of the phase transition: for small values of $`\lambda `$ ($`\xi /a>0.7`$) , instead of the smooth Kosterlitz-Thouless transition there is a first order transition with a discontinuous jump in the vortex density $`v`$ and a larger non-universal drop in the helicity modulus. In particular, for $`\lambda `$ sufficiently small ($`\xi /a1`$) , the density of bound pairs of vortex-antivortex below $`T_c`$ is so low that, $`v`$ drops to zero almost for all temperature $`T<Tc`$. Vortices play a central role in explaining the phase diagram and properties of superfluid systems both neutral and charged. The discovery of high-temperature superconductors has boosted the interest in the dynamics of the vortex lines in the mixed state and opened a new area which regards the physics of vortices as a new state of matter . Superfluid films and Josephson junction arrays in two dimensions are often described by XY-type models in which the unique microscopic degrees of freedom are phases. Thanks to the work of Berezinskii and Kosterlitz and Thouless in the early 70’s we have a fair understanding of the standard XY model in two dimensions. In their theory a thermodynamic phase transition, the Kosterlitz-Thouless (KT) transition is driven by the unbinding vortices ( singular phase configurations ) at a temperature $`T_{KT}`$. In refs. - it was pointed out that the KT phase transition might also apply to thin-film superconductors and some experimental evidence was discussed. In particular, the analysis of the 2 dimensional flux line lattice (FLL) melting transition of ref. is based on a KT-type theory. More recently, it was suggested that some features of the KT theory might be present in under-doped superconducting cuprates . However, the nature of the phase transition driven by vortices in 2D still remains under discussion. By means of Monte Carlo simulations of the XY model with a modified nearest neighbor interaction it was shown that, depending on the value of an additional parameter, continuous as well first-order transitions take place -. The existence of both kinds of phase transitions is in accordance with the richer structure of the 2D Coulomb gas found by Minnhagen and Wallin using self-consistent renormalization group equations and with the tendency towards first-order transition which develops in the case of a strong disorder coupling constant . Our approach to the nature of the phase transition driven by vortices in 2D superfluids/superconductors is a different one based on the effect of amplitude fluctuations. Let us start by noting that right at the ( expected vortex-pair unbinding ) transition the amplitude fluctuations cannot be considered as weak anymore and thus may affect the critical behavior . Therefore, instead of the XY model -with fixed amplitude fields- it is worth investigating the effects of vortices when amplitude fluctuations are not neglected. Hence, in this letter we analyze the more general Landau-Ginzburg-Wilson (LGW) lattice Hamiltonian, in terms of a complex field $`\psi =\|\psi \|\mathrm{exp}i\theta `$ ($`\|\psi \|`$ constant) and two parameters $`K`$ and $`\lambda `$ : $$\beta H=2K\underset{x}{}\underset{\mu }{}\|\psi _x\|\|\psi _{x+a\mu }\|cos(\theta _{x+a\mu }\theta _x)+\underset{x}{}[\lambda (\|\psi _x\|^21)^2+\|\psi _x\|^2],$$ (1) where $`\beta =1/T`$, $`a`$ is the lattice spacing, $`x`$ denotes the lattice sites and the index $`\mu =1,2`$. We will show that the nature of the phase transition of this model diverges dramatically from the KT when the parameter $`\lambda `$ is chosen sufficiently small -which in fact is equivalent to take the coherence length $`\xi a`$\- and that this is connected with the appearance of a sharp jump in the number of vortices. A straightforward discretization of the continuum Landau-Ginzburg Hamiltonian produces the expression : $$\beta H=\beta a^2\underset{x}{}[\underset{\mu =1}{\overset{2}{}}\frac{\mathrm{}^2}{2m}(\psi _{x+a\mu }^c\psi _x^c)^2/a^2+r\|\psi _x^c\|^2+u\|\psi _x^c\|^4],$$ (2) where the superscript $`c`$ in $`\psi `$ denotes the ordinary parameterization in the continuum theory, $`m`$ is the effective mass of the carriers and the coefficients $`r`$ and $`u`$ are analytic functions of the temperature, with $`u>0`$ for stability. Introducing a dimensionless order parameter: $`\overline{\psi }_x=\left(\frac{2u}{r}\right)^{\frac{1}{2}}\psi _x^c`$ and writing $$\beta H=\frac{1}{k_BT}a^2\underset{x}{}[\underset{\mu =1}{\overset{2}{}}(\overline{\psi }_{x+a\mu }\overline{\psi }_x)^2/a^2+\frac{1}{2\xi ^2}(1\|\overline{\psi }_x\|^2)^2],$$ (3) where $`\frac{1}{k_BT}=\frac{\mathrm{}^2\|r\|}{2mu}\beta `$ and $`\xi `$ is the coherence length given by $`\xi ^2=\frac{\mathrm{}^2}{2m\|r\|}`$. Parameterizations (1) and (3) are connected by the relations: $$\frac{\xi }{a}=\left(\frac{K}{\|12\lambda 4K\|}\right)^{\frac{1}{2}},T=\frac{\lambda (\xi /a)^2}{K^2}=\frac{\lambda }{K\|4K+2\lambda 1\|}.$$ (4) In the limit of $`\lambda =\mathrm{}`$ the radial degree of freedom is frozen and this model -sometimes said to describe soft spins with non fixed amplitude- becomes the XY model which is said to describe hard spins with fixed amplitude. A more interesting and less well studied limit is just the opposite i.e. small values of the $`\lambda `$ parameter. By (4) small values of $`\lambda `$ correspond to a large $`\xi `$ in units of $`a`$ i.e. $`\lambda `$ and $`1/\xi ^2`$ are the self-interaction coefficients respectively in (1) and (3) regulating amplitude fluctuations. We have simulated the Hamiltonian (1) using a Monte Carlo algorithm. The calculations were performed on square $`L\times L`$ lattices with periodic boundary conditions (PBC). In order to increase the speed of the simulation we have discretized the $`O(2)`$ global symmetry group to a $`Z(N)`$ and compared the results with previous runs carried out with the full O(2) group in relatively small lattices. For the case of $`Z(60)`$ we found no appreciable differences. Lattices with $`L=10,20,24,32,40`$ and (in some cases) 64 were used. For $`L=10,20,24`$ we thermalized with, usually, 20.000-40.000 sweeps and averaged over another 60.000-100.000 sweeps. For $`L=32,40`$ and 64 larger runs were performed, typically 50.000 sweeps were discarded for equilibration and averaged over 200.000 sweeps. We also performed some more extensive runs near $`K_c`$ for small values of $`\lambda `$. The following quantities were measured: i) The vortex density v. The standard procedure to calculate the vorticity on each plaquette is by considering the quantity $$m=\frac{1}{2\pi }([\theta _1\theta _2]_{2\pi }+[\theta _2\theta _3]_{2\pi }+[\theta _3\theta _4]_{2\pi }+[\theta _4\theta _1]_{2\pi }),$$ (5) where $`[\alpha ]_{2\pi }`$ stands for $`\alpha `$ modulo 2$`\pi `$: $`[\alpha ]_{2\pi }=\alpha +2\pi n`$, with $`n`$ an integer such that $`\alpha +2\pi n(\pi ,\pi ]`$, hence $`m=n_{12}+n_{23}+n_{34}+n_{41}`$. If $`m0`$, there exists a vortex which is assigned to the object dual to the given plaquette. Hence in the case d = 2, $`m`$, the dual of $`m`$, is assigned to the center of the original plaquette $`p`$. The vortex “charge” $`m`$ can take three values: 0, $`\pm 1`$ (the value $`\pm 2`$ has a negligible probability). $`v`$ defined as: $$v=\frac{1}{L^2}\underset{x}{}\|m_x\|,$$ (6) serves as a measure of the vortex density. ii) The energy density $`\epsilon =<H>/L^2`$ and the specific heat $`C_v`$, which were computed to measure the order of the phase transition. iii) The helicity modulus $`\mathrm{\Gamma }`$ which measures the phase-stiffness. For a spin system with PBC the helicity modulus measures the cost in free energy of imposing a “twist” equal to $`L\delta `$ in the phase between two opposite boundaries of the system. $`\mathrm{\Gamma }`$ is obtained in general as a second order derivative of the free energy with respect to $`\delta `$ -which can be regarded as a uniform statistical vector potential- evaluated for $`\delta 0`$. In such a way one gets the following expression , which generalizes the one introduced in ref. , to an order parameter with amplitude as well as phase variations: $$\mathrm{\Gamma }=\frac{1}{N}\{<\underset{<ij>}{}{}_{}{}^{}\psi _i\psi _j\mathrm{cos}(\theta _i\theta _j)>k<[\underset{<ij>}{}{}_{}{}^{}\psi _i\psi _j\mathrm{sin}(\theta _i\theta _j)]^2>\},$$ (7) where the primes denote that the sums are carried out over links along one of the 2 directions (x or y). Fig. 1-(a) shows a plot of $`v`$ vs. $`T`$ for different values of $`\lambda `$ and $`L=40`$. For $`\lambda =0.01`$ ($`\frac{\xi }{a}0.94`$ we observe a sharp jump in the vortex density $`v`$ (triangles up). As long as we increase $`\lambda `$ the jump becomes more smooth and moves to higher values of $`T_c`$ until for $`\lambda =10`$ (+ symbols) we get something very close to the $`\lambda \mathrm{}`$ KT behavior (circles). The increase in the density of vortices when amplitude fluctuations are large is in agreement with the analytical computations of ref. . What is new is the fact that, when $`\xi a`$ ($`\lambda 0.01`$), the transition occurs almost directly from 0 vortex to a plasma of vortices i.e. the bound pairs of vortex-antivortex seem to play no major role. Fig. 1-(b) is a zoom of 1-(a) showing the difference between $`\lambda =0.01`$ and $`\lambda =0.1`$ below $`T_c`$: for $`\lambda =0.1`$ $`v`$ drops to non-zero values due to the existence of vortex-antivortex pairs while for $`\lambda =0.01`$ $`v`$ drops much more sharply to 0 signaling the sudden extinction of vortices. The different behavior of $`v`$ above $`T_c`$ between the large fluctuating amplitude LGW-regime $`\frac{\xi }{a}1`$ and the KT-regime $`\frac{\xi }{a}1`$ is because amplitude fluctuations -governed by the $`\lambda `$ parameter- indeed decrease the energy of vortices enhancing vortex production. The same happens for the XY model with modified interaction ; in fact, the shape modification can be straightforwardly connected to a core energy variation. The scarcity of bound pairs of vortex-antivortex below Tc for the extreme fluctuations regime (very small values of $`\lambda `$ or $`\xi a`$) can be explained in terms of the behavior of their free energy $`F_{pair}=E_{pair}TS_{pair}`$ at $`TT_c`$ from below. Roughly, $`E_{pair}2E_c`$, where $`E_c`$ is the vortex core energy, and $`S_{pair}\mathrm{ln}(\frac{L^2}{\xi ^2})`$. The three quantities $`S_{pair},E_c`$ and $`T_c`$ all decrease as $`\lambda `$ decreases making difficult to disentangle the ”energetic” contribution to $`v`$, proportional to $`\mathrm{exp}[E_{pair}/T_c]`$, from the ”entropic” contribution $`1/\xi ^2`$. For the intermediate range of $`\lambda `$ (or $`\xi `$), although it is not easy to predict the behavior of the energetic factor, the entropy seems to be the main responsible for lowering the density of bound pairs . From a $`\lambda `$ sufficiently small, $`T_c`$ starts to decrease with $`\lambda `$ faster than $`E_{pair}`$ and thus $`\mathrm{exp}[F_{pair}/T]`$ decreases more and more sharply making smaller and smaller the probability of bound pairs of vortex-antivortex. Figures 2-4 show plots of $`\epsilon `$, $`v`$ and $`\mathrm{\Gamma }`$ for $`\lambda =0.01,0.1`$ and 10. For $`\lambda =0.01`$ the transition is clearly first-order: we observe latent heat and discontinuous changes in $`v`$ and $`\mathrm{\Gamma }`$ at $`T=T_c`$. On the other hand, for $`\lambda =10`$ (or $`\frac{\xi }{a}\frac{1}{\sqrt{42}}`$ ) the results are similar to those of the XY model. In particular, we get something close to the KT universal jump $`\mathrm{\Delta }\mathrm{\Gamma }=\frac{2}{\pi }`$. As long as $`\lambda `$ decreases we get a larger non universal jump in $`\mathrm{\Gamma }`$. $`\lambda =0.1`$ corresponds to something in between first order and KT. The double peak structure corresponding to the 2 coexisting phases, characteristic of a first-order transition, is showed in Fig. 5(a) for $`\lambda =0.01`$ and sizes $`L=10,20`$ and 40. The width $`w`$ of each of the peaks clearly scale as $`w\sqrt{(}\frac{1}{L^D})=\frac{1}{L}`$ due to ordinary non-critical fluctuations. Fig 5(b) is a zoom of the right peak centered around the higher energy phase. It shows that $`w0.1(L=40),w0.2(L=20)`$ and $`w0.4(L=10)`$. Therefore, one has a simple model in which the nature of the phase transition depends on the value of one parameter ($`\lambda `$ or $`\xi /a`$) which controls the thermal fluctuations of vortex cores and from which the following picture emerges: 1) For $`\lambda <0.1`$ or $`\frac{\xi }{a}1`$ the density of vortices experiments an abrupt jump which coincides with a first order transition with large latent heat and a non universal jump in $`\mathrm{\Gamma }`$ all at a $`T_c`$ which grows with $`\lambda `$. This is the LGW-regime as opposed to the much more smooth KT-regime. In particular, for sufficiently small values of $`\lambda `$ (for instance $`\lambda =0.01`$) below $`T_c`$ $`v`$ drops to zero i.e. the number of measured bound pairs of vortex-antivortex is negligible compared with the number found in the KT transition. 2) For $`\lambda 10`$ or $`(\frac{\xi }{a})^21`$ we get basically the XY model ($`\psi `$ is fixed = 1 for all K) and the more subtle KT transition (with an unobservable essential singularity in the specific heat at $`T_c`$ and a much more small non-universal maximum above $`T_c`$ and a universal jump in $`\mathrm{\Gamma }`$). The number of vortices and the energy evolve smoothly across the transition. 3) For intermediate values of $`\lambda `$ ($`0.1\lambda <10`$) we have an interpolating regime between LGW and KT. Whether or not a large enough increment of $`\frac{\xi }{a}`$ or $`\lambda `$ to alter the nature of the phase transition driven by vortices can be accomplished by varying some thermodynamic parameter, for instance the pressure, is something which deserves investigation. This change in the nature of the phase transition when amplitude fluctuations of $`\psi `$ are not negligible ($`\lambda 0.1`$) is in agreement with very recent variational computations for the same model . Furthermore, we checked that all the couples of values ($`K_c,\lambda `$) for what we found a first order transition are such that equation (4) gives $`0.7<\xi /a<1.1`$ in complete accordance with figure 6 of ref. from which we can see that for $`\xi /a`$ bigger than 0.5 the RG trajectories cross the first order line of Minnhagen’s generic phase diagram for the two-dimensional Coulomb gas . The jump in $`\mathrm{\Gamma }`$ larger than the universal value seems consistent with experiments on thin films HT<sub>c</sub> superconductors . Finally, our analysis could shed light on the nature of the melting transition for the 2D FLL which remains controversial theoretically as well as experimentally. Experimental works and numerical simulations favor a KT-like transition in some cases and a discontinuous one in others . A recent extensive Monte Carlo simulation found a first order transition at a temperature close to the estimated one assuming a KT melting transition . After all, it is possible that the nature of the melting transition in 2D strongly depends on the particular conditions and details of the studied specimen which in turn translate into different values of $`\frac{\xi }{a}`$. Work supported in part by CSIC, Project No. 052 and PEDECIBA. We are indebted with D. Ariosa, H. Beck, P. Curty, E. Dagotto and G. Gonzalez-Sprinberg for valuable discussions.
warning/0006/hep-th0006165.html	ar5iv	text	# 1 Introduction ## 1 Introduction The AdS/CFT correspondence - connects $`𝒩=4`$ supersymmetric $`SU(N)`$ Yang-Mills theory in four dimensions at large $`N`$ and strong ’t Hooft coupling $`\lambda =g_{YM}^2N`$ with type $`IIB`$ supergravity on the $`AdS_5\times S^5`$ background based on a perturbatively defined action. The correspondence works by comparison of series expansions in powers of $`\frac{1}{N^2}`$. At leading order many predictions have been verified, and at next order, results such as concerning anomalies, nonrenormalization theorems and $`\frac{1}{N^2}`$-corrections to field dimensions for composite fields and structure constants of the SYM<sub>4</sub> field algebra have been obtained -. In this context the evaluation of AdS graphs, that represent the holographic image of the AdS perturbation expansion in powers of $`\frac{\alpha ^{}}{R^2}=\lambda ^{\frac{1}{2}}`$, confronts us with serious technical problems whose difficulty goes much beyond the corresponding CFT flat space graphs. Partly with techniques developed first for CFT in flat space, the exchange graph was calculated and studied in a series of works . The results of all such calculations were finally expressed in terms of generalized hypergeometric functions. However, in some cases the field dimensions had to be specialized to small natural numbers. Due to these difficulties, we advocate another approach in this work. We present Green functions as multiple ”Mellin-Barnes integrals” <sup>4</sup><sup>4</sup>4Inverse Mellin transforms and Barnes integrals are equivalent over a meromorphic function $`\mathrm{\Phi }`$. This function $`\mathrm{\Phi }`$ is defined as the integral over a positive function on a compact domain. Usually one would expand this integral into a series of ratios of gamma functions, so that $`\mathrm{\Phi }`$ obtains poles from the gamma functions and from the divergence of the series. The latter are difficult to work out <sup>5</sup><sup>5</sup>5Except for the functions, say, $`{}_{2}{}^{}F_{1}^{}(1)`$ and $`{}_{3}{}^{}F_{2}^{}(1)`$ almost nothing is known. Thus we would like to extract the poles of $`\mathrm{\Phi }`$ by another method. The relevant poles of $`\mathrm{\Phi }`$, namely those to the right of the Mellin-Barnes contours, originate from the divergence (infinity) of the integrand at certain faces or intersections of faces of the regular polyhedral integration domain. So guessing them is not difficult. These poles form sequences which are integrally spaced and tend to $`+\mathrm{}`$. Of course at the end all Mellin-Barnes integration contours are shifted to $`+\mathrm{}`$, so that we find series expansions again. In Section 2 we discuss this method and typical results from the point of view of unitarity of Green functions and operator product expansions. Important information on the structure of the field algebra is obtained this way. Since AdS/CFT correspondence also implies (supposedly) a correspondence between both field algebras (all orders of $`\frac{1}{N^2}`$ included), the AdS conformal field theory as the holographic picture of supergravity and flat space CFT must therefore already show a partial correspondence on the level of the meromorphic functions $`\mathrm{\Phi }`$. We demonstrate that this is in fact true for the box graph. In Section 3 we study once again the exchange graph as a simple example of the previously developed method. In Section 4 we treat the box graph with arbitrary field dimensions <sup>6</sup><sup>6</sup>6This arbitrariness is essential with our method. We do not give all the details of the lengthy analysis. A few remarks are added in Section 5. ## 2 Critical exponents and unitarity We discuss here the connection between unitarity, operator product expansions and the ”critical exponents”, that we shall introduce now. Consider a four-point function in flat space CFT<sub>d</sub>. Its Green function $`G`$ can be split into a covariant multiplier and an invariant function $`\stackrel{~}{G}`$ $$G(x_1,x_2,x_3,x_4)=(x_{12}^2)^{\frac{1}{2}(\mu _1+\mu _2\mu _3+\mu _4)}(x_{13}^2)^{\frac{1}{2}(\mu _1\mu _2+\mu _3\mu _4)}(x_{23}^2)^{\frac{1}{2}(\mu _1\mu _2\mu _3+\mu _4)}(x_{34}^2)^{\mu _4}\stackrel{~}{G}(u,v)$$ (2.1) where $$x_{ij}=x_ix_j$$ (2.2) and $$u=\frac{x_{14}^2x_{23}^2}{x_{12}^2x_{34}^2},v=\frac{x_{13}^2x_{24}^2}{x_{12}^2x_{34}^2}$$ (2.3) are conformal invariant variables. If we intend an operator product expansion in the channel $$(1,\mathrm{\hspace{0.17em}4})(2,\mathrm{\hspace{0.17em}3})$$ we must let $$u0,v1.$$ (2.4) The function $`\stackrel{~}{G}(u,v)`$ can in turn be decomposed as $$\stackrel{~}{G}(u,v)=\underset{k=1}{\overset{K}{}}u^{\gamma _k}F_k(u,v),$$ (2.5) where $`F_k`$ are holomorphic functions in the neighborhood of (2.4) and possess the Taylor expansion $$F_k(u,v)=\underset{m,n=0}{\overset{\mathrm{}}{}}\frac{u^n(1v)^m}{n!m!}c_{mn}^{(k)}.$$ (2.6) The $`\gamma _k`$ are the ”critical exponents”. Of course, the $`\gamma _k`$ are, due to possible changes in the covariant multiplier (2.1), defined up to a common additive constant. So, what is the physical information encoded in these exponents? Consider the exchange of the scalar field of dimension $`\delta `$ in the channel <sup>7</sup><sup>7</sup>7We call the exchange channel the ”direct” channel $$(1,\mathrm{\hspace{0.17em}4})(2,\mathrm{\hspace{0.17em}3})$$ as described by Fig.2 where the dimension $`\delta `$ is assumed to be generic. Note, that the CFT<sub>d</sub> covariant vertex functions $$𝑑y\underset{i=1}{\overset{n}{}}((yx_i)^2)^{\mu _i}$$ (2.7) are necessarily ”unique”, i.e. they satisfy the condition $$\underset{i=1}{\overset{n}{}}\mu _i=d.$$ (2.8) A full vertex, such as in Fig.2, can always be resolved in three unique vertices (Fig.3) in an unambiguous fashion. This fact can be readily used to compute the Green function corresponding to Fig.2. The result is explicitly known and can be represented as $$\stackrel{~}{G}(u,v)=\underset{k=1}{\overset{2}{}}u^{\gamma _k}F_k(\delta ;u,v)$$ (2.9) with $$\gamma _1=\frac{1}{2}(\delta \mu _1\mu _4)$$ (2.10) $$\gamma _2=\frac{1}{2}(d\delta \mu _1\mu _4)$$ (2.11) and, after an appropriate renormalization, $$F_2(\delta ;u,v)=F_1(d\delta ;u,v).$$ (2.12) On the other hand, for the holographic image of the AdS exchange graph Fig.4, termed ”Witten graph”, we obtain $$\stackrel{~}{G}_W(u,v)=\underset{k=1}{\overset{3}{}}u^{\gamma _k}F_{W,k}(\delta ;u,v)$$ (2.13) with $`2\gamma _1+(\mu _1+\mu _4)`$ $`=\mu _1+\mu _4,`$ (2.14) $`2\gamma _2+(\mu _1+\mu _4)`$ $`=\mu _2+\mu _3,`$ (2.15) $`2\gamma _3+(\mu _1+\mu _4)`$ $`=\delta .`$ (2.16) Of course (2.10) and (2.16) are identical. On the other hand there are striking differences. In CFT jargon the $`k=2`$ term in (2.9) is called ”shadow term” of the $`k=1`$ term. Its appearance is a consequence of conformal harmonic analysis on $`𝐑_𝐝`$ and the equivalence of scalar representations with dimension $`\delta `$ and $`d\delta `$. Only if $$\delta \frac{d}{2}+1$$ (2.17) a scalar $`\underset{¯}{\text{field}}`$ of dimension $`d\delta `$ exists and we have two equivalent formulations of the same CFT<sub>d</sub>: each external leg of a Green function with dimension $`\delta `$ can by amputation be transformed into a leg with dimension $`d\delta `$ and vice versa. This shadow term is absent in (2.13). Instead, there are two terms $`k1,2`$, which are obviously connected with the exchange of some fields of dimension $$\mu _1+\mu _4+n(\mu _2+\mu _3+n),n𝐍_\mathrm{𝟎}.$$ Now we remember that unitarity of the S-matrix in perturbative quantum field theory is usually formulated by Cutkosky’s rule : cutting a graph (Fig.5) through internal lines and replacing these by a sum over the corresponding states (with appropriate normalization) gives an absorptive part of the original Green function. In CFT, we can reduce these states by operator product expansion to states created by conformal blocks of fields. In Fig.2 there is one conformal block, namely the conformal field of dimension $`\delta `$ and all its derivative fields. The same is true for Fig.4 and the part $`k=3`$, (eqn.(2.16)). The part $`k=1`$ (eqn.(2.14)) involves an infinite number of conformal tensor fields of rank $`l`$ with dimension $$\mu _1+\mu _4+l+2t$$ and their derivative fields. In fact, there are two parameters $`l`$ (rank) and $`t`$ (twist) to label all blocks exchanged. The same is true for $`k=2`$. The fact that for $`k=3`$ only one block is exchanged is reflected in the analytic property of $`F_{W,3}(u,v)`$. Thus we conclude that each critical exponent corresponds to an infinite tower of conformal blocks, that this tower is determined by a Cutkosky cut acting on internal and external lines and that $`2\gamma _k+\mu _1+\mu _4`$ is in fact the dimension of the lowest dimensional scalar field in the tower, which in turn can be understood as ”composite field” of the fields belonging to the lines cut. Thus the difference between CFT<sub>d</sub> and AdS<sub>d+1</sub> theory is in the exchange graphs: 1. there is no shadow term in AdS<sub>d+1</sub>; 2. there are terms from cutting external lines in AdS<sub>d+1</sub>. As was argued in the shadow term in CFT<sub>d</sub> and the external line terms in AdS<sub>d+1</sub> are necessary to guarantee analytic behavior in the crossed channel. Indeed, it turns out that such differences between CFT<sub>d</sub> and AdS<sub>d+1</sub> seem to arise $`\underset{¯}{\text{only}}`$ in the exchange graphs <sup>8</sup><sup>8</sup>8The usual notation is ”one-particle reducible graphs” in the direct channel. Next we consider a CFT<sub>d</sub> box graph with four unique vertices (Fig.6). The uniqueness conditions imply certain constraints on the dimensions of the external and internal fields, e.g. $$\mu _1+\mu _3=\mu _2+\mu _4=2d\underset{i=1}{\overset{4}{}}\lambda _i.$$ (2.18) This box graph Green function is explicitly known and $$\stackrel{~}{G}(u,v)=\underset{k=1}{\overset{3}{}}u^{\gamma _k}F_k(u,v)$$ (2.19) with $`\gamma _1`$ $`=0,`$ (2.20) $`\gamma _2`$ $`={\displaystyle \frac{1}{2}}(\lambda _1+\lambda _3\mu _1\mu _4),`$ (2.21) $`\gamma _3`$ $`={\displaystyle \frac{1}{2}}(\mu _2+\mu _3\mu _1\mu _4),`$ (2.22) Now we have Cutkosky cuts through the external pairs of lines as well $`(\gamma _1,\gamma _3)`$. We note that the box graph with non-unique vertices (full vertices) has not been calculated yet. Since the critical exponents determine the towers of conformal blocks that are coupled to the pairs of external and internal fields in the direct channel, they enter the structure of the field algebra. If the Maldacena conjecture in the strong version is correct, the large $`\lambda `$$`\frac{1}{N}`$-expanded SYM<sub>d</sub> with gauge group $`SU(N)`$ and $`𝒩=4`$ supercharges has the same field algebra as the holographic image of the AdS<sub>d+1</sub> supergravity with coupling constants of order $`\frac{1}{N^k}`$, $`k2𝐍`$. Therefore the results (2.19) - (2.22) should hold in the case of the Witten graph Fig.7 as well. We shall prove in the sequel, that this is correct indeed. ## 3 The singularity analysis of conformally covariant Green functions We aim at a direct determination of the critical exponents $`\gamma _k`$ (2.5) before attempting the explicit evaluation of integral representations. The Taylor coefficients $`c_{mn}^{(k)}`$ (2.6) are then finally represented as integrals which eventually can be evaluated numerically. Since analytic continuation of the integral representations in the parameters (field and space dimension) off the domain of absolute convergence is always tacitly understood, the integrals must necessarily be transformed into absolutely convergent expressions by substraction regularization methods before the numerics can be performed. The method of analyzing conformal Green functions developed by us consists of several steps: 1. We derive a multi-parametric Mellin-Barnes integral representation, where the integrand $`\mathrm{\Phi }`$ depends meromorphically on the Mellin-Barnes parameters and the field and space dimensions. This function $`\mathrm{\Phi }`$ is itself given as an integral of a positive function over a compact polyhedral domain $`𝕂_n`$ in $`𝐑_𝐧`$ with possible zeros and infinities on the boundary on $`𝕂_n`$. $`𝕂_n`$ is the $`n`$-dimensional generalization of the regular tetrahedron $`𝕂_3`$ or the regular triangle $`𝕂_2`$. $`𝕂_n`$ is bounded by $`(n+1)`$ faces $`𝕂_{n1}`$, which intersect in edges $`𝕂_{n2}`$ etc. 2. If the integrand is $`+\mathrm{}`$ on a face or a lower dimensional intersection $`𝕂_r`$, then poles may appear in the Mellin-Barnes parameters on the ”right” side of the Mellin-Barnes contours. 3. Two Mellin-Barnes parameters are connected with the kinematical variables $`u`$ and $`1v`$ (2.3) by the powers $$u^{\sigma _1}(v1)^{\sigma _2}.$$ (3.1) The pole positions of $`\mathrm{\Phi }`$ in $`\sigma _2`$ lie in $`𝐍_\mathrm{𝟎}`$ and the shift of the $`\sigma _2`$ integration contour to $`+\mathrm{}`$ gives simple power series in $`1v`$. The pole positions in $`\sigma _1`$ lie in different sequences $$\underset{𝑘}{}\{\gamma _k+𝐍_\mathrm{𝟎}\}$$ (3.2) which leads us to the series representations (2.5), (2.6). 4. Since the zero of an analytically continued integral is difficult to recognize (zeros can $`\underset{¯}{\text{only}}`$ arise after analytic continuation since the original integrand is a positive function) the list of candidates for the exponents $`\{\gamma _k\}`$ is generally too long. We can reduce this list by different arguments, e.g. a ”beta-function argument” and a symmetry argument. As a nontrivial example of describing our method, we choose the holographic image of the AdS<sub>d+1</sub> graph in Fig.4. Due to conformal invariance, a Green function can be completely reconstructed if three of its $`n3`$ variables are fixed to the values, say $$x_1=0,x_2=\mathrm{},x_3\text{arbitrary unit vector}$$ We shall exploit this fact by letting $`x_3\mathrm{}`$, but keeping translational and scale invariance $$\underset{x_3\mathrm{}}{\text{lim}}(x_3^2)^{\mu _3}G(x_1,x_2,x_3,x_4)=(x_{12}^2)^{\mathrm{\Delta }\mu }\stackrel{~}{G}(u,v)$$ (3.3) with $$u=\frac{x_{14}^2}{x_{12}^2},v=\frac{x_{24}^2}{x_{12}^2}$$ (3.4) and $$\mathrm{\Delta }\mu =\frac{1}{2}(\mu _1+\mu _2\mu _3+\mu _4).$$ (3.5) Denoting the bulk variables of $`\text{AdS}_{d+1}`$ by $`w_1,w_2,w_3,\mathrm{}`$ and boundary variables by $`\stackrel{}{x}`$, $`\stackrel{}{y}`$, $`\stackrel{}{z}`$, …we have for the bulk-to-boundary propagators $`(i\{1,2,3,4\})`$ $$K_{\mu _i}(w,x_i)=c_{\mu _i}\left(\frac{w_0}{w_0^2+(\stackrel{}{w}\stackrel{}{x})^2}\right)^{\mu _i}$$ (3.6) where $$c_{\mu _i}=\frac{\mathrm{\Gamma }(\mu _i)}{\pi ^{\frac{d}{2}}\mathrm{\Gamma }(\nu _i)},(\nu _i=\mu _i\frac{1}{2}d)$$ (3.7) and $$\underset{x_3\mathrm{}}{\text{lim}}(x_3^2)^{\mu _3}K_{\mu _3}(w,x_3)=c_{\mu _3}w_0^{\mu _3}$$ (3.8) For the bulk-to-bulk propagator we use the Mellin-Barnes integral representation $$G_\lambda (w,w^{})=\frac{1}{2\pi i}_i\mathrm{}^{+i\mathrm{}}𝑑s\mathrm{\Gamma }(s)e^{i\pi s}\frac{\mathrm{\Gamma }(\lambda +2s)}{\mathrm{\Gamma }(\stackrel{~}{\nu }+s+1)}\frac{1}{2\pi ^{\frac{d}{2}}}\left[\frac{w_0w_0^{}}{w_0^2+w_{}^{}{}_{0}{}^{2}+(\stackrel{}{w}\stackrel{}{w^{}})^2}\right]^{\lambda +2s}$$ (3.9) with $`\stackrel{~}{\nu }=\lambda \frac{1}{2}d`$. The graph of interest (Fig.4) is, up to coupling constants, factorials and symmetry factors, represented by the integral $$𝑑\mu (w)𝑑\mu (w^{})G_\lambda (w,w^{})\underset{i\{1,4\}}{}K_{\mu _i}(w,x_i)\underset{j\{2,3\}}{}K_{\mu _j}(w^{},x_j)$$ (3.10) where $`d\mu `$ is the invariant $`\text{AdS}_{d+1}`$ measure $$d\mu (w)=\frac{dw_0d\stackrel{}{w}}{w_0^{d+1}}$$ (3.11) The integration starts by using a $`\mathrm{\Gamma }`$-function auxiliary integration for each denominator in (3.6), (3.9) $$\frac{1}{x^{2\mu }}=\frac{1}{\mathrm{\Gamma }(\mu )}_0^{\mathrm{}}𝑑tt^{\mu 1}e^{tx^2}$$ (3.12) distributing the parameters $`\{t_i\}_{1,2,4}`$ to the $`K_{\mu _i}`$ and $`r`$ to $`G_\lambda `$. Then $`w_0`$ and $`w_0^{}`$ can be integrated giving $$\underset{i\{1,2\}}{}\frac{1}{2}\mathrm{\Gamma }(\frac{1}{2}\alpha _i)\eta _i^{\frac{1}{2}\alpha _i}$$ (3.13) with $$\eta _1=r+t_1+t_4,\eta _2=r+t_2$$ (3.14) and $$\alpha _1=\mu _1+\mu _4+\lambda +2sd,\alpha _2=\mu _2+\mu _3+\lambda +2sd$$ (3.15) The $`\stackrel{}{w},\stackrel{}{w}^{}`$ integration is Gaussian and gives $$(\frac{\pi ^2}{\text{det}A})^{\frac{d}{2}}\text{exp}\{\chi ^TA^1\chi D\}$$ (3.16) with $`A`$ $`=\left(\begin{array}{cc}\eta _1& r\\ r& \eta _2\end{array}\right),`$ (3.17) $`D`$ $`={\displaystyle \underset{i}{}}t_ix_i^2,`$ (3.18) $`\chi `$ $`=\left(\begin{array}{c}t_1x_1+t_4x_4\\ t_2x_2\end{array}\right)`$ (3.19) The exponent (3.16) can be written as a quadratic form $$\frac{1}{\text{det}A}\underset{i<j}{}\beta _{ij}(x_ix_j)^2$$ (3.20) where $$\beta _{12}=rt_1t_2,\beta _{14}=\eta _2t_1t_4,\beta _{24}=rt_2t_4$$ (3.21) Since we aim at an expression of the type (2.6), we can use (3.4) to write (3.20) as $$x_{12}^2\frac{\beta _0}{\text{det}A}\{1+\frac{\beta _{14}}{\beta _0}u+\frac{\beta _{24}}{\beta _0}(v1)\}$$ (3.22) with $$\beta _0=\beta _{12}+\beta _{24}=rt_2(t_1+t_4)$$ (3.23) Following Symanzik , the second and third term in (3.22) are represented by Mellin-Barnes integrals $$e^x=\frac{1}{2\pi i}_i\mathrm{}^{+i\mathrm{}}𝑑\sigma \mathrm{\Gamma }(\sigma )x^\sigma $$ (3.24) Finally we perform one integration by introducing scaled parameters $$T=r+\underset{i\{1,2,4\}}{}t_i,t_i=T\tau _i,i\{1,2,4\},r=T\rho ,$$ (3.25) so that $$\text{det}A=T^2[\rho (1\rho )+\tau _2(1\rho \tau _2)].$$ (3.26) The remaining parameter integrals can then be summed up into a meromorphic function. $`\mathrm{\Phi }(\sigma _1,\sigma _2,s)`$ $`=\mathrm{\Gamma }(\sigma _1)\mathrm{\Gamma }(\sigma _2)\mathrm{\Gamma }(s){\displaystyle _{𝕂_2}}𝑑\tau _1𝑑\tau _2𝑑\tau _4𝑑\rho \delta (1\tau _1\tau _2\tau _4\rho )\tau _1^{\mu _11}\tau _2^{\mu _21}\tau _4^{\mu _41}\rho ^{\lambda +2s1}`$ $`\times (1\tau _2)^{\frac{1}{2}\alpha _1}(\rho +\tau _2)^{\frac{1}{2}\alpha _2}[\rho \tau _2(\tau _1+\tau _4)]^{\mathrm{\Delta }\mu \sigma _2}[\tau _1\tau _4(\rho +\tau _2)]^{\sigma _1}({\displaystyle \frac{\tau _4}{\tau _1+\tau _4}})^{\sigma _2}`$ $`\times [\rho (1\rho )+\tau _2(1\rho \tau _2)]^{\frac{1}{2}d+\mathrm{\Delta }\mu }`$ (3.27) and this enters a threefold Mellin-Barnes integral $`\stackrel{~}{G}(u,v)`$ $`={\displaystyle \frac{1}{8\pi ^{\frac{3}{2}d}}}{\displaystyle \frac{\mathrm{\Gamma }(\mu _3)}{_{i=1}^4\mathrm{\Gamma }(\nu _i)}}(2\pi i)^3{\displaystyle \underset{i\mathrm{}}{\overset{+i\mathrm{}}{}}}𝑑\sigma _1𝑑\sigma _2𝑑s{\displaystyle \frac{\mathrm{\Gamma }(\frac{1}{2}\alpha _1)\mathrm{\Gamma }(\frac{1}{2}\alpha _2)\mathrm{\Gamma }(\mathrm{\Delta }\mu +\sigma _2+\sigma _2)}{\mathrm{\Gamma }(\stackrel{~}{\nu }+s+1)}}`$ $`\times e^{i\pi s}u^{\sigma _1}(v1)^{\sigma _2}\mathrm{\Phi }(\sigma _1,\sigma _2,s)`$ (3.28) with $`\mathrm{\Delta }\mu =\frac{1}{2}(\mu _1+\mu _2\mu _3+\mu _4)`$, see (3.5). Such a representation (3), (3) of any four-point function for $`\text{CFT}_d`$ or $`\text{AdS}_{d+1}`$ field theory is the starting point for our singularity analysis, leading to the critical exponents. In this particular case we can simplify the integral representation (3) by integrating over $`\xi `$ in $$\tau _1=\tau \xi ,\tau _4=\tau (1\xi )$$ (3.29) $`\mathrm{\Phi }(\sigma _1,\sigma _2,s)`$ $`={\displaystyle \frac{\mathrm{\Gamma }(\mu _1+\sigma _1)\mathrm{\Gamma }(\mu _4+\sigma _1+\sigma _2)}{\mathrm{\Gamma }(\mu _1+\mu _4+2\sigma _1+\sigma _2)}}\mathrm{\Gamma }(\sigma _1)\mathrm{\Gamma }(\sigma _2)\mathrm{\Gamma }(s){\displaystyle _{𝕂_2}}𝑑\tau 𝑑\tau _2𝑑\rho \delta (1\tau \tau _2\rho )`$ $`\times \tau ^{\mu _1+\mu _4\mathrm{\Delta }\mu +\sigma _11}(1\tau )^{\frac{1}{2}\alpha _2+\sigma _1}\tau _2^{\mu _2\mathrm{\Delta }\mu \sigma _11}(1\tau _2)^{\frac{1}{2}\alpha _1}`$ $`\times \rho ^{\lambda +2s\mathrm{\Delta }\mu \sigma _11}[\rho (1\rho )+\tau _2(1\rho \tau _2)]^{\frac{1}{2}d+\mathrm{\Delta }\mu }`$ (3.30) Here $`\sigma _2`$ has vanished from the integral into the factor in front. Except for the factor $`\mathrm{\Gamma }(\sigma _2)`$, there is no pole to the right of the $`\sigma _2`$ Mellin-Barnes contour. This a general feature since (see(3.23)) in $`𝕂_n`$ $`0{\displaystyle \frac{\beta _{24}}{\beta _0}}={\displaystyle \frac{\beta _{24}}{\beta _{12}+\beta _{24}}}\mathrm{\hspace{0.25em}1}`$ (3.31) There are obviously poles from the faces $`\tau _2=0`$ and $`\rho =0`$ in $`\sigma _1`$, arising from the Mittag-Leffler expansion $$t^{\mu 1}\mathrm{\Theta }(t)\underset{\text{poles only}}{}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n\delta ^{(n)}(t)}{n!(\mu +n)}$$ (3.32) with positions $`\mu 𝐍_\mathrm{𝟎}`$. Including the poles in $`\sigma _1`$ from the factor $`\mathrm{\Gamma }(\sigma _1)`$, we have three possibilities: $`(n𝐍_\mathrm{𝟎})`$ 1. $$\sigma _1=n$$ (3.33) 2. $$\text{from}\tau _2=0:\sigma _1=\mu _2\mathrm{\Delta }\mu +n$$ (3.34) 3. $$\text{from}\rho =0:\sigma _1=\lambda +2s\mathrm{\Delta }\mu +n$$ (3.35) In the cases (1.) and (2.) we get the critical exponents $`\gamma _1`$ $`=0,`$ (3.36) $`\gamma _2`$ $`={\displaystyle \frac{1}{2}}(\mu _2+\mu _3\mu _1\mu _4)`$ (3.37) whereas case (3.) necessitates knowledge of the pole positions in $`s`$. One possibility is that these poles are produced by $`\mathrm{\Gamma }(s)`$, then $`\sigma _1`$ $`=\lambda \mathrm{\Delta }\mu +n,`$ (3.38) $`\gamma _3`$ $`=\lambda \mathrm{\Delta }\mu `$ (3.39) There is another candidate for poles in $`\sigma _1`$, namely the intersection of the faces (2) and (3): $$\tau _2=\rho =0,\tau =1.$$ (3.40) We use the parameters $$\rho =\omega \psi ,\tau _2=\omega (1\psi ),\tau =1\omega $$ (3.41) The behavior of the integrand at $`w0`$ is given by $$_0𝑑\omega \omega ^{(\mu _2\mathrm{\Delta }\mu \sigma _1)+(\lambda +2s\mathrm{\Delta }\mu \sigma _1)+(\frac{1}{2}\alpha _2+\sigma _1)+(\frac{1}{2}d+\mathrm{\Delta }\mu )1}$$ (3.42) The exponent is $$\frac{1}{2}\lambda +s\frac{1}{2}(\mu _1+\mu _4)\sigma _11$$ (3.43) and gives rise to poles in $`\sigma _1`$ at 1. $$\tau _2=\rho =0:\sigma _1=\frac{1}{2}\lambda +s\frac{1}{2}(\mu _1+\mu _4)+n$$ (3.44) If the $`s`$ poles are from $`\mathrm{\Gamma }(s)`$, we get from (3.44) $`\sigma _1`$ $`={\displaystyle \frac{1}{2}}(\lambda \mu _1\mu _4)+n,`$ (3.45) $`\gamma _4`$ $`={\displaystyle \frac{1}{2}}(\lambda \mu _1\mu _4)`$ (3.46) But there exist other $`s`$-poles. If we consider (3.44), set $`n=0`$ and insert the delta function following from (3.32) into (3), there remains the $`\psi `$-integral, see (3.41) $$\begin{array}{c}_0^1𝑑\psi \psi ^{\lambda +2s\mathrm{\Delta }\mu \sigma _11}(1\psi )^{\mu _2\mathrm{\Delta }\mu \sigma _11}\|_{\sigma _1=\frac{1}{2}\lambda +s\frac{1}{2}(\mu _1+\mu _4)}\hfill \\ \hfill =\frac{1}{\mathrm{\Gamma }(\mu _3)}\mathrm{\Gamma }(\frac{1}{2}(\mu _2+\mu _3\lambda )s)\mathrm{\Gamma }(\frac{1}{2}(\lambda +\mu _3\mu _2)+s)\end{array}$$ (3.47) which shows, that there exist relevant $`s`$-poles from the first factor in the numerator. For arbitrary $`n`$ in (3.44) the poles lie at $$s+n=\frac{1}{2}(\mu _2+\mu _3\lambda )+n^{},n^{}𝐍_\mathrm{𝟎}$$ (3.48) If we consider (3.35) at $`n=0`$ and insert it together with the delta function (3.32) into (3), then the integral turns into a beta-function $$\begin{array}{c}_0^1𝑑\tau _2\tau _2^{(\mu _2\mathrm{\Delta }\mu \sigma _11)+(\frac{1}{2}\alpha _2+\sigma _1)+(\frac{1}{2}d+\mathrm{\Delta }\mu )}(1\tau _2)^{(\mu _1+\mu _4\mathrm{\Delta }\mu +\sigma _11)\frac{1}{2}\alpha _1+(\frac{1}{2}d+\mathrm{\Delta }\mu )}\hfill \\ \hfill =\frac{1}{\mathrm{\Gamma }(0)}\mathrm{\Gamma }(\frac{1}{2}(\mu _2\mu _3\lambda )s)\mathrm{\Gamma }(\frac{1}{2}(\lambda \mu _2+\mu _3)+s)\end{array}$$ (3.49) The denominator is unchanged if we let $`n`$ in (3.35) assume arbitrary values from $`𝐍_\mathrm{𝟎}`$. Thus the denominator of the beta-function lets the singularity (3.35) vanish, implying that (3.38), (3.39) do not exist either. Only in exceptional cases do we get control over the zeros when we can perform an integral completely. Often the integral is a beta-function, then we call our way of proof ”the beta-function argument”. More effort is needed to evaluate integrals in terms of functions $`{}_{p+1}{}^{}F_{p}^{}(1)`$ in which case the zeros are also controllable. A simple but surprisingly powerful argument to eliminate whole sequences of poles comes from the symmetry of the graph (Fig.4). We define this symmetry to consist of those mappings of the graph on itself: 1. which lead to the same graph after an appropriate relabelling of the external coordinates and the field dimensions; 2. leave $`u`$ and $`v`$ invariant. In the case of Fig.4, this leads to a group $`Z_2\times Z_2`$, generated by the reflections $`S_1`$ $`:12,\mathrm{\hspace{0.25em}3}4`$ $`S_2`$ $`:14,\mathrm{\hspace{0.25em}3}2`$ (3.50) While the Green function $`G(x_1,x_2,x_3,x_4)`$ is invariant under $`Z_2\times Z_2`$ by definition, the invariant function $`\stackrel{~}{G}`$ is not. Let $`S_i(\stackrel{~}{G})`$ denote the function obtained by applying (3) to the dimensions in $`\stackrel{~}{G}`$, then from (2.1) we obtain $$\stackrel{~}{G}(u,v)=u^{\delta _i}v^{ϵ_i}S_i(\stackrel{~}{G})(u,v)$$ (3.51) with $`S_1:\delta _1`$ $`={\displaystyle \frac{1}{2}}(\mu _2+\mu _3\mu _1\mu _4),ϵ_1`$ $`={\displaystyle \frac{1}{2}}(\mu _1+\mu _3\mu _2\mu _4)`$ (3.52) $`S_2:\delta _2`$ $`=0,ϵ_2`$ $`=ϵ_1`$ (3.53) Inserting (3.51) into (2.5), we see that the labels $`\{k\}`$ of $`\gamma _k`$ are submitted to a representation of $`Z_2\times Z_2`$: $`S_i\sigma _i`$, so that: $`S_i(\gamma _k)+\delta _i`$ $`=\gamma _{\sigma _i(k)}`$ (3.54) $`S_i(G_k)v^{ϵ_i}`$ $`=G_{\sigma _i(k)}`$ (3.55) Holomorphy of $`G_k`$ at $`v=1`$ is obviously not touched by (3.55). Applying (3.54) to the graph Fig.4, we find $`\sigma _1(1)`$ $`=2;\sigma _1(2)`$ $`=1,`$ $`\sigma _2(1)`$ $`=1;\sigma _2(2)`$ $`=2`$ (3.56) and $$\sigma _{1,2}(4)=4$$ (3.57) whereas $`\gamma _3`$ does not fit into any representation. ## 4 The AdS box graph Now we turn to the box graph Fig.7. In terms of bulk-to-bulk propagators $`G_\lambda `$ and bulk-to-surface propagators $`K_\mu `$, the Green function is given by the integral $$G(x_1,x_2,x_3,x_4)=\underset{i=1}{\overset{4}{}}d\mu (w_i)K_{\mu _i}(x_i,w_i)G_{\lambda _i}(w_i,w_{i+1}),w_5=w_1$$ (4.1) Again we consider the limit (3.3), (3.8). Due to the four bulk-to-bulk propagators, the invariant Green function $`\stackrel{~}{G}(u,v)`$ has the form of a sixfold Mellin-Barnes integral<sup>9</sup><sup>9</sup>9All the techniques and notations used are the same as in the preceeding section. $`\stackrel{~}{G}(u,v)`$ $`={\displaystyle \frac{1}{2^8\pi ^{2d}}}{\displaystyle \frac{\mathrm{\Gamma }(\mu _3)}{_{i=1}^4\mathrm{\Gamma }(\nu _i)}}(2\pi i)^6{\displaystyle \underset{i\mathrm{}}{\overset{+i\mathrm{}}{}}}𝑑\sigma _1𝑑\sigma _2{\displaystyle \underset{i\mathrm{}}{\overset{+i\mathrm{}}{}}}\{{\displaystyle \underset{i=1}{\overset{4}{}}}ds_i{\displaystyle \frac{\mathrm{\Gamma }(\frac{1}{2}\alpha _i)}{\mathrm{\Gamma }(\stackrel{~}{\nu }_i+s_i+1)}}\}`$ $`\times \mathrm{\Gamma }(\mathrm{\Delta }\mu +\sigma _1+\sigma _2)e^{i\pi _is_i}u^{\sigma _1}(v1)^{\sigma _2}\mathrm{\Phi }(\sigma _1,\sigma _2,s_1,s_2,s_3,s_4)`$ (4.2) where $$\alpha _i=\mu _i+\lambda _i+2s_i+\lambda _{i1}+2s_{i1}d,(\lambda _0=\lambda _4,s_0=s_4)$$ (4.3) and the meromorphic function $`\mathrm{\Phi }`$ is given by $$\begin{array}{c}\mathrm{\Phi }(\sigma _1,\sigma _2,s_1,s_2,s_3,s_4)=\underset{i=1}{\overset{2}{}}\mathrm{\Gamma }(\sigma _i)\underset{j=1}{\overset{4}{}}\mathrm{\Gamma }(s_j)_{𝕂_6}(\underset{i=1(3)}{\overset{4}{}}d\tau _i\tau _i^{\mu _i1})(\underset{j=1}{\overset{4}{}}d\rho _j\rho _j^{\lambda _j+2s_j1}ϵ_j^{\frac{1}{2}\alpha _j})\hfill \\ \hfill \times \delta (1\underset{i=1(3)}{\overset{4}{}}\tau _i\underset{j=1}{\overset{4}{}}\rho _j)\delta (A)^{\frac{d}{2}+\mathrm{\Delta }\mu }f_0^{\mathrm{\Delta }\mu \sigma _1\sigma _2}f_1^{\sigma _1}f_2^{\sigma _2}\end{array}$$ (4.4) Here $$ϵ_i=\tau _i+\rho _i+\rho _{i1},(\tau _3=0,\rho _0=\rho _4)$$ (4.5) and the remaining functions $`f_0,f_1,f_2`$ and $`\delta (A)`$ can be represented best with the help of elementary symmetric polynomials $`S_2(1,2,3)`$ $`=\rho _1\rho _2+\rho _1\rho _3+\rho _2\rho _3`$ $`S_3(1,2,3,4)`$ $`=\rho _1\rho _2\rho _3+\rho _1\rho _3\rho _4+\rho _2\rho _3\rho _4+\rho _1\rho _2\rho _4`$ (4.6) namely $`f_0`$ $`=\tau _2[\tau _1\tau _4S_2(1,2,3)+(\tau _1+\tau _4)S_3(1,2,3,4)]`$ (4.7) $`f_1`$ $`=\tau _1\tau _4[S_3(1,2,3,4)+\tau _2\rho _4(\rho _2+\rho _3)]`$ (4.8) $`f_2`$ $`=\tau _2\tau _4[S_3(1,2,3,4)+\tau _1\rho _2\rho _3]`$ (4.9) $`\delta (A)`$ $`=\tau _1\tau _2\tau _4(\rho _2+\rho _3)+\tau _1\tau _2S_2(2,3,4)+\tau _1\tau _4S_2(1,2,3)`$ $`+\tau _2\tau _4(\rho _1+\rho _4)(\rho _2+\rho _3)+(\tau _1+\tau _2+\tau _4)S_3(1,2,3,4)`$ (4.10) The function $`\delta (A)`$ originates from the determinant in the Gaussian integration. It is obvious that $$0\frac{f_2}{f_0}\mathrm{\hspace{0.25em}1}\text{on}𝕂_6$$ (4.11) so that the only relevant poles in $`\sigma _2`$ arise from $`\mathrm{\Gamma }(\sigma _2)`$. The analysis of the pole positions in $`\sigma _1`$ is rather involved. In the sequel $`n_0𝐍_\mathrm{𝟎}`$ holds throughout. There is one face of type $`𝕂_5`$ producing a singularity: 1. $`\tau _2=0`$: poles appear at $`\sigma _1`$ $`={\displaystyle \frac{1}{2}}(\mu _2+\mu _3\mu _1\mu _4)+n_0`$ (4.12) $`\gamma _1`$ $`={\displaystyle \frac{1}{2}}(\mu _2+\mu _3\mu _1\mu _4)`$ (4.13) There are two faces of $`𝕂_4`$ type leading to poles. 1. $`\rho _1=\rho _2=0`$: we introduce the parameters $$\rho _1=\rho \xi ,\rho _2=\rho (1\xi )$$ (4.14) and let $`\rho 0`$. This gives pole positions $$\sigma _1=\lambda _1+\lambda _2+2(s_1+s_2)\mathrm{\Delta }\mu +n_0$$ (4.15) If the pole positions of $`s_1,s_2`$ are chosen from $`𝐍_\mathrm{𝟎}`$, we get $$\gamma _2=\lambda _1+\lambda _2\mathrm{\Delta }\mu $$ (4.16) The other case is 1. $`\rho _1=\rho _3=0`$: this case is treated analogously to case (II). We find poles at $$\sigma _1=\lambda _1+\lambda _3+2(s_1+s_3)\mathrm{\Delta }\mu +n_0$$ (4.17) If the pole positions of $`s_1,s_3`$ are from $`𝐍_\mathrm{𝟎}`$, we find $$\gamma _3=\lambda _1+\lambda _3\mathrm{\Delta }\mu $$ (4.18) Now we come to the intersections $`𝕂_5𝕂_4`$ and $`𝕂_4𝕂_4^{}`$ of $`𝕂_3`$ type. 1. $`\rho _1=\rho _2=\rho _3=0`$: we choose as parameters $$\rho _i=\rho \xi _i,i\{1,2,3\},\underset{i}{}\xi _i=1$$ (4.19) and let $`\rho 0`$. Pole positions are $$\sigma _1=\lambda _1+\frac{1}{2}(\lambda _2+\lambda _3)+2s_1+s_2+s_3\frac{1}{2}(\mu _1+\mu _2+\mu _4)+n_0$$ (4.20) If the $`\{s_i\}_{i=1}^3`$ have poles in $`𝐍_\mathrm{𝟎}`$, we get $$\gamma _4=\lambda _1+\frac{1}{2}(\lambda _2+\lambda _3)\frac{1}{2}(\mu _1+\mu _2+\mu _4)$$ (4.21) 2. $`\rho _1=\rho _2=\tau _2=0`$: we choose as parameters $$\rho _i=\rho \xi _i,i\{1,2\},\tau _2=\rho \xi _3$$ (4.22) and let $`\rho 0`$. The poles of $`\sigma _1`$ appear at $$\sigma _1=\frac{1}{2}(\lambda _1+\lambda _2\mu _1\mu _4+\mu _3)+s_1+s_2+n_0$$ (4.23) Provided the poles of $`s_1,s_2`$ are in $`𝐍_\mathrm{𝟎}`$, we find $$\gamma _5=\frac{1}{2}(\lambda _1+\lambda _2\mu _1\mu _4+\mu _3)$$ (4.24) However, if we perform some of the integrations after insertion of the delta function $`\delta ^{(0)}(\rho )`$ corresponding to the pole (4.29) by (3.32), we obtain a beta function with denominator $`\mathrm{\Gamma }(2n_0)`$. So these poles (V) cancel completely. 3. $`\rho _1=\rho _3=\tau _2=0`$: We proceed as in the case (V) and get as pole positions $$\sigma _1=\lambda _1+\lambda _3+\frac{1}{2}(\mu _2+\mu _3\mu _1\mu _4)\frac{1}{2}d+2(s_1+s_3)+n_0$$ (4.25) which, if the poles of $`s_1,s_3`$ are in $`𝐍_\mathrm{𝟎}`$, gives $$\gamma _6=\lambda _1+\lambda _3+\frac{1}{2}(\mu _2+\mu _3\mu _1\mu _4)\frac{1}{2}d$$ (4.26) 4. Finally, there is one $`𝕂_2`$ face: $`\tau _2=\rho _1=\rho _2=\rho _3=0`$: coordinates are $$\rho _i=\rho \xi _i,i\{1,2,3\},\tau _2=\rho \chi _4,\underset{i}{}\xi _i=1$$ (4.27) and we let $`\rho 0`$. We get the pole positions $$\sigma _1=\frac{1}{2}(\lambda _1+\lambda _3\mu _1\mu _4)+s_1+s_3+n_0$$ (4.28) If $`s_1,s_3`$ have poles in $`𝐍_\mathrm{𝟎}`$, we obtain $$\gamma _7=\frac{1}{2}(\lambda _1+\lambda _3\mu _1\mu _4).$$ (4.29) The symmetry group of the graph Fig.7 is the same as that of Fig.4: $`Z_2\times Z_2`$. It acts on the $`\lambda _i`$ as $`S_1(\lambda _i)`$ $`=\lambda _i,i1,3,S_1(\lambda _2)`$ $`=\lambda _4,S_1(\lambda _4)`$ $`=\lambda _2`$ (4.30) $`S_2(\lambda _i)`$ $`=\lambda _i,i2,4,S_2(\lambda _1)`$ $`=\lambda _3,S_2(\lambda _3)`$ $`=\lambda _1`$ (4.31) This rules out all $`\gamma `$’s, except $`\gamma _1,\gamma _7`$ and of course $`\gamma _0=0`$, which originates from the $`\sigma _1`$ poles of $`\mathrm{\Gamma }(\sigma _1)`$. Thus the AdS box graph has the same critical exponents as the CFT box graph Fig.6. The poles in $`s_4`$ are all from $`\mathrm{\Gamma }(s_4)`$. The other variables $`(s_1,s_2,s_3)`$produce poles of the function $`\mathrm{\Phi }`$ (4.4) that can be ordered in (triple) sequences $$\{(\nu _1+n_1,\nu _2+n_2,\nu _3+n_3),\nu _i\text{fixed},n_i𝐍_\mathrm{𝟎}\text{running}\}$$ (4.32) In the two tables below we list all possible triples $`(\nu _1,\nu _2,\nu _3)`$ and their connection (”origin”) with the $`\sigma _1`$ singularities (I) - (VII). The entries in the tables originating from the cases VI or VII are marked by (\*). The corresponding $`n_i`$ runs over $`\frac{1}{2}𝐍_\mathrm{𝟎}`$ (not $`𝐍_\mathrm{𝟎}`$). ## 5 Concluding Remarks We have proved that for the box graphs of CFT<sub>d</sub> and AdS<sub>d+1</sub> supergravity, we obtain the same critical exponents, namely those which are determined from the ”Cutkosky rule” with external lines included. We suggest that this behavior is also shown by other one-particle-irreducible graphs. Each critical exponent $`\gamma _k`$ belongs to one or more sequence of poles in the Mellin-Barnes parameters $`(s_1,s_2,s_3)`$, each of which is generated by a triple $`(\nu _1,\nu _2,\nu _3)`$ (see (4.32)), and each sequence contributes to the coefficient $`c_{mn}^{(k)}`$ in (2.6). The larger the number of $`\gamma _k,\nu _1,\nu _2,\nu _3`$ that are nonzero, the smaller the number of remaining integrations. More details on this can be found in . ## Acknowledgements The authors thank A. C. Petkou for interesting discussions during the initial stages and one of us (W. R.) thanks the staff of the Werner-Heisenberg-Institut in Munich for their hospitality during the final stage of this work.
warning/0006/hep-ex0006019.html	ar5iv	text	# 1 Summary of the selected event sample for the two tagging modes and the three cross-check modes. ## <sup>∗∗</sup>List of Authors Koya Abe,<sup>(24)</sup> Kenji Abe,<sup>(15)</sup> T. Abe,<sup>(21)</sup> I. Adam,<sup>(21)</sup> H. Akimoto,<sup>(21)</sup> D. Aston,<sup>(21)</sup> K.G. Baird,<sup>(11)</sup> C. Baltay,<sup>(30)</sup> H.R. Band,<sup>(29)</sup> T.L. Barklow,<sup>(21)</sup> J.M. Bauer,<sup>(12)</sup> G. Bellodi,<sup>(17)</sup> R. Berger,<sup>(21)</sup> G. Blaylock,<sup>(11)</sup> J.R. Bogart,<sup>(21)</sup> G.R. Bower,<sup>(21)</sup> J.E. Brau,<sup>(16)</sup> M. Breidenbach,<sup>(21)</sup> W.M. Bugg,<sup>(23)</sup> D. Burke,<sup>(21)</sup> T.H. Burnett,<sup>(28)</sup> P.N. Burrows,<sup>(17)</sup> A. Calcaterra,<sup>(8)</sup> R. Cassell,<sup>(21)</sup> A. Chou,<sup>(21)</sup> H.O. Cohn,<sup>(23)</sup> J.A. Coller,<sup>(4)</sup> M.R. Convery,<sup>(21)</sup> V. Cook,<sup>(28)</sup> R.F. Cowan,<sup>(13)</sup> G. Crawford,<sup>(21)</sup> C.J.S. Damerell,<sup>(19)</sup> M. Daoudi,<sup>(21)</sup> S. Dasu,<sup>(29)</sup> N. de Groot,<sup>(2)</sup> R. de Sangro,<sup>(8)</sup> D.N. Dong,<sup>(13)</sup> M. Doser,<sup>(21)</sup> R. Dubois, I. Erofeeva,<sup>(14)</sup> V. Eschenburg,<sup>(12)</sup> E. Etzion,<sup>(29)</sup> S. Fahey,<sup>(5)</sup> D. Falciai,<sup>(8)</sup> J.P. Fernandez,<sup>(26)</sup> K. Flood,<sup>(11)</sup> R. Frey,<sup>(16)</sup> E.L. Hart,<sup>(23)</sup> K. Hasuko,<sup>(24)</sup> S.S. Hertzbach,<sup>(11)</sup> M.E. Huffer,<sup>(21)</sup> X. Huynh,<sup>(21)</sup> M. Iwasaki,<sup>(16)</sup> D.J. Jackson,<sup>(19)</sup> P. Jacques,<sup>(20)</sup> J.A. Jaros,<sup>(21)</sup> Z.Y. Jiang,<sup>(21)</sup> A.S. Johnson,<sup>(21)</sup> J.R. Johnson,<sup>(29)</sup> R. Kajikawa,<sup>(15)</sup> M. Kalelkar,<sup>(20)</sup> H.J. Kang,<sup>(20)</sup> R.R. Kofler,<sup>(11)</sup> R.S. Kroeger,<sup>(12)</sup> M. Langston,<sup>(16)</sup> D.W.G. Leith,<sup>(21)</sup> V. Lia,<sup>(13)</sup> C. Lin,<sup>(11)</sup> G. Mancinelli,<sup>(20)</sup> S. Manly,<sup>(30)</sup> G. Mantovani,<sup>(18)</sup> T.W. Markiewicz,<sup>(21)</sup> T. Maruyama,<sup>(21)</sup> A.K. McKemey,<sup>(3)</sup> R. Messner,<sup>(21)</sup> K.C. Moffeit,<sup>(21)</sup> T.B. Moore,<sup>(30)</sup> M. Morii,<sup>(21)</sup> D. Muller,<sup>(21)</sup> V. Murzin,<sup>(14)</sup> S. Narita,<sup>(24)</sup> U. Nauenberg,<sup>(5)</sup> H. Neal,<sup>(30)</sup> G. Nesom,<sup>(17)</sup> N. Oishi,<sup>(15)</sup> D. Onoprienko,<sup>(23)</sup> L.S. Osborne,<sup>(13)</sup> R.S. Panvini,<sup>(27)</sup> C.H. Park,<sup>(22)</sup> I. Peruzzi,<sup>(8)</sup> M. Piccolo,<sup>(8)</sup> L. Piemontese,<sup>(7)</sup> R.J. Plano,<sup>(20)</sup> R. Prepost,<sup>(29)</sup> C.Y. Prescott,<sup>(21)</sup> B.N. Ratcliff,<sup>(21)</sup> J. Reidy,<sup>(12)</sup> P.L. Reinertsen,<sup>(26)</sup> L.S. Rochester,<sup>(21)</sup> P.C. Rowson,<sup>(21)</sup> J.J. Russell,<sup>(21)</sup> O.H. Saxton,<sup>(21)</sup> T. Schalk,<sup>(26)</sup> B.A. Schumm,<sup>(26)</sup> J. Schwiening,<sup>(21)</sup> V.V. Serbo,<sup>(21)</sup> G. Shapiro,<sup>(10)</sup> N.B. Sinev,<sup>(16)</sup> J.A. Snyder,<sup>(30)</sup> H. Staengle,<sup>(6)</sup> A. Stahl,<sup>(21)</sup> P. Stamer,<sup>(20)</sup> H. Steiner,<sup>(10)</sup> D. Su,<sup>(21)</sup> F. Suekane,<sup>(24)</sup> A. Sugiyama,<sup>(15)</sup> A. Suzuki,<sup>(15)</sup> M. Swartz,<sup>(9)</sup> F.E. Taylor,<sup>(13)</sup> J. Thom,<sup>(21)</sup> E. Torrence,<sup>(13)</sup> T. Usher,<sup>(21)</sup> J. Va’vra,<sup>(21)</sup> R. Verdier,<sup>(13)</sup> D.L. Wagner,<sup>(5)</sup> A.P. Waite,<sup>(21)</sup> S. Walston,<sup>(16)</sup> A.W. Weidemann,<sup>(23)</sup> E.R. Weiss,<sup>(28)</sup> J.S. Whitaker,<sup>(4)</sup> S.H. Williams,<sup>(21)</sup> S. Willocq,<sup>(11)</sup> R.J. Wilson,<sup>(6)</sup> W.J. Wisniewski,<sup>(21)</sup> J.L. Wittlin,<sup>(11)</sup> M. Woods,<sup>(21)</sup> T.R. Wright,<sup>(29)</sup> R.K. Yamamoto,<sup>(13)</sup> J. Yashima,<sup>(24)</sup> S.J. Yellin,<sup>(25)</sup> C.C. Young,<sup>(21)</sup> H. Yuta.<sup>(1)</sup> (The SLD Collaboration) <sup>(1)</sup>Aomori University, Aomori, 030 Japan, <sup>(2)</sup>University of Bristol, Bristol, United Kingdom, <sup>(3)</sup>Brunel University, Uxbridge, Middlesex, UB8 3PH United Kingdom, <sup>(4)</sup>Boston University, Boston, Massachusetts 02215, <sup>(5)</sup>University of Colorado, Boulder, Colorado 80309, <sup>(6)</sup>Colorado State University, Ft. Collins, Colorado 80523, <sup>(7)</sup>INFN Sezione di Ferrara and Universita di Ferrara, I-44100 Ferrara, Italy, <sup>(8)</sup>INFN Laboratori Nazionali di Frascati, I-00044 Frascati, Italy, <sup>(9)</sup>Johns Hopkins University, Baltimore, Maryland 21218-2686, <sup>(10)</sup>Lawrence Berkeley Laboratory, University of California, Berkeley, California 94720, <sup>(11)</sup>University of Massachusetts, Amherst, Massachusetts 01003, <sup>(12)</sup>University of Mississippi, University, Mississippi 38677, <sup>(13)</sup>Massachusetts Institute of Technology, Cambridge, Massachusetts 02139, <sup>(14)</sup>Institute of Nuclear Physics, Moscow State University, 119899 Moscow, Russia, <sup>(15)</sup>Nagoya University, Chikusa-ku, Nagoya, 464 Japan, <sup>(16)</sup>University of Oregon, Eugene, Oregon 97403, <sup>(17)</sup>Oxford University, Oxford, OX1 3RH, United Kingdom, <sup>(18)</sup>INFN Sezione di Perugia and Universita di Perugia, I-06100 Perugia, Italy, <sup>(19)</sup>Rutherford Appleton Laboratory, Chilton, Didcot, Oxon OX11 0QX United Kingdom, <sup>(20)</sup>Rutgers University, Piscataway, New Jersey 08855, <sup>(21)</sup>Stanford Linear Accelerator Center, Stanford University, Stanford, California 94309, <sup>(22)</sup>Soongsil University, Seoul, Korea 156-743, <sup>(23)</sup>University of Tennessee, Knoxville, Tennessee 37996, <sup>(24)</sup>Tohoku University, Sendai, 980 Japan, <sup>(25)</sup>University of California at Santa Barbara, Santa Barbara, California 93106, <sup>(26)</sup>University of California at Santa Cruz, Santa Cruz, California 95064, <sup>(27)</sup>Vanderbilt University, Nashville,Tennessee 37235, <sup>(28)</sup>University of Washington, Seattle, Washington 98105, <sup>(29)</sup>University of Wisconsin, Madison,Wisconsin 53706, <sup>(30)</sup>Yale University, New Haven, Connecticut 06511.
warning/0006/astro-ph0006238.html	ar5iv	text	# The Highest Redshift Radio Galaxies ## 1 Why the highest redshift radio galaxies are interesting Within standard Cold Dark Matter scenarios the formation of galaxies is a hierarchical and biased process. Large galaxies are thought to grow through the merging of smaller systems, and the most massive objects form in over–dense regions, which will eventually evolve into the clusters of galaxies seen today (e.g. Ref. ). It has also been suggested that the first massive black holes may grow in similar hierarchical fashion together with their parent galaxies (e.g. Ref. ) or, because of time scale constraints, may precede galaxy formation and be primordial (e.g. Ref. ). It is therefore of great interest to find the progenitors of the most massive galaxies and their AGN (active massive black holes) at the highest possible redshifts and to study their properties and cosmological evolution. Radio sources are convenient beacons for pinpointing massive elliptical galaxies, at least up to redshifts $`z1`$ (Ref. ; Ref. ). The near–infrared ‘Hubble’ $`Kz`$ relation for such galaxies appears to hold up to $`z=5.2`$, despite large K–correction effects and morphological changes (Ref. ; Fig. 1). This suggests that radio sources may be used to find massive galaxies and their likely progenitors out to very high redshift through near–IR identification. While optical, ‘color–dropout’ techniques have been successfully used to find large numbers of ’normal’ young galaxies (without dominant AGN) at redshifts even surpassing those of quasars and radio galaxies (Ref. ), the radio and near–infrared selection technique has the additional advantage that it is unbiased with respect to the amount of dust extinction. High redshift radio galaxies are therefore also important laboratories for studying the large amounts of dust (e.g. Ref. ) and molecular gas (Ref. ), which are observed to accompany the formation of the first forming massive galaxies. Indeed, a significant part of the scientific rationale for building future large mm-arrays is based on the expectation that to understand galaxy formation will ultimately require understanding their cold gas and dusty environments. Finally, it has been claimed that the (co–moving) space densities of the most powerful radio galaxies and quasars were much higher near $`z2`$, but that they drop off precipitously at even higher redshifts (Ref. ; Ref. ). However, using recently completed studies of moderately faint radio galaxies (Ref. ) it has been argued that here is no such evidence for a redshift cut–off and that these previous results have been biased due to unknown radio K–correction, and thus radio spectral index, trends and associated selection effects. ## 2 How to find the Highest Redshift Radio Galaxies The near–infrared ‘Hubble’ $`Kz`$ relation for radio galaxies (Fig. 1) provides a convenient tool for finding radio galaxies at ever larger redshifts. This was shown convincingly for the first time by Lilly (Ref. ) who found that one of the faintest near-IR radio source identifications in a complete, flux limited sample of $`70`$ objects, B3 0924+34, was a redshift $`z=3.395`$. Unfortunately in complete, flux–limited samples the vast majority of the sources will be relatively nearby, or at only modest redshifts. Lilly (Ref. ) in his survey found only 1 / 70 radio galaxies at $`z>3`$, and McCarthy et al. in a similar but $`7`$ times larger survey also only found one (Ref. ). However, one can pre-select very good HzRG candidates from the radio catalogs, before even going to the telescope, by choosing sources with ultra–steep radio spectra or ‘red radio color’. It is already known more than 20 years that the identification fraction of radio sources on the POSS plates decreases with increasing spectral index (Ref. ). One had to wait for the much more sensitive CCD detectors before further progress in identifying ultra-steep spectrum (USS) sources could be made. One of the first HzRGs which was then found, using the Kitt Peak 4m, was the radio galaxy 4C41.17, at $`z=3.800`$ (Ref. ). This source was the record holder for many years, until it was by-passed, using the same USS method, by 8C 1435+635 ($`z=4.25`$; Ref. ) and 6C 0140+326 ($`z=4.41`$; Ref. ). Together with graduate student C. De Breuck and colleagues at Leiden Observatory we therefore defined the ‘ultimate’ USS source sample by using several new, large radio surveys (De Breuck et al. 2000$`a`$ \[astro-ph/0002297\]). The sample consists of 669 sources with extremely steep radio continuum spectra ($`\alpha 1.3`$; Table 1; Fig. 3), at 10 – 100 times lower flux density limits than has been possible before (Ref. ; Ref. ; Ref. ). To identify these sources we first looked at the POSS and found that approximately $``$ 15% of the sources could be identified, usually with moderately bright galaxies in galaxy clusters. This identification fraction appears to be independent of spectral index (Fig. 3), at least for $`\alpha 1.3`$, in support of the idea that these are mostly foreground objects. The USS selection proved to be extremely efficient. Attempts to obtain optical identifications of USS sources using 3m–4m–class telescopes ($`R\mathrm{}<24`$) were largely unsuccessful. Also near–IR imaging would be very difficult, given the typically expected $`RK4`$ values of high redshift radio galaxies . We therefore decided to entirely skip the optical identification program at Lick Observatory and go straight to near–IR imaging at the Keck I telescope. ## 3 Morphological evolution of the highest redshift radio galaxies When we started our near–IR imaging program at Keck our first order of business was to observe high redshift radio galaxies with known redshifts $`z>1.9`$ to investigate their morphological evolution and to obtain more accurate photometry to study the Hubble $`Kz`$ diagram at the highest redshifts. We obtained near–IR images of 15 HzRGs with $`1.9<z<4.4`$ with the Near Infrared Camera (NIRC, Ref. ) at the Keck I telescope. The images show that there is strong morphological evolution at rest–frame optical ($`\lambda _{\mathrm{rest}}>4000`$Å) wavelengths (Ref. ; Fig. 4). At the highest redshifts, $`z>3`$, the rest–frame visual morphologies exhibit structure on at least two different scales: relatively bright, compact components with typical sizes of $`1`$ $`^{\prime \prime }`$($``$10 kpc) surrounded by large–scale ($``$ 50 kpc) diffuse emission. The brightest components are often aligned with the radio sources, and their individual luminosities are $`M_B20`$ to $`22.`$ For comparison, present–epoch L galaxies and, perhaps more appropriately, ultraluminous infrared starburst galaxies, have, on average, $`M_B21.0`$. The total, integrated rest–frame B–band luminosities are $`35`$ magnitudes more luminous than present epoch $`L_{}`$ galaxies. At lower redshifts, $`z<3`$, the rest–frame optical morphologies become smaller, more centrally concentrated, and less aligned with the radio structure. Galaxy surface brightness profiles for the $`z<3`$ HzRGs are much steeper than those of at $`z>3`$. We attempted to fit the $`z<3`$ surface brightness profiles with a de Vaucouleurs r<sup>1/4</sup> law and with an exponential law, the forms commonly used to fit elliptical and spiral galaxy profiles, respectively. We demonstrate the fitting for our best resolved object at $`z<3`$, 3C 257 at $`z=2.474`$ (Fig. 4). Within the limited dynamical range of the data, both functional forms fit the observed profiles—neither is preferred. Interestingly, despite this strong morphological evolution the $`Kz`$ Hubble diagram for the most luminous radio galaxies remains valid even at the highest redshifts, where a large fraction of the K-band continuum is due to a radio–aligned component. Having established that the $`Kz`$ diagram for high redshift radio galaxies holds even at the highest known redshifts we embarked on our identification program of USS selected sources. Our typical method of observation would be to begin with 16 x 1 minute exposures (1 minute consisting of 2 or 3 co–added frames), start a second 16 x 1 minute run while reducing the first set of observations using DIMSUM. (DIMSUM is the Deep Infrared Mosaicing Software package, developed by P. Eisenhardt, M. Dickinson, A. Stanford, and J. Ward, which is available as a contributed package in IRAF.) If we could identify our target we would break off our second observation, or, if the identification was faint, would let it finish and then go on to the next target. This ‘on–line’ way of observing turned out to be very efficient and resulted in a 100% identification rate with good photometric magnitudes and has provided excellent high redshift radio galaxy candidates using the Hubble $`Kz`$ diagram (Fig. 1; Fig. 5). Often more than a dozen near–IR identifications could be obtained this way in a single night. ## 4 Spectroscopy of the highest redshift radio galaxies As with our near–IR imaging program, our first spectroscopic observations, using the Low Resolution Imaging Spectrograph (LRIS, Ref. ) were made of high redshift radio galaxies with known redshifts. The main purpose, initially, was to determine the origin of the radio–aligned optical / near–IR features using spectro–polarimetry. As is now well–known the rest–frame optical continua of high redshift radio galaxies are often clumpy and aligned with their associated radio sources (Ref. ; Ref. ). This suggested that there must be a causal connection between their optical morphological appearance and the collimated outflow and/or ionizing radiation from their AGN. The most popular explanations for such an alignment effect are scattered light from hidden or mis–aligned quasars, jet–induced star formation or nebular continuum emission. Evidence for each of these processes has been found. In particular, $`z1`$ and $`z2.5`$ most high redshift radio galaxies are strongly polarized, indicating that a large fraction of the optical continuum is due to scattered light from hidden or mis–aligned quasars (Ref. ; Ref. ;Ref. ). However, deep spectropolarimetry observations of two $`z>3.5`$ radio galaxies (4C 41.17 at $`z=3.800`$ and 6C J1908+722 at $`z=3.534`$) show no polarized continua but instead show evidende for absorption lines from young hot stars (Ref. ). It suggests that at the highest redshifts radio galaxy hosts are dominated by massive starbursts, possibly induced by radio jets (Ref. ;Ref. ) and not by scattered light from their AGN. Subsequently our spectroscopic observations focused on the newly identified USS high redshift radio galaxy candidates. At the present time we have observed and analyzed 34 USS high redshift radio galaxies with the following results. Only 5 of the sources have $`z<2`$, 8 have $`2<z<3`$, 9 have $`3<z<4`$ and 3 sources have $`z>4`$, including one at $`z>5`$. At least 3 sources were not detected in optical continuum, despite $``$ 1 hr or longer integrations with LRIS. All we know of these objects is that they are detected in the near-IR at $`K21`$, and have a radio source identified with them. They may be extremely obscured, or at record high redshifts, with Ly-$`\alpha `$ redshifted to near-IR wavelengths ($`z>8`$). Future observations with near–IR spectrographs may tell. We also found 6 sources with only a continuum detection and no emission–lines. These were all extremely compact USS sources, and may be moderately high redshift ($`1<z<3`$) BL Lac objects, ‘emission–line free quasars’ (cf. Ref. ), or even pulsars (which typically have $`\alpha _{radio}1.6`$, Ref. , and are faint optically). The high redshift radio galaxy spectra, when obtained with sufficient spectral resolution, show nearly all very strong blueward asymmetries in the Ly-$`\alpha `$ emission lines (Ref. ; Ref. ; De Breuck et al. 2000$`b`$ (in preparation); Fig. 6. This is almost certainly due to the presence of cold gas (HI) and dust in the vicinity of the radio galaxies, not just because of cosmological Ly-$`\alpha `$ ‘forest’ absorption in the foreground (although this will contribute as well). Spatially resolved emission line regions show that this absorption can occur over the entire region (up to $``$ 50 kpc; Fig. 6), and is strongest in the smallest radio sources (Ref. ). There is much additional evidence for the presence of large amounts of cold gas in dust in high redshift radio galaxies . Many of the highest redshift radio galaxies have been detected at sub–mm wavelengths, both in continuum (e.g.Ref. ; Archibald et al. 2000 \[astro-ph/0002083\]) and molecular lines (e.g.Ref. ). These observations indicate total dust masses of 10<sup>8</sup> – 10<sup>9</sup> $`\mathrm{M}_{}`$, and star formation rates of more than 1000 $`\mathrm{M}_{}`$/yr. Thus high redshift radio galaxies indeed appear to be massive forming systems. One object deserves special mention: 6C J1908+722 at $`z=3.53`$ (Ref. ; Fig. 8). The source shows very broad absorption lines in several of its UV resonant lines (CIV, SIV, NV, Ly-$`\alpha `$ ; Fig. 8). This was interpreted as being caused by outflow, similar to the classical Broad Absorption Line quasars. However, it is interesting to note that this Broad Absorption Line Radio Galaxy (BALRAG) is hosted by a $`1.5\times 10^{13}`$ $`\mathrm{L}_{}`$Ultra Luminous Infrared Galaxy with $`1.5\times 10^8`$ $`\mathrm{M}_{}`$ in dust, $`5\times 10^{10}`$ $`\mathrm{M}_{}`$ in molecular gas, and has an estimated star formation rate of $`1500`$ $`\mathrm{M}_{}`$/yr (Ref. ). The observed velocity range of the gas is large (530 $`\mathrm{km}\mathrm{s}^1`$), and could be even larger: for another high redshift radio galaxy , 4C60.07 at $`z=3.788`$, Papadopoulos et al. find that the molecular gas is distributed over at least two major components, with a total velocity range $`>1000`$ $`\mathrm{km}\mathrm{s}^1`$. Thus it could very well be that the broad, rest–frame UV, absorption lines in 6C J1908+722 may be due to absorption within the parent galaxy. The large BAL velocity range (Fig. 8) could then be caused by a number of cold gas components in the foreground to 6C J1908+722, and which could be falling in or merging with the galaxy. In that case one would expect that the BAL system would be resolved at higher spectral resolution and observations at Keck to test this are planned. ## 5 The Highest Redshift Radio Galaxies ### 5.1 TN J1338$``$1942 at $`z=4.11`$ The first $`z>4`$ USS radio galaxy discovered by us was TN J1338$``$1942. The initial identification was made with the ESO 3.6m at R–band, and subsequent spectroscopy with that same telescope showed that the radio galaxy has a redshift of $`z=4.11\pm 0.02`$, based on a strong detection of Ly-$`\alpha `$ , and weak confirming $`\mathrm{IV}`$ $`\lambda `$ 1549 and He $`\mathrm{II}`$ $`\lambda `$ 1640 (Ref. ). Subsequently we obtained a deep K–band image (rest–frame B–band) at Keck, shown in Fig. 10 overlaid with a VLA radio image (Ref. ). The Ly-$`\alpha `$ and rest–frame optical emission appear co–spatial with the brightest radio hotspot of this very asymmetric radio source. Such asymmetric radio sources are not uncommon, even in the local Universe, and are usuallly thought to be due to strong interaction of one of its radio lobes with very dense gas. A similar asymmetric radio/optical/emission-line morphology has also been seen in the $`z=3.800`$ radio galaxy 4C41.17, where it has been interpreted as being caused by jet-induced star formation (Ref. ;Ref. ; Bicknell et al. 2000 \[astro-ph/9909218\]). With the Keck K-band the identification and astrometry for TN J1338$``$1942 secured we next obtained a high signal–to–noise, medium resolution (5.5 Å FWHM) spectrum using the VLT Antu telescope (Ref. ; Fig. 10). The spectrum of TN J1338$``$1942 is dominated by a bright Ly-$`\alpha `$ line ($`W_{\mathrm{Ly}\alpha }^{\mathrm{rest}}`$ = 210Å) which shows deep and broad ($`1400`$ $`\mathrm{km}\mathrm{s}^1`$) blue–ward absorption, and relatively bright ($`F_{1400}2\mu `$Jy) UV-continuum. In fact, at optical wavelengths, TN J1338$``$1942 turned out to be the most luminous of its kind (Table 2). If all the UV continuum in TN J1338$``$1942 would be due to young O–B stars the implied SFR, based on the optical data alone and without correction for extinction, would be several hundred $`\mathrm{M}_{}`$/yr, similar to 4C41.17. TN J1338$``$1942 might be another example of a very HzRG in which jet-induced star formation might occur. The Ly-$`\alpha `$ is spatially extended by $``$ 4$`^{\prime \prime }`$ (30 kpc) and has a spectral profile that is very asymmetric with a deficit towards the blue. This blue-ward asymmetry is probably due to absorption of the Ly-$`\alpha `$ photons by cold gas in a turbulent halo surrounding the radio galaxy. Using a simple model, and fitting the Ly-$`\alpha `$ profile with a Gaussian emission function and a single Voigt absorption function, De Breuck et al. estimate that the neutral hydrogen column density must be in the range $`3.513\times 10^{19}`$ cm<sup>-2</sup>, and a total mass of $`210\times 10^7M_{}`$. The bright optical continuum and high S/N data also allowed a measurement of the Lya forest continuum break ( Ly-$`\alpha `$ ’discontinuity’, $`D_A`$), and the Lyman limit. The measured value, $`D_A=0.37\pm 0.1`$, is $`0.2`$ lower than the values found for quasars at comparable redshifts. This might perhaps be due due to a bias towards large $`D_A`$ introduced in high–redshift quasar samples that are selected on the basis of large color gradients. The true space density of optically selected quasars, – and Lyman break galaxies –, may have been underestimated and the average HI column density along cosmological lines of sight might have been overestimated. Because of their radio-based, non–color selection, $`z>4`$ radio galaxies may be excellent objects for investigating $`D_A`$ statistics. ### 5.2 TN J0924$``$2201 at $`z=5.19`$ TN J0924$``$2201 is one of the steepest spectrum sources in our USS sample ($`\alpha _{365\mathrm{M}\mathrm{H}\mathrm{z}}^{1.4\mathrm{GHz}}=1.63`$) and therefore was one of our primary targets for near–IR identification. A deep K–band image at Keck showed indeed a very faint ($`K=21.3\pm 0.3`$), multi–component object at the position of the small (1.2$`^{\prime \prime }`$) radio source (Fig. 12). The expected redshift on the basis of the $`Kz`$ diagram was $`z>5`$, and spectroscopic observations at Keck showed that this was indeed the case (Fig. 12), based on a single emission line at $`\lambda 7530`$ Å which we identified as Ly-$`\alpha `$ at $`z=5.19`$ (Ref. ; none of the $`z>5`$ galaxies have more than one line detection). Among all radio selected high redshift radio galaxies TN J0924$``$2201 is fairly typical in radio luminosity, equivalent width and velocity width (Table 2). It does have the steepest radio spectrum, consistent with the $`\alpha z`$ relationship for powerful radio galaxies (e.g. Ref. ), and also has the smallest linear size. The latter may be evidence of its ‘inevitable youthfulness’ or a dense confining environment, neither of which would be surprising because of its extreme redshift (Ref. ; Ref. ). Among the radio selected high redshift radio galaxies TN J0924$``$2201 appears underluminous in Ly-$`\alpha `$ , together with 8C 1435$`+`$63, which might be caused by absorption in an exceptionally dense cold and dusty medium. Evidence for cold gas and dust in several of the most distant high redshift radio galaxies has been found from sub–mm continuum and CO–line observations (e.g.Ref. ; Ref. ). The second highest redshift radio galaxy currently known listed in Table 2 is VLA J123642+621331 at $`z=4.42`$ (Ref. ). This source was not USS selected and provides a view on the possible selection effects of our USS high redshift radio galaxies . The source is an asymmetric double and although its radio luminosity is about a factor 1000 times lower than that of its much more luminous brothers at similar redshifts, it is still radio loud, with a radio luminosity close to the FRI / FRII break at 408 MHz ($`P_{408}3.2\times 10^{33}`$ erg s<sup>-1</sup> Hz<sup>-1</sup>). Its radio spectrum is steep ($`\alpha _{8.4GHz}^{1.4GHz}1.0`$, using the flux densities given by Waddington et al. ), but not as steep as our USS selected high redshift radio galaxies , and the Ly-$`\alpha `$ luminosity is a factor 5 – 10 times less. Apart from the luminosity these properties are not hugely different from expected on the basis of radio selection. It suggests that less extreme steep spectrum selected samples ($`\alpha <1.0`$) at much lower flux densities ($`\mathrm{}<1`$ mJy) might be used to find many more high redshift radio galaxies at very high redshifts, although with lower efficiency, we suspect, than USS selected samples. Our observations of TN J0924$``$2201 extend the Hubble $`Kz`$ diagram for powerful radio galaxies to $`z=5.19`$, as shown in Fig. 1. Simple stellar evolution models are shown for comparison. Despite the enormous $`k`$–correction effect (from $`U_{\mathrm{rest}}`$ at $`z=5.19`$ to $`K_{\mathrm{rest}}`$ at $`z=0`$) and strong morphological evolution (from radio–aligned to elliptical structures), the $`Kz`$ diagram remains a powerful phenomenological tool for finding radio galaxies at extremely hy redshifts. Deviations from the $`Kz`$ relationship may exist Ref. ; but see Ref. ), and scatter in the $`Kz`$ values appears to increase with redshift. Part of this may be due to lack of S/N or contamination by strong line emission in some of the measurements. The clumpy $`U_{\mathrm{rest}}`$ morphology resembles that of other high redshift radio galaxies (Ref. ; Ref. ) and if it is dominated by star light we derive a SFR of $``$200 M yr<sup>-1</sup>, without any correction for extinction, which may be a factor of several. TN J0924$``$2201 may be a massive, active galaxy in its formative stage, in which the SFR is boosted by jet–induced star formation. For comparison other, ‘normal’ star forming galaxies at $`z>5`$ have 10 – 30 times lower SFR ($`620`$ $`\mathrm{M}_{}`$/yr). At the time of its discovery, December 1998, TN J0924$``$2201 was the most distant AGN known, surpassing even quasars for the first time since their discovery 36 years ago. The recent, serendipitous discovery of a color–selected $`z=5.50`$ quasar (Stern et al. 2000 \[astro-ph/0002338\]) returned the record to optically selected quasars. The presence of AGN at such early epochs in the Universe ($`<`$1 Gyr in most cosmogonies) poses interesting challenges to common theoretical wisdom, which assumes, at least for radio loud AGN, that they are powered by massive (billion solar mass), active black holes. The question how these can form so shortly after the putative Big Bang may prove even more challenging then that of the formation of galaxies (e.g.Ref. ). ## ACKNOWLEDGMENTS WvB thanks his many collaborators for stimulating discussions and fun observing runs. Special thanks to C. De Breuck, who has done much of the work as part of his thesis research, and A. Dey for permission to use Fig. 6, Fig. 8 and Fig. 8. The work by W.v.B. at the University of California Lawrence Livermore National Laboratory was performed under the auspices of the US Department of Energy under contract W-7405-ENG-48. W.v.B. also acknowledges support from several NASA grants in support of high redshift radio galaxy research with HST. This figure is available as fig4.jpg This figure is available as fig5.jpg This figure is available as fig6.jpg This figure is available as fig11.jpg
warning/0006/nlin0006019.html	ar5iv	text	# Chaotic Advection near 3-Vortex Collapse ## I Introduction The understanding of the advection of passive tracers is of fundamental interest for many different fields, ranging from a pure mathematical problem, to transport or mixing related ones. One area extensively studied in the last decade is the so called chaotic advectionAref84 -Crisanti92 . This phenomenon, resulting from a chaotic nature of Lagrangian trajectories, enhances the mixing of tracers in laminar flows, while in the absence of chaotic advection the mixing relies on the much less efficient mechanism of molecular diffusion. Chaotic advection in geophysical flows is one of the important areas of application, where the advected quantities vary from the ozone in the stratosphere to various pollutants in the atmosphere and ocean, or such scalar quantities as temperature or salinity. The interest in the geophysical flows increases the practical significance of two-dimensional models and more specifically the advection in the system of vortices provenzale99 -Meleshko93 . In addition to the large scale geophysical flows, $`2D`$ decaying turbulence is another example, where the inverse cascade of energy generates coherent structures (vortices) which dominate the evolution of the flow Benzi86 -Carnevale91 . This type of problems represents only one facet of the interest inherent to the advection in few-vortices systems. Another facet is related to a transport of advected particles. It is known from different observations and numerous models, that the transport of advected particles is anomalous and, in one or another way, can be linked to the Levy-type processes and their generalizations Chernikov90 -Kovalyov2000 . Although these results pertain to fairly simple flows and models, there are speculations relating the chaotic advection in low-dimensional flows to particle dispersion in turbulent flows (see for example a discussion in zaslavsky93.2 ). The interest to the chaotic advection in a three-vortex system is special not only for the reasons mentioned above. The three-vortex system is integrable, and its dynamics can be described in an explicit analytical form. An addition of a tracer (that can be regarded as another vortex of vanishing circulation) brings the number of particles in the system to four, which is a minimum number, from which the point vortex chaos begins Novikov78 ; ArefPomph80 ; the relative simplicity of the system permits to study anomalous transport in considerable detail. A discussion on the importance of three-vortex systems can be found in Aref99 . In this article, we investigate dynamical and statistical properties of the advection in flows produced by three point vortices with different signs. For specific conditions on both the initial position of the vortices and their strength, the collapse of the three vortices to a single point is then possible. We first summarize the motivations for this work. It is known that for systems involving a large number of vortices, the local density of vortices fluctuates, these fluctuations are related to situations when few vortices are close to each other and form a cluster. The number of vortices included in theses cluster vary, but the typical number is around 2-3 Weiss98 ; Novikov75 ; Sire99 ; Zabusky96 , while cluster involving 4 vortices or more are much less probable. A given cluster exists only for a definite finite time $`\mathrm{\Delta }t`$. During its existence, the notion of space-time structure of the cluster (or its configuration) can be introduced, and hence the influence of this structure on typical tracers motion and transport can be studied. Among different clusters, long time transport properties will be most influenced by clusters of vortices with the largest lifetime. The typical size of a cluster can be approximately defined by the minimum distance obtained between two vortices. The closer the vortices are to each other, the stronger is their mutual interaction and the longer the cluster lives. In order to change the inter-vortex distances it is necessary to consider a minimum of three-vortex interactions Aref79 ; Aref99 ; Novikov75 . And to bring all vortices close together (strong interaction) the configuration between the three vortices has to be close to to a configuration leading to a collapse of the vortices Zabusky96 . A typical lifetime estimation of these type of vortices can be assumed to be linked to the period of the motion of three vortices which are evolving on a close to collapse course. It is shown in this paper that the closer the vortices are to a collapse configuration the larger the period of the motion is, and that the growth of the period is exponential with respect to the closeness to the collapse configuration (see Appendix). We then expect that some clusters may have arbitrary long lifetime. Having in mind to shed some light on long time transport properties of systems involving many vortices, we decided to consider the simpler but nevertheless crucial case of transport properties in 3-vortex systems for close to collapse motion; this situations also provides an extreme situation best suited to test a possible universal transport property for three vortex flows. A simple way to define how close a motion is from the exact collapse course is through the deviations of the systems initial parameters from the conditions necessary for collapse. For instance one condition for collapse is $`K=0`$ (where $`K`$ is related to the angular momentum of the system), a close to collapse motion then necessarily satisfies $`\|K\|=\epsilon 1`$, and $`\epsilon `$ is the “distance” from the collapse condition. Hence a special attention should be made to the notion of “close to collapse motion”, as in the whole paper we are referring as close for a small “distance” of the system in parameter space from the collapse conditions. And we insist that for “close to collapse” motions in the previous sense, no collapse of the vortices occurs and the minimum distance of approach between the vortices is finite. In this paper the motion of the vortices is chosen periodic, which retains the possibility of investigating long time transport properties of these types of flows. To conduct our study, we use the methodology and the results of our previous works. The dynamics of tracers is analyzed in a spirit similar to KZ98.1 , where the structure of the advection pattern (chaotic sea, resonant islands, stochastic layers, coherent cores, etc) was investigated numerically and analytically for the case of three identical vortices. Transport properties of the advection in that particular case (where no collapse or near-collapse vortex motion is possible) were found to be anomalous in KZ2000 . When vortex circulations have different signs, their dynamics may change considerably; the collapse phenomenon being one of the most striking examples. The motion of vortices in the vicinity of the collapse (when the collapse conditions are just slightly violated) was studied in LKZ2000 , where different routes in parameter space all leading to the collapse were outlined. While the kinetics of advected particles in non-collapsing three-vortex flows was described in KZ2000 , the situation with the advection in the near-collapse flows remained unclear. In this article we describe different topological structures of the advection pattern, depending on how far is the vortex system from the collapse conditions. We can speculate, that the characteristics of the transport of 3-vortex system can be extended to the case of many-vortex flows, since the advected particle finds itself quite often in the vicinity of a cluster of 3-4 vortices, that define the advection kinetics for a fairly long time-span. For our study, we chose three different cases, corresponding to a specific route to collapse: a far from collapse situation, an intermediate one, and one near the collapse. In these three cases the collapse of the vortices is never reached and the motion of the vortices is periodic, which allows the study of advection for large times and a better understanding of transport in three vortex flows, as we consider here an extreme case with bounded motion, namely the vicinity of collapse configuration. In the Section 2, we present the basic equations of the point vortex dynamics, and discuss the collapse conditions in a three-vortex system. These results are based on the previous studies of the point vortex collapse Aref79 -Kimura90 . We use the notations introduced in LKZ2000 , and develop some argumentation related to our choice of approaching the initial configuration of the vortex system corresponding to the collapse configuration. Advection equations are introduced in the Section 3, where we present different tools used to investigate the dynamic properties of tracers. We focus on Poincaré sections, topology of the phase space, and trajectories, while in Section 4 we present different statistical results. They include velocity distributions, distributions of displacements and its moments, distribution of the Poincaré recurrences, etc. An important point of Section 4 is to understand better how the different regions of phase space influence the statistical properties of trajectories of advected particles. On the basis of the obtained statistical information we consider kinetics of advected particles in Section 5. Fractional kinetics is involved in that description, and corresponding scalings and characteristic exponents are estimated on the basis of the results of Section 4. Motivation for the “3/2 law” of the transport is speculated, as well as the multi-fractal structure of kinetics. Finally, in the Conclusion we discuss different implications of the obtained results and a possibility to exploit them for the analysis of the advection in multi-vortex systems. ## II Near-collapse vortex dynamics The evolution of a system of $`N`$ point vortices can be described by a Hamiltonian system of $`N`$ interacting particles (see for instance Lamb45 ). The nature of the interaction depends on the geometry of the domain occupied by the fluid. For the case of an unbounded plane, the system’s evolution writes $$k_l\dot{z}_l=2i\frac{H}{\overline{z}_l},\dot{\overline{z}}_l=2i\frac{H}{(k_lz_l)},(l=1,\mathrm{},N)$$ (1) with the Hamiltonian $$H=\frac{1}{2\pi }\underset{l>m}{}k_lk_m\mathrm{ln}\|z_lz_m\|=\frac{1}{4\pi }\mathrm{ln}\mathrm{\Lambda },$$ (2) $`z_l=x_l+iy_l`$ is the complex coordinate of the vortex $`l`$, $`k_l`$ its strength and the couple $`(k_lz_l,\overline{z})`$ are the conjugate variables of the Hamiltonian $`H`$, to whom a new energy parameter is associated $$\mathrm{\Lambda }e^{4\pi H}=\underset{lm}{}\|z_lz_m\|^{k_lk_m},$$ (3) in order to simplify subsequent formulae. The resulting complex velocity field $`v`$ is given by the sum of the individual vortex contributions: $$v(z,t)=\frac{1}{2\pi i}\underset{l=1}{\overset{N}{}}k_l\frac{1}{\overline{z}\overline{z}_l(t)}.$$ (4) When $`z_l`$ evolves according to (1), $`v`$ provides a solution of the two-dimensional Euler equation, describing the dynamics of a singular distribution of vorticity $$\omega (z)=\underset{l=1}{\overset{N}{}}k_l\delta \left(zz_l(t)\right).$$ (5) in an ideal incompressible two-dimensional fluid. The motion equations (1) have, besides the energy, three other conserved quantities resulting from the translational and rotational invariance of $`H`$: $$Q+iP=\underset{l=1}{\overset{N}{}}k_lz_l,L^2=\underset{l=1}{\overset{N}{}}k_l\|z_l\|^2.$$ (6) It can be easily verified, that there are three independent first integrals in involution: $`H`$, $`Q^2+P^2`$ and $`L^2`$, from which it follows, that the motion of three vortices is always integrable. An analysis of possible regimes of the motion of three vortices and their classification can be found in Aref79 ; Tavantzis88 . Among the different types of motion, there is an important special case known as vortex collapse. This motion is available when the sum of inverse vortex circulations (harmonic mean) is zero, $$\underset{l=1}{\overset{3}{}}\frac{1}{k_l}=0$$ (7) and vortex positions positions are such, that the modified angular momentum $`K`$ computed in the reference frame for which the center of vorticity is placed at the origin, vanishes $$K\left(\underset{l=1}{\overset{3}{}}k_l\right)L^2(Q^2+P^2)=\underset{lm}{\overset{3}{}}k_lk_m\|z_iz_j\|^2=0$$ (8) Note, that the conditions (7) and (8) do not specify the motion uniquely, rather, they define a range of energies, for which the collapse is possible (once the values of $`k_l`$ satisfying (7) are given). When both of the conditions are satisfied, the solutions of (1) are singular and all three vortices collide at the center of vorticity in a finite time. Depending on the orientation of the vortex triangle, the collapse time $`t_c`$ can be either positive or negative. The first possibility, $`t_c>0`$, corresponds to an actual collapse, while for $`t_c<0`$ the vortex configuration expands without bounds. These two cases are exact images of each other under the time-reversal symmetry, and we refer to both of them as vortex collapse, keeping in mind, that for $`t_c<0`$ the collapse singularity lies backward in time. During the motion, the vortex configuration stays similar to the initial one, meanwhile the area of the triangle, formed by the three vortices grows/decreases linearly in time. We refer the reader to Aref79 ; Tavantzis88 ; LKZ2000 for more details. When the collapse conditions (7,8) are satisfied only approximately, the motion has a specific “near-collapse” type characterized by an emergence of new scales of distances and velocities, which differ significantly from the initial ones LKZ2000 . A detailed analysis of the near-collapse dynamics of three vortices was performed in LKZ2000 for the case when two of the vortices have the same strength. A multitude of motion regimes was found in the vicinity of collapse; different regimes are distinguished by the way the collapse conditions (7,8) are violated, i.e. whether the combinations $`1/k_l`$ and $`K`$ are greater, less or equal to zero. Moreover, for some of these combinations, the motion type also depends on the energy $`\mathrm{\Lambda }`$ of the vortex configuration; for example, when $`1/k_l<0`$ and $`K>0`$ there exist critical energies $`\mathrm{\Lambda }_{c_1}`$ and $`\mathrm{\Lambda }_{c_2}`$, such that in the range $`\mathrm{\Lambda }_{c_1}>\mathrm{\Lambda }>\mathrm{\Lambda }_{c_2}`$ the motion has two branches of periodic motion, for $`\mathrm{\Lambda }=\mathrm{\Lambda }_{c_2}`$ the motion is aperiodic, and for $`\mathrm{\Lambda }>\mathrm{\Lambda }_{c_2}`$ there is only one periodic branch left. A classification of the near-collapse motion regimes, including the behavior of their length and time scales as the collapse is approached, can be found in LKZ2000 . In order to carry out a detailed study of the advection for close to collapse situations, we restrict our consideration to a specific one-parameter family of near-collapsing vortex systems, defined as follows: 1. Only one of the collapse conditions, the strength condition (7), is violated; the second condition (8) is satisfied, i.e. $`K=0`$. 2. The two positive vortices are identical, by an appropriate choice of time units their strength can be put to 1; in other words, we fix $`k_1=k_2=1`$. The circulation of the third vortex is negative, $`k_3k`$, ($`k>0`$). In this situation, the collapse happens, when the strength of the third vortex reaches a critical value $`k=k_c1/2`$. The “distance” of the system from the collapse can be measured by the amount of the strength detuning $`\delta `$, $$\delta k_ck=1/2k$$ (9) We will be considering only the case $`\delta >0`$. 3. The energy of the vortex configuration $`\mathrm{\Lambda }`$ is fixed to a constant value $`\mathrm{\Lambda }=0.9`$. The first two conditions ensure the relative vortex motion to be periodic, and the choice of energy is such, that in the approach to collapse ($`\delta 0`$), the maximum inter-vortex distance grows very slowly (no noticeable changes in the range of $`\delta `$ considered below) while the minimum inter-vortex distance rapidly approaches zero. It gives us a convenient “collapse in a box” setting, where vortices initially separated by a distance of order one, are brought arbitrary close (controlled by $`\delta `$) together, which may be of interest for studies of transport in many vortex-systems. For instance, the ability to bring vortices closer to each other by orders of magnitude makes these 3-vortex processes a dominant interaction mechanism in a rare gas of vortex patches Zabusky96 . Before proceeding to the discussion of the advection, we will briefly summarize the results from LKZ2000 , pertaining to our case. The dynamics of a three-vortex system with two identical vortices can be mapped to a one-dimensional Hamiltonian system describing a motion of a particle of a unit mass and zero total energy in an effective potential, that depends on the strength of the third vortex $`k`$, and the constants of motion $`\mathrm{\Lambda }`$ and $`K`$ (see appendix). The dynamical variable $`X`$ of this one-dimensional system is equal to the squared distance between the two positive vortices $$X\|z_2z_1\|^2,$$ (10) the other two distances, can be found from the expressions for $`\mathrm{\Lambda }`$ and $`K`$. And in our special case ($`K=0)`$ the motion is confined between two single roots of the potential, $$X_{min}<X<X_{max},$$ (11) where $$X_{min}\left(\frac{1k}{2k}\right)^{k/\delta }\mathrm{\Lambda }^{1/2\delta },X_{max}\left(\frac{1}{2k}\right)^{k/\delta }\mathrm{\Lambda }^{1/2\delta },$$ (12) which implies that $`X(t)`$ is a periodic function of time, and consequently the other inter-vortex distances are too, i.e. the relative motion of vortices in our one-parameter family is always periodic. As the collapse is approached ($`k1/2`$, $`\delta 0`$), the motion tends to cover all length scales: $`X_{min}0`$, $`X_{max}\mathrm{}`$, and its period diverges as $$T\frac{1}{\delta }\mathrm{\Lambda }^{1/2\delta }.$$ (13) (see the appendix for the derivation). The influence of this behavior on the properties of advection is investigated in the next section. ## III Dynamics of the advection A passive particle (tracer) follows the flow according to the advection equation $$\dot{z}=v(z,t)$$ (14) where $`z(t)`$ represent the tracer trajectory, and $`v`$ is the velocity field. In the case of a point vortex system, the velocity field is given by Eq. (4). The incompressibility of the flow allows to write the advection equation (14) in a Hamiltonian form: $$\dot{z}=i\frac{\mathrm{\Psi }}{\overline{z}},\dot{\overline{z}}=i\frac{\mathrm{\Psi }}{z}$$ (15) where a stream function $$\mathrm{\Psi }(z,\overline{z},t)=\frac{1}{2\pi }\underset{l=1}{\overset{3}{}}k_l\mathrm{ln}\|zz_l(t)\|$$ (16) acts as a Hamiltonian. This system in non-autonomous, since the stream function depends on time through the vortex coordinates $`z_l(t)`$. The character of this dependence (periodic or not) is important for the further analysis. Below we will show, that although (16) is quasiperiodic, it can be made periodic by an appropriate coordinate transformation, which means, that the advection in our system has a 1 1/2 degrees of freedom Hamiltonian dynamics. Indeed, as was mentioned in the previous section, the relative vortex motion is periodic, i.e. the vortex triangle repeats its shape after a time $`T`$. This does not imply a periodicity of the absolute motion, since the triangle is rotated by some angle $`\mathrm{\Theta }`$ during this time, see Fig. 1. In general, $`\mathrm{\Theta }`$ is incommensurate with $`2\pi `$, rendering a quasiperiodic time dependence of $`z_l(t)`$. Let us consider a reference frame, rotating around the center of vorticity with an angular velocity $$\mathrm{\Omega }\mathrm{\Theta }/T.$$ (17) In this co-rotating reference frame, vortices return to their original positions in one period of relative motion $`T`$, see Fig. 2, their new coordinates $$\stackrel{~}{z}ze^{i\mathrm{\Omega }t},$$ are periodic functions of time. In the co-rotating frame the advection equation retains its Hamiltonian form with a new stream function $`\stackrel{~}{\mathrm{\Psi }}`$ which acquires an extra (rotational energy) term $$\stackrel{~}{\mathrm{\Psi }}\mathrm{\Psi }+\mathrm{\Omega }^2/2\|z\|^2.$$ (18) An advantage of this new frame is that $`\stackrel{~}{\mathrm{\Psi }}`$ is time-periodic: $$\stackrel{~}{\mathrm{\Psi }}(\stackrel{~}{z},\stackrel{~}{\overline{z}},t+T)=\stackrel{~}{\mathrm{\Psi }}(\stackrel{~}{z},\stackrel{~}{\overline{z}},t)$$ (19) and well-developed techniques for periodically forced Hamiltonian systems can be used to study its solutions. Note, that the one-period rotation angle $`\mathrm{\Theta }`$ is defined modulo $`2\pi `$, making the choice of the co-rotating frame non-unique. We remove this ambiguity by requiring the negative vortex to make no full revolutions around the center of vorticity in the co-rotating frame (as in Fig. 2). This particular choice of $`\mathrm{\Omega }`$ is inconsequential for the further analysis. The one-period rotation angles, relative motion periods and angular velocities of the co-rotating frame are presented in the Table 1. As the collapse is approached ($`k1/2`$), $`T`$ grows rapidly (compare to the formula (13)), and the vortices make more and more turns per period, an acceleration of the rotation speed is also observed. We start our analysis of the advection by numerically constructing Poincaré sections of tracer trajectories (in the co-rotating frame). A Poincaré section is defined as an orbit of a period-one (Poincaré) map $`\widehat{P}`$ $$\widehat{P}z_0=\stackrel{~}{z}(T,z_0)=e^{i\mathrm{\Theta }}z(T,z_0)$$ (20) where $`\stackrel{~}{z}(t,z_0)`$ denotes a solution $`\stackrel{~}{z}(t)`$ with an initial condition $`\stackrel{~}{z}(t=0)=z_0`$. Plots of Poincaré sections for three different values, $`k=0.2`$, $`k=0.3`$ and $`k=0.41`$ are shown in Figures 3-5. Vortex and tracer trajectories were computed using a symplectic fifth-order Gauss-Legendre scheme McLachlan92 . Exact conservation of Poincaré invariants by the symplectic scheme suppresses numerical diffusion, yielding high-resolution phase space portraits. The Poincaré sections presented in Figures 3-5 show an intricate mixture of regions with chaotic and regular tracer dynamics, typical for periodically forced Hamiltonian systems. All three phase portraits share common features with the advection patterns, found in a flow due to three identical point vortices KZ98.1 ; NeufeldTel : the stochastic sea is bounded by a more or less circular domain, there are a number of islands inside it, where the tracer’s motion is predominantly regular. In particular, all three vortices are surrounded by robust near-circular islands, known as vortex cores. Contrary to the case of three identical vortices, where the tracer dynamics is integrable for a special value of vortex energy (when the vortices form a steadily rotating isosceles triangle) and has a near-integrable character in the vicinity, the tracer motion in the near-collapse flow family considered here is always strongly chaotic, the stochastic sea remains a principal element of the advection pattern for any $`k`$. The phase portraits indicate, that the degree of chaotization increases with the approach to collapse. For instance, in the far from collapse case $`k=0.2`$ the islands of regular motion inside the stochastic sea occupy a considerable area, and as $`k`$ increases, their share drastically diminishes. A decrease in the radii of the vortex cores is of particular interest, since they are the robust structures, which also appear in many-vortex systems. The upper bound on the core radii can be obtained from the minimum approach inter-vortex distances. The minimum distance between the two positive vortices (see (12)) is $$R_1^{min}=\left(\frac{1k}{2k}\right)^{k/2\delta }\mathrm{\Lambda }^{1/4\delta },$$ (21) at the same moment the distance between the negative vortex and one of the positive two also reaches its minimum, which is $$R_2^{min}=\frac{1}{2}\left(\sqrt{\frac{2k}{k}}1\right)R_1^{min}.$$ (22) At this moment the vortices are collinear. Since the sum of the core radii of two vortices cannot exceed the minimum distance between them, we get an upper bound for the radius of the positive vortex core, $`R_{core}^+`$ and for that of the negative one, $`R_{core}^{}`$, in terms of $`R_1^{min}`$, $`R_2^{min}`$: $$R_{core}^+=\mathrm{min}(\frac{1}{2}R_1^{min},\frac{1}{1+\sqrt{k}}R_2^{min})$$ $$R_{core}^{}=\frac{\sqrt{k}}{1+\sqrt{k}}R_2^{min}$$ (23) where we took into account, that the minimum distance between the positive and the negative vortex is shared between the corresponding cores depending on relative strength, and choose to define the boundary as the point between the two vortices with minimum speed. The core radii, measured directly from the phase portraits in Figures 345, and their upper bounds, obtained from (23) are listed in the Table 2. The presence of islands where the motion is regular, in the stochastic sea, is known to alter the transport properties of a physical system. This phenomenon is known as “stickiness”; when a passive particle, traveling in the stochastic sea, gets close to an island, it is likely to stick to this island for a while and mimic a regular trajectory of a trapped particle. And since with each island, a whole hierarchy of smaller islands around islands is present, the particle can stick for arbitrary long times, which affects on the whole the transport properties of the system. In KZ2000 , stickiness has been exhibited by measuring recurrence times of particles to a given part of phase space, and plotting the particles position on the map with a color accordingly to their return times. Particles, that have long return times, all stick to a particular island, and do not jump from one to another. Taking these facts into account, we use another way to visualize stickiness. Indeed, sticking particles have all long coherent time behavior, which reflects in quantities such as their angular speed, or intrinsic speed. We decide then to compute the average intrinsic speed over a definite amount of time of an ensemble of particles and record it. The measured quantity is the following $$V_i(m,n)=\frac{1}{nT}_{t_0+nmT}^{t_0+n(m+1)T}v_i(t)𝑑t,$$ (24) where $`n`$ is the number of periods over which the speed is averaged, $`m`$ keeps track of the elapsed time, and $`v_i(t)`$ is the instantaneous speed at time $`t`$ of the particle $`i`$. We then define the distribution of such averaged velocities as $$\rho (V;n,m)=\frac{1}{N_p}\underset{i}{}\delta \left(VV_i(m,n)\right),$$ (25) and smooth it over an interval to obtain a continuous curve. In fact, as can be observed in Fig. 6, for which the speed of particles is plotted versus time, we can notice that after a brief period of time the distribution seems stationary, meaning that $`\rho (V,n,m)`$ is independent of $`m`$ and therefore we can average its value over $`m`$, which in practice allows better statistics. Figure 6 is already informative, as we can notice some darker stripes, which means that some special average velocities are favored. This fits with the picture of some particles sticking to some island for a long time (at least $`>10T`$ here). To obtain this data we computed the trajectories over $`10000`$ periods for a sample of $`253`$ particles and recorded every $`n=10`$ periods. We use the stationarity property and plot the distribution $`\rho `$ versus the speed $`V`$ for the three different cases $`k=0.2`$, $`0.3`$, $`0.41`$, these are represented respectively in the figures Fig. 7, Fig. 8, and Fig. 9. We notice that the dark stripes observed in Fig. 6, correspond to peaks in the density probability. In order to characterize the origin of these peaks, we plotted in Fig. 10, Fig. 11, and Fig. 12, the position of the particles in the phase space contributing to the peaks in the distribution function. As anticipated, each of the observed peak correspond to a specific region of the phase space around some island. Concerning the influence of collapse, we notice that the area occupied by the contributing particles decreases as the critical condition is approached. We may then speculate that the transport properties of the three different system, which we discuss in the next section, may differ in a substantive way. ## IV Anomalous statistical properties of tracers Deterministic description of the motion of a passive particle in the mixing region is impossible, since a local instability produces exponential divergence of trajectories, and after a short time, the position of a tracer would be completely unpredictable. Even the outcome of an idealized numerical experiment is non-deterministic in this situation, since a round-off error is creeping slowly but steadily from the smallest to the observable scale. For this reason, long-time behavior of tracer trajectories in the mixing region is usually studied within a probabilistic approach. In the absence of long-term correlations, a kinetic description, which uses Fokker-Plank-Kolmogorov equation and leads to Gaussian statistics, Zaslavsky92 works fairly well in many cases. Yet, in the present case, the complex topology of the advection pattern, illustrated by the Poincaré sections in Figures 3-5, indicates that one should anticipate anomalous statistical properties of the tracers in the chaotic sea. Singular zones around KAM islands usually produce long-time correlations, which may result in essential changes in the particle kinetics. Although in some cases these “memory effects” can be accounted for by the modification of the diffusion coefficient in the FPK equation Chirikov79 Rechester80 , often their influence is more profound Zaslavsky93 dCN98 Castiglione99 KZ2000 , and leads to a super-diffusive behavior with faster than linear growth of the particle displacement variance: $$(xx)^2t^\mu $$ (26) where the transport exponent $`\mu `$ exceeds the Gaussian value: $`\mu >1`$. In this section we analyze the statistical properties of tracers in the chaotic region for three vortex flow geometries, introduced above: far from collapse ($`k=0.2`$, Fig. 3), intermediate ($`k=0.3`$, Fig. 4), and close to collapse ($`k=0.41`$, Fig. 5). A plot of a time series of the arclength versus time $`s_i(t)=v_i(t)𝑑t`$ for a set of typical tracers trajectories (Fig. 13) reveals an intermittent character of tracer motion: random pieces of trajectory are interrupted by regular flights, some of which are fairly long. To remain consistent with previous work, we focus our interest on the character of tracer rotation, and for that matter, we define its azimuthal coordinate in the center of vorticity reference frame $$\theta (t)\text{Arg}z$$ (27) to be a continuous function of time, i.e. $`\theta (t)(\mathrm{},\mathrm{})`$ keeps track of the number of revolutions performed by a tracer. Mean advection angle $`\theta (t)`$ (here $``$ denotes ensemble average) grows linearly with time: $$\theta (t)=\omega t,$$ (28) the values of the average rotation frequency $`\omega `$ for the three cases are listed in Table 3. The growth of the variance, $$\sigma ^2(t)(\theta (t)\theta (t))^2$$ (29) is faster than linear for all three cases: angular tracer diffusion is anomalous. From log-log plots of $`\sigma ^2(t)`$ versus time in Fig. 14, one may see, that in order to describe the growth of the variance with a power law $$\sigma ^2(t)t^\mu $$ (30) one has to introduce different transport exponents for different time ranges. The values of these exponents, obtained by linear fits of corresponding parts of the graphs of Fig. 14, are given in Table 3. Below, the first time range (with the exponent $`\mu _1`$) will be referred to as “short times”, and the second one ($`\mu _2`$) “long times”. Recently Castiglione99 , two types of anomalous diffusion were distinguished by the behavior of the moments, other than variance. The case when the evolution of all of the moments can be described by a single self-similarity exponent $`\nu `$ according to $$\|xx\|^qt^{q\nu }$$ (31) was called “weak anomalous diffusion”, whereas the case when $`\nu `$ in (31) is not constant, i.e. $$\|xx\|^qt^{q\nu (q)}$$ (32) was named “strong anomalous diffusion”. The importance of this distinction comes from the fact, that in the weak case the PDF must evolve in a self-similar way: $$P(x,t)=t^\nu f(\xi ),\xi t^\nu (xx)$$ (33) whereas non-constant $`\nu (q)`$ in (32) precludes such self-similarity. Note, that a self-similar PDF evolution can have a more general form than (33), with $`t^\nu `$ replaced by an arbitrary decaying function of time $`g(t)`$: $$P(x,t)=g(t)f(g(t)(xx)),$$ (34) which means, that if for different time decades $`g(t)`$ has different asymptotics, the self-similarity exponent $`\nu `$ will change from one decade to another. This variation of $`\nu `$ with time (in particular, differences of transport exponents for short and long times in Table 3) is not related to the type (strong or weak) of anomalous diffusion. We have performed the measurements of a set of moments of tracer angular PDF (including non-integer values of $`q`$) defined as: $$M_q(t)\|\theta (t)\theta (t)\|^q$$ (35) for the three vortex geometries. Time evolution of each moment was fitted by a power law: $$M_q(t)t^{\mu (q)}$$ (36) separately for short, and for long times. The results are summarized in Figures 15, 16, 17 (short times) and 18, 19, 20 (long times), where the exponents $`\mu (q)`$ are plotted versus the moment number $`q`$. In all cases, the apparent absence of a single linear fit indicates the presence of strong anomalous diffusion. This property was also found in KZ2000 (by comparison of the scaling properties of the central part of tracer PDF with the behavior of the variance) in a flow due to three vortices of equal strength. Thus, strong anomalous diffusion is a generic property of advection in three vortex flows. Our results show, that $`\mu (q)`$ is well approximated by a piecewise linear function of the form: $$\mu (q)=\{\begin{array}{ccc}\nu q& \text{for}& q<q_c\\ qc& \text{if}& q>q_c\end{array}$$ (37) where $`c`$ is a constant, and $`q_c`$ is a crossover moment number $`q_c=c/(1\nu )`$. In Castiglione99 , where this form was introduced, it was found, that it fits fairly well the numerically obtained values of $`\mu (q)`$ in all cases of strong anomalous diffusion, considered there, although a theoretical example of a system with arbitrary (concave) $`\mu (q)`$ was mentioned. Note, that deviations from (37), occurring in the crossover region $`qq_c`$, are probably a result of finite observation time, and the form (37) might be precise in the limit $`t\mathrm{}`$. As we have mentioned, the non-constant $`\nu (q)`$ in (32) is incompatible with the self-similar evolution of tracer distribution. Let us introduce an “almost self-similar distribution” $$P(x,t)=\{\begin{array}{ccc}t^\nu f(t^\nu x)& \text{for}& xvt\\ 0& \text{if}& x>vt\end{array}$$ (38) where the exact self-similarity (33) is broken only in the time-dependence of the cut-off. This modification of the exact self-similarity relation (33) takes into account the fact, that tracer speed is bounded by a maximum speed $`v`$. If $`f(\xi )`$ decays sufficiently fast at infinity (e.g. like a Gaussian), the cutoff behavior is irrelevant and the moments of (38) follow (31), but if $`f(\xi )`$ has a power tail, the almost self-similar distribution (38) will give the piecewise linear form (37) for the moments. If $`f(\xi )\xi ^\beta `$ for large $`\xi `$, than low moments $`M_q`$ with $`q<\beta 1`$ will be determined by the central, self-similar part of the distribution, and high moments ($`q>\beta 1`$) by the cutoff value, $$M_q(t)\{\begin{array}{ccc}t^{\nu q}& \text{for}& q<\beta 1\\ t^{q(1\nu )(\beta 1)}& \text{if}& q>\beta 1\end{array}$$ (39) which is equivalent to (37) with $$c=(1\nu )(\beta 1)$$ (40) We may conclude, that the piecewise linear dependence of the exponent $`\mu (q)`$ on the moment number $`q`$ is a signature of an almost self-similar evolution of tracer distribution with a long-tailed $`f(\xi )`$. The constant $`c`$ in (37) is related to the self-similarity exponent $`\nu `$ and power law decay exponent $`\beta `$ of $`f(\xi )`$ by (40). Another consequence of the intermittent character of tracer motion is an anomalous distribution of recurrences of the Poincaré map of tracer trajectories. To define recurrences, we take a region $`B`$ in the chaotic sea, and register all returns of a Poincaré map trajectory into $`B`$. The length of a recurrence is a time interval between two successive returns. In a system with perfect mixing, the PDF of recurrence lengths obeys a Poissonian law, provided $`B`$ is small enough, and decay of the long-recurrence tail of the distribution is exponential for any $`B`$. Recurrence distributions for tracers in all three cases ($`k=0.2`$, $`k=0.3`$ and $`k=0.41`$) are shown in Figures 21, 22, 23. The plots show, that all distributions have long tails, indicating, that between the returns tracers are being trapped in long flights of highly correlated motion. The form of the graphs suggests, that long recurrences are distributed according to a power law $$P(t)t^\gamma $$ (41) The values of the exponent $`\gamma (\delta )`$ are: $$\gamma (0.2)=2.2\gamma (0.3)=2.4\gamma (0.41)=3.1.$$ (42) Note, that while the collapse configuration is approached, the value of the exponent increases, which may be interpreted as an improvement in mixing properties of the flow. This agrees with the changes in the structure of Poincaré section (Fig. 10-12): the closer to collapse we get, the bigger part of the chaotic domain is occupied by a well-mixed area, and the smaller is the role of the singular zones around KAM islands. In fact one can try to find out the influence of the different islands on transport by using the distributions illustrated in Figures 7-9. Indeed, each island corresponds to a specific peak. We recompute the moments of the distribution in the far from collapse case $`k=0.2`$, for the modified data set, where the trajectories, corresponding to a specific peak are discarded. The result is presented in Fig. 24. We notice that the cores do affect the transport, but their influence is essentially visible for the high moments, while the slow particles trapped in the outer rim, are mainly responsible for the low moments; we also notice that we do not observe the change in slope of the strong anomalous behavior anymore, and conclude that the strong anomalous feature is due to the interplay of the different structures in the phase plane. ## V Kinetics of advected particles In some of the previous publications (see, for example, Chernikov90 ZEN97 -ZasEdel2000 ) it was clearly indicated that the properties of anomalous transport are sensitive to phase space topology. More specifically, if we use the fractional kinetic equation Zaslavsky92 ; ZEN97 in the form $$\frac{^\beta P(\theta ,t)}{t^\beta }=𝒟\frac{^\alpha P(\theta ,t)}{\|\theta \|^\alpha }$$ (43) to describe distributions $`P(\theta ,t)`$ of rotations over angle $`\theta `$, then the transport coefficient $`𝒟`$ and exponents $`(\alpha ,\beta )`$ depend on the presence of different structures such as boundaries of the domain, islands, cantori, etc. The results of Section 4 show the stickiness of trajectories of advected particles to the boundary of the domain and to boundaries of islands. This phenomenon is similar to what has been observed in KZ98.1 for the same-sign vortices. Our goal of this section is to estimate the values of the exponents $`\alpha ,\beta `$. Figures 10-12 demonstrate stickiness of trajectories to specific structures with a filamentation of sticky domains along stable/unstable manifolds. In fact, different sticky domains generate different intermittent scenarios with some associated values of $`(\alpha ,\beta )`$ ZEN97 ; ZasEdel2000 . As a result, the real kinetics is multi-fractional and can be characterized by a set of values of $`(\alpha ,\beta )`$ or, more precisely, by a spectral function of $`(\alpha ,\beta )`$ in the same sense as the spectral function for multi-fractals Hentschel84 -Jensen85 . Figures 7-9 show that trajectories, sticking to different structures (islands), have different angular velocities (compare to peaks in Figures 7-9). Due to this, different asymptotics to the distribution function $`P(\theta ,t)`$ and different values of $`(\alpha ,\beta )`$ will appear for different time intervals. In other words, for a considered time interval one can expect a specific “intermediate asymptotics” for $`P(\theta ,t)`$ and, correspondingly, different pairs $`(\alpha ,\beta )`$. Different classes of universality for the values $`(\alpha ,\beta )`$ were discussed in ZasEdel2000 . Below we will apply some of these results. Multiplying (43) by $`\|\theta \|^\alpha `$ and integrating it over $`\|\theta \|`$ we obtain $$\|\theta \|^\alpha t^\beta $$ (44) or, in the case of self-similarity the transport exponent $`\mu `$ from the equation $$\|\theta \|^2t^\mu $$ (45) can be estimated as $$\mu =2\beta /\alpha $$ (46) Expression (45) should be considered with some reservations since the second and higher moment may diverge. For a finite time $`t<t_{\text{max}}`$ particles reach a distance (angular rotation) $`\theta <\theta _{max}`$, which makes all moments finite. Typically all $$\theta _{max}=\omega _{max}t$$ (47) and $`\omega _{max}`$ (maximal angular velocity) can be reached only at the boundary of the domain of chaotic motion (see Figs. 6-8). Using notations (31), (32), and $$\mu (2)\mu $$ (48) we can present the results for the transport exponents $`\mu `$ in Table 4. They are almost the same independently of how far is the control parameter $`k`$ from its critical value $`k_c=1/2`$ (the collapse condition). For large values of $`q`$ we have $`\mu `$ close to $`\mu =2`$ which corresponds to ballistic dynamics with $`\alpha \beta 1`$. This result can be well understood from the stickiness of trajectories to the cores (see Figs. 10-12 in the black color). As it follows from distributions in Figs. 7-9, the particles that stick to the cores are the fastest ones, and they just define the large moments values. The value of $`\mu `$ for $`q=2`$ is defined mainly by meso-structures in the middle of Figs. 10-12 (light gray). A typical property of these structures is existence of islands with well resolved filamentations due to the vicinity of the structures to a bifurcation. The latter is evident from the sharp corners of islands, which may indicate a parabolic type singular point Melnikov96 . A corresponding effective Hamiltonian, describing dynamics near a singular point, has a form Karney83 ; Melnikov96 ; ZEN97 ; Rom-Kedar99 : $$H_{\text{eff}}=a_1(\mathrm{\Delta }P)^2+a_2\mathrm{\Delta }QQ_3(\mathrm{\Delta }Q)^3$$ (49) where $`(P,Q)`$ are generalized momentum and coordinate and $`(\mathrm{\Delta }P,\mathrm{\Delta }Q)`$ are their corresponding deviations from the singular point $`(P_0,Q_0)`$: $$\mathrm{\Delta }P=PP_0,\mathrm{\Delta }Q=QQ_0.$$ (50) Particularly, it may be $$Q=\theta ,P=\dot{\theta }.$$ (51) Depending on the coefficients $`a_j`$ and on the meaning of variables $`(P,Q)`$, which may be different from (51), one can describe singularity due to bifurcations for different types of dynamical modes: accelerator mode Karney83 ; Melnikov96 , “blinking island” mode Melnikov96 , ballistic mode Rom-Kedar99 , etc. For all these situations, universality of the Hamiltonian (49) permits estimation of the exponents $`(\alpha ,\beta )`$ in (43) A trajectory that approaches the vicinity of the singular point (or, simply, a corner of the island boundary), behaves intermittently and escapes the near-separatrix boundary layer. The phase volume of the escaping trajectories is $$\delta \mathrm{\Gamma }=\delta P\delta Q$$ (52) where $`\delta P,\delta Q`$ are values $`\mathrm{\Delta }P,\mathrm{\Delta }Q`$ related to the escaping particles. From (49) we can estimate $$\delta P_{\text{max}}\delta Q^{3/2}$$ (53) and from (52), (53) $$\delta \mathrm{\Gamma }\delta Q^{5/2}$$ (54) Escaping from the boundary layer means growth of the “radial” variable $`\delta Q`$ with time, i.e. for an initial time interval $`\delta Qt`$, and consequently, $$\delta \mathrm{\Gamma }t^{5/2}.$$ (55) From (55) we conclude for the escape probability density to leave the boundary layer at time instant $`t`$ within interval $`dt`$: $$\psi (t)1/\delta \mathrm{\Gamma }t^{5/2}.$$ (56) It was shown in Montroll84 that under special conditions the exponent $`\gamma `$ for the trapping time asymptotic distribution $$\psi (t)t^\gamma $$ (57) can be linked to fractal time dimension. Moreover, $`\gamma `$ is related to the kinetic equation (43) as ZEN97 $$\beta =\gamma 1.$$ (58) For the considered case we have $`\beta =3/2`$. For the spatial distribution of particles, the simplest situation occurs when the diffusion process has Gaussian type and, consequently, $`\alpha =2`$. In the case of the presence of hierarchical set of islands, $`\alpha `$ can be defined through scaling properties of the island areas. In the considered situation random walk is more or less uniform but trajectories are entangled near stable/unstable manifolds, i.e. in the light gray areas of Figs. 10-12. That means that $`\alpha 2`$ although it is not exactly 2. Finally, we arrive to: $$\mu =2\beta /\alpha 3/2$$ (59) in correspondence to observations in the Table 4. The value (59) was also discussed in ZasEdel2000 as one of possible universal values for the transport exponent $`\mu `$. We need to comment that it is not worthwhile to try to obtain $`\mu `$ with a higher accuracy since a specific value of $`\mu `$ has no meaning due to multi-fractal nature of transport ZasEdel2000 . It is also important that we have considered such values of the control parameter $`k`$ for which there exists a strong filamentation. That guarantees a possibility of using Eq. (49) and the following analysis. ## VI Conclusion We have considered the dynamical and statistical properties of the passive particle advection in a family of flows, created by three point vortices of different signs. In all three particular cases, investigated numerically, tracer advection was strongly chaotic: advection patterns, visualized via Poincare sections of tracer trajectories are dominated by a well developed stochastic sea, occupying most of the area around the center of vorticity. With the approach of the vortex system to the collapse configuration, the degree of tracer chaotization increases: the stochastic sea grows, expanding outward and consuming some of the inner resonant islands. The statistics of the tracers in the chaotic region is non-Gaussian. Anomalous diffusion (faster than linear growth of variance) with different time and space scales was found in all three cases, as well as non-Poissonian distributions of Poincaré recurrences (with power law decay of long recurrence probability). We did not find normal transport regimes, if such regime exist, they are confined to narrow windows in the parameter domain. Transport anomalies are caused by the phenomenon of stickiness of the chaotic trajectories to the highly structured boundaries of the chaotic region. In the cases considered, three important types of boundaries can be distinguished: external border of the chaotic sea, boundaries of the resonant islands inside the chaotic sea, and boundaries of the vortex cores. Each of these influences various aspects of tracer statistics, analysis of their separate contributions shows, that the vortex cores, that rotate with the fastest rate, determine the high moments of the tracer distribution, while the external boundary, being the slowest, but the most sticky, dominate the low moments. Vortex cores appeared in simulations Babiano ; NeufeldTel ; KZ98.1 , their origin and sizes were derived in KZ98.1 for a system of three identical vortices; particularly it was shown that the cores are the islands of stability filled by invariant curves and extremely thin stochastic layers. As the control parameter $`k`$ approaches the collapse value $`k_c=1/2`$, the sizes of the vortex cores noticeably decrease. An upper bound of the core radii, obtained from the minimum distance of vortex approach to each other, gives a good estimation for both positive and negative vortex core size. Although the transport possesses multi-fractal features, it can be successively described by a fractional kinetic equation with characteristic exponents $`\alpha 2`$ and $`\beta 3/2`$. A corresponding moments dependence is $$\|\theta \|^\alpha t^\beta .$$ (60) The transport can be characterized by strong intermittency which manifests itself in strong deviation from (60) for higher moments, i.e. $$\|\theta \|^{2m}t^{\mu (m)}$$ (61) with $`\mu 3/2`$ for $`m=1`$ and $`\mu 2m`$ for large values of $`m`$. The latter corresponds to a strong influence of ballistic regime of tracer dynamics. We note that the value $`\beta 3/2`$ for $`\alpha 2`$ has also been observed for flows generated by three identical vortices KZ2000 and since this value remains for 3-vortex flows with extreme stress (vicinity of collapse) we may reasonably speculate that for all periodic (bounded) three vortex flows $`\beta 3/2`$. We would like to point out that the present work by analyzing the role played in transport by the different structures involved in the flow using various techniques, and by confirming a typical value of the second moment exponent should be of interest for the analysis of more realistic and complicated systems involving many vortices and coherent structures such as geophysical fluid dynamics. ## Acknowledgments We would like to thanks N. J. Zabusky for usefull discussions regarding the special influence of vortex collapse in passive tracers dynamics. This work was supported by the US Department of Navy, Grant No. N00014-96-1-0055, and the US Department of Energy, Grant No. DE-FG02-92ER54184. This research was supported in part by NSF cooperative agreement ACI-9619020 through computing resources provided by the National Partnership for Advanced Computational Infrastructure at the San Diego Supercomputer Center. ## Exponential Period Growth In this appendix, we recall some earlier results presented in LKZ2000 , and compute an asymptotic of the period growth as a function of $`\delta =1/2k`$. It has been shown for the case of three vortices with two identical ones that the relative motion of these vortices can been described using an one-dimensional effective Hamiltonian $$H_{eff}(\dot{X},X;\mathrm{\Lambda },K,k)P^2/2+V(X)=0,$$ (62) with Hamiltonian equations $$\dot{X}=H_{eff}/PP,\dot{P}=H_{eff}/X,$$ (63) where $`X=R_1^2`$ is the square of the distance between the two positive vortices, and the potential $`V`$ has the following form $$V(X)\frac{[(K(1k)X)^24k^2Y][(XK)^24k^2Y]}{8\pi ^2k^2Y^2},Y=(\mathrm{\Lambda }X)^{1/k}.$$ (64) Let us now estimate the period of the relative motion in our case. In this paper we choose a situation with $`K=0`$, using then the transformation $$U=X^2/4k^2Y,$$ (65) the potential (64) becomes, $$V(U)=\frac{2}{\pi ^2}k^2\left(1k\right)^2\left(U\frac{1}{\left(1k\right)^2}\right)\left(U1\right),$$ (66) given the fact that $`H_{eff}=0`$, the motion is confined to the negative regions of the potential which leads to $$1U\frac{1}{(1k)^2}.$$ (67) Consequently we obtain from (67) the boundaries for $`X`$ during the motion. We then compute the period of the relative motion using the effective potential $$T=2_{X_1}^{X_2}\frac{dX}{\sqrt{V(X)}}.$$ (68) In the limit $`\delta 0`$ ($`k1/2`$), we obtain $$V(U)\frac{1}{8\pi ^2}(U4)(U1),$$ (69) and using the inverse transformation, $$X=(4k^2\mathrm{\Lambda }^{1/k}U)^{1/(21/k)}(\mathrm{\Lambda }^2U)^{1/4\delta },$$ (70) we express the period in terms of $`U`$, which leads to $$T\frac{1}{4\delta }_1^4\frac{(\mathrm{\Lambda }^2U)^{11/4\delta }}{\sqrt{(U1)(4U)}}𝑑U.$$ (71) As $`\delta 0`$ ($`k1/2`$), it is the numerator in (71), which defines the asymptotic behavior, since it is decreasing function of $`U`$, the dominant term is from $`U=1`$ which leads to $$T\frac{1}{\delta }\mathrm{\Lambda }^{1/2\delta }.$$ (72) Since the condition $`\mathrm{\Lambda }=0.9<1`$, is verified (and is necessary for the collapse to happen LKZ2000 ), we have an exponential growth of the period as the vortex-collapse configuration is approached.
warning/0006/hep-ph0006029.html	ar5iv	text	# 1 Introduction ## 1 Introduction The evidence in favour of oscillations of solar and atmospheric neutrinos with square mass differences orders of magnitude smaller than the square masses of the other fermions strongly supports SO(10) gauge unification, where the see-saw mechanism naturally accounts for this large mass gap. In this framework it is important to establish the order of magnitude of the elements of the $`\nu _R`$’s Majorana mass matrix,which is related to the scale of the spontaneous symmetry breaking of the SO(10) generator, B - L . In the following section we shall show the uncertainties for the $`\nu _L`$ effective Majorana mass matrix,which follow not only from the presence of different solutions to the ”solar neutrino problem”, but also from not knowing the mass of the lightest neutrino mass eigenstate and the relative sign of the three mass eigenvalues. We shall also show that an expansion in the eigenvalues of the Dirac neutrino mass matrix, assumed to be hierarchically ordered according to the family as the other fermions, is unable to reproduce with the first two terms the phenomenological properties of the $`\nu _L`$ effective Majorana mass-matrix, if small mixing angles for the Dirac mass matrices are assumed for leptons in analogy with what we know about the CKM matrix for quarks. This fact motivates us to take in the third section a vanishing value for $`M_{33}^R`$, . Within the simplifying assumption of a trivial mixing matrix for the Dirac leptons, we are able to implement that condition for all the different solutions of the ”solar neutrino problem”. In most cases we find solutions with the two lower masses almost degenerate and at the order of magnitude of the highest one and with the highest matrix element of $`M^R`$,$`M_{23}^R`$, in the range $`10^{11}10^{12}`$ GeV in good agreement with the scale of the spontaneous breaking of B-L, $`2.710^{11}`$ GeV in the model with SU(4)xSU(2)xSU(2) intermediate symmetry.We find also solutions with the two lower masses at the order $`\sqrt{\mathrm{\Delta }m_{sun}^2}`$ and in that case the order of magnitude of the largest matrix element of $`M^R`$ depends on the solution to the ”solar neutrino problem” and on the value of $`U_{e3}`$. Finally we give our conclusions. ## 2 The effective $`\nu _L`$ Majorana mass matrix. The experimental information from solar and atmospheric neutrinos is such to ”almost” allow a model-independent determination of the neutrino mass matrix, within the hypothesis of the existence of only the three active Majorana neutrinos and CP symmetry. The second hypothesis is expected to be at most an approximation, since CP violation in the neutral kaon system most probably comes from a phase in the CKM matrix, as is confirmed by the indication of a large CP-violating asymmetry in the $`J/\psi K^S`$ channel, and a similar phase may be present in the lepton sector, for which, with Majorana neutrinos, two other CP-violating phases are also allowed. The incompleteness in the determination of the neutrino mass matrix arises from the following facts: * Up to now there is only an upper limit for $`\|U_{e3}\|^2(.05)`$ . * Oscillations depend on $`\mathrm{\Delta }m^2`$’s and cannot by themselves determine the masses of the neutrino eigenstates. * There are sign ambiguities in the mass-eigenvalues. Moreover there are five solutions of the “neutrino solar problem” with the oscillation parameters reported in Table I. \| Scenario \| $`\frac{\mathrm{\Delta }m_s^2}{eV^2}`$ \| $`sin^2(2\theta _s)`$ \| C.L. \| \| --- \| --- \| --- \| --- \| \| LMA \| 2.7 $`10^5`$ \| 0.79 \| 68 % \| \| SMA \| 5.0 $`10^6`$ \| 0.72 $`10^2`$ \| 64 % \| \| LOW \| 1.0 $`10^7`$ \| 0.91 \| 83% \| \| $`VAC_S`$ \| 6.5 $`10^{11}`$ \| 0.72 \| 90 % \| \| $`VAC_L`$ \| 4.4 $`10^{10}`$ \| 0.90 \| 95 % \| Table I For atmospheric neutrino oscillations the relevant parameters are : $$\mathrm{\Delta }m_a^23.510^3(eV)^2\mathrm{sin}^2(2\theta _a)1$$ (2.1) We shall use this information to determine the effective $`\nu _L`$ Majorana mass matrix. By assuming CP symmetry, which implies a real symmetric $`3\times 3`$ Majorana neutrino mass matrix $`M^L`$ : $$M^L=\left(m_{ij}\right)_{i,j=1,2,3}$$ $$=\mathrm{\Sigma }\|m_i>m_i<m_i\|$$ (2.2) where the eigenvalues $`m_i`$ are real, not necessarily positive, numbers. Our convention is $$\|m_3\|>\|m_2\|>\|m_1\|$$ (2.3) and we may take $`m_3>0`$, since an overall change of sign of $`M^L`$ has no physical consequences. We define $`\|m_3>`$ $`=`$ $`\mathrm{sin}\psi \|\nu _e>+\mathrm{cos}\psi (\mathrm{cos}\theta \|\nu _\mu >+\mathrm{sin}\theta \|\nu _\tau >)`$ $`\|m_2>`$ $`=`$ $`\mathrm{sin}\chi \|v_1>+\mathrm{cos}\chi \|v_2>`$ $`\|m_1>`$ $`=`$ $`\mathrm{cos}\chi \|v_1>+\mathrm{sin}\chi \|v_2>`$ (2.4) where $`\|v_1>`$ $`=`$ $`\mathrm{cos}\psi \|\nu _e>\mathrm{sin}\psi (\mathrm{cos}\theta \|\nu _\mu >+\mathrm{sin}\theta \|\nu _\tau >)`$ $`\|v_2>`$ $`=`$ $`\mathrm{sin}\theta \|\nu _\mu >+\mathrm{cos}\theta \|\nu _\tau >`$ (2.5) The difference $`m_3^2m_1^2`$ has to be identified with $`\mathrm{\Delta }m_a^2`$, while for $`\mathrm{\Delta }m_s^2`$ one has two options, either $`m_3^2m_2^2`$ or $`m_2^2m_1^2`$. The current interpretation of the atmospheric neutrino anomaly due to a practically maximal mixing between $`\nu _\mu `$ and $`\nu _\tau `$, with $`\nu _e`$ playing a marginal role, would imply with the first choice having $`\nu _e`$ components almost completely along the two nearly degenerate higher-mass eigenstates. Since this situation is unnatural in the $`SO(10)`$ framework, which we will describe in the second part of this paper,we take the option $`\mathrm{\Delta }m_s^2=m_2^2m_1^2`$. From eqs. 2.2, 2.4 and 2.5, it is easy to get: $`m_{11}`$ $`=`$ $`m_3\mathrm{sin}^2\psi +m_2\mathrm{sin}^2\chi \mathrm{cos}^2\psi +m_1\mathrm{cos}^2\chi \mathrm{cos}^2\psi `$ $`m_{22}`$ $`=`$ $`m_3\mathrm{cos}^2\theta \mathrm{cos}^2\psi +\mathrm{cos}^2\theta \mathrm{sin}^2\psi (m_2\mathrm{sin}^2\chi +m_1\mathrm{cos}^2\chi )`$ $`+`$ $`sin^2\theta (m_2\mathrm{cos}^2\chi +m_1\mathrm{sin}^2\chi )\mathrm{sin}2\theta \mathrm{sin}2\chi {\displaystyle \frac{\mathrm{sin}\psi }{2}}(m_2m_1)`$ $`m_{33}`$ $`=`$ $`m_3\mathrm{sin}^2\theta \mathrm{cos}^2\psi +sin^2\theta \mathrm{sin}^2\psi (m_2\mathrm{sin}^2\chi +m_1\mathrm{cos}^2\chi )`$ $`+`$ $`\mathrm{cos}^2\theta (m_2\mathrm{cos}^2\chi +m_1\mathrm{sin}^2\chi )+\mathrm{sin}2\theta \mathrm{sin}2\chi {\displaystyle \frac{\mathrm{sin}\psi }{2}}(m_2m_1)`$ $`m_{12}`$ $`=`$ $`\mathrm{cos}\psi [\mathrm{sin}\psi \mathrm{cos}\theta (m_3m_2\mathrm{sin}^2\chi m_1\mathrm{cos}^2\chi )`$ $`+\mathrm{sin}\theta {\displaystyle \frac{\mathrm{sin}2\chi }{2}}(m_2m_1)]`$ $`m_{13}`$ $`=`$ $`\mathrm{cos}\psi [\mathrm{sin}\psi \mathrm{sin}\theta (m_3m_2\mathrm{sin}^2\chi m_1\mathrm{cos}^2\chi )`$ $`\mathrm{cos}\theta {\displaystyle \frac{\mathrm{sin}\mathrm{\hspace{0.17em}2}\chi }{2}}(m_2m_1)]`$ $`m_{23}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}\mathrm{\hspace{0.33em}2}\theta }{2}}[m_3\mathrm{cos}^2\psi m_2\mathrm{cos}^2\chi m_1\mathrm{sin}^2\chi `$ $`+\mathrm{sin}^2\psi (m_2\mathrm{sin}^2\chi +m_1\mathrm{cos}^2\chi )]`$ $`+`$ $`\mathrm{cos}2\theta {\displaystyle \frac{\mathrm{sin}2\chi }{2}}\mathrm{sin}\psi (m_2m_1)`$ The angles $`\theta `$ and $`\chi `$ should be identified with the mixing angles for atmospheric and solar neutrino oscillations, respectively. The experimental value found for the atmospheric neutrino oscillations, $`\mathrm{sin}^22\theta _a1`$, and the upper limit on $`\|U_{e3}\|^2`$ imply $`\theta =\frac{\pi }{4}`$ and a small value for $`\mathrm{sin}\psi `$. With hierarchical neutrino masses $$m_3>>\|m_2\|>>\|m_1\|$$ (2.7) one has $`m_3^2`$ $``$ $`\mathrm{\Delta }m_a^2`$ $`m_2^2`$ $``$ $`\mathrm{\Delta }m_s^2`$ (2.8) and by taking $`\psi =0`$ for the contributions proportional to $`m_2`$ and neglecting $`\|m_1\|<<\sqrt{\mathrm{\Delta }m_s^2}`$, one gets: $`M^L`$ $`=`$ $`{\displaystyle \frac{m_3}{2}}\mathrm{cos}^2\psi \left(\begin{array}{ccc}2\mathrm{tan}\psi ^2& \sqrt{2}\mathrm{tan}\psi & \sqrt{2}\mathrm{tan}\psi \\ \sqrt{2}\mathrm{tan}\psi & 1& 1\\ \sqrt{2}\mathrm{tan}\psi & 1& 1,\end{array}\right)`$ (2.12) $`+`$ $`{\displaystyle \frac{m_2}{2}}\mathrm{cos}^2\chi \left(\begin{array}{ccc}2\mathrm{tan}^2\chi & \sqrt{2}\mathrm{tan}\chi & \sqrt{2}\mathrm{tan}\chi \\ \sqrt{2}\mathrm{tan}\chi & 1& 1\\ \sqrt{2}\mathrm{tan}\chi & 1& 1,\end{array}\right).`$ (2.16) We now consider the possibility, often advocated, that the two lower masses are almost degenerate and larger than $`\sqrt{\mathrm{\Delta }m_s^2}`$, but smaller than $`\sqrt{\mathrm{\Delta }m_a^2}`$ . In that case, one has still $`\mathrm{\Delta }m_a^2=m_3^2`$, but we can neglect $`\mathrm{\Delta }m_{sun}^2`$ with respect to $`m_1^2`$ and $`m_2^2`$. The first term in eq. 2.16 remains the same, while the other one takes a different form, depending on the relative sign of $`m_2`$ and $`m_1`$.We will consider $`m_1m_2<0`$, since this assumption will be necessary in the following and get $$\frac{m_2\mathrm{cos}2\chi }{2}\left(\begin{array}{ccc}1& \sqrt{2}\mathrm{tan}2\chi & \sqrt{2}\mathrm{tan}2\chi \\ \sqrt{2}\mathrm{tan}2\chi & 1& 1\\ \sqrt{2}\mathrm{tan}2\chi & 1& 1,\end{array}\right)$$ (2.17) Again for the term proportional to $`m_2`$ in eq. 2.17 we made the approximation to take $`\psi =0`$. We now consider the possibility advocated by Georgi and Glashow , with the motivation of having a sizeable neutrino contribution to the hot dark matter, of almost degenerate square masses for the neutrinos, larger than their differences. We have four options, according to the relative signs of the $`m_i`$’s. By taking again $`m_1<0<m_2`$ and $`\psi =0`$ for the part proportional to $`\mathrm{\Delta }m_a^2`$, we get : $`M^L`$ $`=`$ $`m_3\mathrm{cos}^2\chi \left(\begin{array}{ccc}\mathrm{tan}^2\chi 1& \sqrt{2}\mathrm{cos}\psi \mathrm{tan}\chi & \sqrt{2}\mathrm{cos}\psi \mathrm{tan}\chi \\ \sqrt{2}\mathrm{cos}\psi \mathrm{tan}\chi & 1& \mathrm{tan}^2\chi \\ \sqrt{2}\mathrm{cos}\psi \mathrm{tan}\chi & \mathrm{tan}^2\chi & 1,\end{array}\right)`$ (2.21) $`+`$ $`m_3\mathrm{cos}^2\chi \mathrm{sin}\psi \left(\begin{array}{ccc}0& \sqrt{2}\mathrm{cos}\psi & \sqrt{2}\mathrm{cos}\psi \\ \sqrt{2}\mathrm{cos}\psi & 2\mathrm{tan}\chi & 0\\ \sqrt{2}\mathrm{cos}\psi & 0& 2\mathrm{tan}\chi ,\end{array}\right)`$ (2.25) $`+`$ $`2m_3\mathrm{cos}^2\chi sin^2\psi \left(\begin{array}{ccc}1& 0& 0\\ 0& \frac{1}{2}& \frac{1}{2}\\ 0& \frac{1}{2}& \frac{1}{2},\end{array}\right)`$ (2.29) $`+`$ $`{\displaystyle \frac{\mathrm{\Delta }m_a^2\mathrm{cos}2\chi }{4m_3}}\left(\begin{array}{ccc}2& \sqrt{2}\mathrm{tan}2\chi & \sqrt{2}\mathrm{tan}2\chi \\ \sqrt{2}\mathrm{tan}2\chi & 1& 1\\ \sqrt{2}\mathrm{tan}2\chi & 1& 1,\end{array}\right).`$ (2.33) Let us consider $`m_{\nu _e\nu _e}`$ $`=`$ $`m_3(\mathrm{cos}2\chi +2\mathrm{cos}^2\chi \mathrm{sin}^2\psi `$ (2.34) $`+`$ $`\mathrm{sin}^2\chi {\displaystyle \frac{\mathrm{\Delta }m_s^2}{2m_3^2}}+{\displaystyle \frac{\mathrm{\Delta }m_a^2}{2m_3^2}}\mathrm{cos}^2\chi )`$ for which there is an experimental upper limit coming from the study of neutrinoless double beta-decay $`\|m_{\nu _e\nu _e}\|<.2`$ $`eV`$ . To get the cancellation between $`\mathrm{cos}2\chi `$ and $`2\mathrm{cos}^2\chi \mathrm{sin}2\psi `$, $`\chi `$ should be near to $`\frac{\pi }{4}`$, more precisely : $$\mathrm{sin}^2\psi =\frac{\mathrm{tan}^2\chi 1}{2}$$ (2.35) But the r.h.s. of eq. 2.35 takes at least the central value .22 (for $`\mathrm{sin}^22\theta _s=.91`$), larger than the upper limit for $`\mathrm{sin}^2\psi `$, which implies $$\|m_3\|\|\frac{m_{\nu _e\nu _e}}{2\mathrm{cos}^2\chi \mathrm{sin}^2\psi \mathrm{cos}2\chi }\|.85eV$$ (2.36) The situation is similar for the case with $`m_1>0>m_2`$. In conclusion many options are open for the neutrino mass matrix and even an exact determination of the mixing matrix and of the square mass differences is unable to resolve the ambiguity associated to the value of $`\|m_1\|`$ and to the relative signs of the masses .Also the same Majorana mass matrix for the $`\nu _L`$’s may be obtained with different choices of the Dirac lepton matrices and of $`M^R`$. It has been suggested since a long time that the peculiar properties of neutrino mixing, almost maximal for atmospheric, rather large for all the solutions to the “solar neutrino problem”, except the MSW small angle solution, should be related to the fact that, differently from the case for the charged fermions, their masses arise from the see-saw mechanism .The later has the indisputable merit of providing a reason for the small value of the neutrino masses, especially in the framework of $`SO(10)`$ unified theories, which predict the existence of left-handed antineutrinos with Majorana masses related to the spontaneous breaking of $`BL`$ symmetry . Unified $`SO(10)`$ theories are suitable for the study of fermion masses as all the fermions of one family belong to a single representation, the spinorial 16 . Indeed, by classifying the Higgs doublet responsible for the breaking of the electroweak symmetry in the vector representation (10), one obtains,together with the celebrated( but actually not so successfull) relationship , $$\frac{m_\tau }{m_b}(unification)=1$$ (2.37) while the analogous relationship, $`\frac{m_{\nu _\tau }}{m_t}(unification)=1`$, is turned by the see-saw mechanism into the intriguing prediction of very small neutrino masses. To correct $`\frac{m_\mu }{m_s}=\frac{m_e}{m_d}=1`$, some component of the electroweak Higgs, at least along the 126, should be introduced $`(16\times 16)_S`$ $`=`$ $`126+10`$ $`(16\times 16)_A`$ $`=`$ $`120`$ (2.38) Without adopting a particular scheme, we shall limit ourselves to assume that the Dirac neutrino mass matrix, once diagonalized by the biunitary transformation, gives rise to a hierarchical relationship for the elements of the diagonal matrix similar to that existing for the other fermions, and the matrix corresponding to the CKM, which gives the misalignment with the corresponding leptons, has small non diagonal matrix elements. $`M^L`$ $`=`$ $`m^D(M^R)^1(m^D)^T`$ $`m^D`$ $`=`$ $`U_Lm_{diag}U_R^+`$ (2.39) with $$m_{diag}=diag(\mu _1,\mu _2,\mu _3)$$ (2.40) and we assume $`\mu _1<<\mu _2<<\mu _3`$. From eqs. 2.39 it is easy to derive: $$M^L=U_Lm_{diag}U_R^+(M^R)^1U_R^{}m_{diag}U_L^T$$ (2.41) which is simplified by defining $$\stackrel{~}{N}=U_R^+M^{R^1}U_R^{}$$ (2.42) into $$M^L=U_Lm_{diag}\stackrel{~}{N}m_{diag}U_L^T$$ (2.43) With our hypothesis of CP conservation, $`U_R`$ and $`U_L`$ are, in fact, real orthogonal matrices and $`\stackrel{~}{N}`$ is symmetric real. We may develope $`M^L`$ as a quadratic form in the $`\mu _i`$’s: $`M_{\mathrm{}\mathrm{}^{}}^L`$ $`=`$ $`(U_L)_\mathrm{}3(U_L)_\mathrm{}^{}3\stackrel{~}{N}_{33}\mu _3^2`$ (2.44) $`+`$ $`((U_L)_\mathrm{}2(U_L)_\mathrm{}^{}3+(U_L)_\mathrm{}3(U_L)_\mathrm{}^{}2)\stackrel{~}{N}_{23}\mu _2\mu _3`$ $`+`$ $`(U_L)_\mathrm{}2(U_L)_\mathrm{}^{}2\stackrel{~}{N}_{22}\mu _2^2+\mathrm{}`$ Should the matrix elements of $`\stackrel{~}{N}`$ be of the same order, it would be a reasonable approximation to take only the terms proportional to $`\mu _3^2`$, which would give for $`M^L`$ the expression: $$\stackrel{~}{N}_{33}\mu _3^2\left(\begin{array}{cccc}(U_L)_{13}^2& (U_L)_{13}(U_L)_{23}& (U_L)_{13}(U_L)_{33}& \\ (U_L)_{13}(U_L)_{23}& (U_L)_{23}^2& (U_L)_{23}(U_L)_{33}\\ (U_L)_{13}(U_L)_{33}& (U_L)_{23}(U_L)_{33}& (U_L)_{33}^2\end{array}\right)$$ (2.45) The matrix defined by eq. 2.45 should have eigenvalues $`\mu _3^2\stackrel{~}{N}_{33}`$, corresponding to the eigenvector $$\left(\begin{array}{c}U_{13}\\ U_{23}\\ U_{33}\end{array}\right),$$ (2.46) and twice the eigenvalue 0. The eigenvector in eq. 2.46 should be identified with: $$\left(\begin{array}{c}\mathrm{sin}\psi \\ \mathrm{cos}\psi \mathrm{cos}\theta \\ \mathrm{cos}\psi \mathrm{sin}\theta \end{array}\right),$$ (2.47) which, for $`\theta =\frac{\pi }{4}`$, has almost equal second and third components, implying $`U_{23}U_{33}`$, in disagreement with our assumption that the mixing angles of Dirac neutrinos with the corresponding leptons are small. Let us consider also the term proportional to $`\mu _2\mu _3`$, and, for brevity, define: $`M`$ $`=`$ $`\stackrel{~}{N}_{33}\mu _3^2`$ $`\stackrel{~}{M}`$ $`=`$ $`\stackrel{~}{N}_{23}\mu _2\mu _3.`$ (2.48) In such a case, the neutrino mass matrix would be: $`M`$ $`\left(\begin{array}{cccc}(U_L)_{13}^2& (U_L)_{13}(U_L)_{23}& (U_L)_{13}(U_L)_{33}& \\ (U_L)_{13}(U_L)_{23}& (U_L)_{23}^2& (U_L)_{23}(U_L)_{33}\\ (U_L)_{13}(U_L)_{33}& (U_L)_{23}(U_L)_{33}& (U_L)_{33}^2\end{array}\right)`$ $`+`$ $`\stackrel{~}{M}`$ $`\left(\begin{array}{ccc}2(U_L)_{13}(U_L)_{12}& (U_L)_{13}(U_L)_{22}+(U_L)_{12}(U_L)_{23}& (U_L)_{12}(U_L)_{33}+(U_L)_{13}(U_L)_{32}\\ (U_L)_{13}(U_L)_{22}+(U_L)_{12}(U_L)_{23}& 2(U_L)_{22}(U_L)_{23}& (U_L)_{22}(U_L)_{33}+(U_L)_{23}(U_L)_{32}\\ (U_L)_{12}(U_L)_{33}+(U_L)_{13}(U_L)_{32}& (U_L)_{22}(U_L)_{33}+(U_L)_{23}(U_L)_{32}& 2(U_L)_{32}(U_L)_{33}\end{array}\right)`$ which has a vanishing eigenvalue and the other two given by: $$\frac{M\pm \sqrt{M^2+4\stackrel{~}{M}^2}}{2}$$ (2.50) By taking $`M`$ positive, one should identify $$\mathrm{\Delta }m_a^2=\frac{1}{4}[M+\sqrt{M^2+4\stackrel{~}{M}^2}]^2,$$ (2.51) while for $`\mathrm{\Delta }m_s^2`$ ($`<<`$ $`\mathrm{\Delta }m_a^2`$ ) one should have $$\mathrm{\Delta }m_s^2=\frac{1}{4}[\sqrt{M^2+4\stackrel{~}{M}^2}M]^2$$ (2.52) $`{\displaystyle \frac{\mathrm{\Delta }m_s^2}{\mathrm{\Delta }m_a^2}}{\displaystyle \frac{\mathrm{\Delta }m_s^2\mathrm{\Delta }m_a^2}{(\mathrm{\Delta }m_a^2)^2+(\mathrm{\Delta }m_s^2)^2}}=`$ $`{\displaystyle \frac{2\stackrel{~}{M}^4}{2M^4+8\stackrel{~}{M}^2M^2+4\stackrel{~}{M}^4}}\left({\displaystyle \frac{\stackrel{~}{M}}{M}}\right)^4`$ (2.53) In the most favourable case, large angle MSW, it would give: $$\left(\frac{\stackrel{~}{M}}{M}\right)=\left(\frac{1.610^5}{3.510^3}\right)^{\frac{1}{4}}\frac{1}{3.8}$$ (2.54) Let us see how large the non-diagonal matrix elements of $`U_L`$ should be in order to give rise to a neutrino mass matrix, which has the property of having almost equal matrix elements in $`22`$ and $`33`$ positions.We should have $$M(U_L)_{23}^2+2\stackrel{~}{M}(U_L)_{23}(U_L)_{22}=M(U_L)_{33}^2+2\stackrel{~}{M}(U_L)_{33}(U_L)_{32}$$ (2.55) which is impossible to obtain with small values of the non-diagonal matrix elements, since $$(U_L)_{23}(M(U_L)_{23}+2\stackrel{~}{M}(U_L)_{22})<<(U_L)_{33}(M(U_L)_{33}+2\stackrel{~}{M}(U_L)_{32})$$ (2.56) So the expansion with only the two terms is not able to reproduce a neutrino mass matrix consistent with the experimental information. Therefore, we must now consider further terms, i.e., those proportional to $`\mu _2^2`$ and to $`\mu _1\mu _3`$. If we define: $`\sigma `$ $`=`$ $`\stackrel{~}{N}_{22}\mu _2^2`$ $`\tau `$ $`=`$ $`\stackrel{~}{N}_{13}\mu _1\mu _3`$ (2.57) in the limit $`(U_L)_{ij}=\delta _{ij}`$, we would just find the matrix proposed by Stech $$\left(\begin{array}{ccc}0& 0& \tau \\ 0& \sigma & \mu \\ \tau & \mu & M\end{array}\right)$$ (2.58) which, to reproduce the maximal mixing for $`\nu _\mu `$ and $`\nu _\tau `$, requires the equality of the terms proportional to $`\mu _2^2`$ and $`\mu _3^2`$, respectively. From this discussion it follows that the term proportional to $`\mu _3^2`$ does not play a dominant role in the effective $`\nu _L`$ Majorana mass matrix not only with respect to the one proportional to $`\mu _2\mu _3`$, but also with respect to the one proportional to $`\mu _2^2`$.For a diagonal lepton Dirac matrix the term $`\mu _3^2`$ appears in the expression of $`M_{33}^R`$, multiplied by$`(M_{}^{L}{}_{}{}^{1})_{33}`$. So to reduce the contribution of $`\mu _3^2`$ we take a small value for $`(M_{}^{L}{}_{}{}^{1})_{33}`$. In the following section we shall impose the vanishing of $`(M_{}^{L}{}_{}{}^{1})_{33}`$ and consequently of $`(M^R)_{33}`$ with a diagonal neutrino Dirac matrix given by: $$m_\nu ^D=\frac{m_\tau }{m_b}diag(m_u,m_c,m_t)$$ (2.59) ## 3 See-Saw model in SO(10) with diagonal neutrino Dirac matrix and $`M_{33}^R=0`$ We have seen in the previous section that $`M^L`$, derived from a Dirac neutrino mass matrix with the properties of the quark mass matrix, hierarchical relationship for its eigenvalues and small mixing angles, is not dominated by the term proportional to $`(m_{\nu _\tau }^D)^2`$,which is as important as the term proportional to $`(m_{\nu _\mu }^D)^2`$.A way to implement this property, the non-dominance of the term proportional to $`(m_{\nu _\tau }^D)^2`$, is found by considering the inverse of the first of eqs. 2.39: $$M^R=m_\nu ^D(M^L)^1m_\nu ^D.$$ (3.1) with $`m_\nu ^D`$ given by eq. 2.59: $`({\displaystyle \frac{m_b}{m_\tau m_u}})^2(M^R)_{11}`$ $`=`$ $`{\displaystyle \frac{1}{m_1}}\mathrm{cos}^2\psi \mathrm{cos}^2\chi {\displaystyle \frac{1}{m_2}}\mathrm{cos}^2\psi \mathrm{sin}^2\chi {\displaystyle \frac{1}{m_3}}\mathrm{sin}^2\psi `$ $`{\displaystyle \frac{(\frac{m_b}{m_\tau })^2}{m_um_c}}(M^R)_{12}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}m_1}}\mathrm{cos}\psi \mathrm{cos}\chi (\mathrm{cos}\chi \mathrm{sin}\psi +\mathrm{sin}\chi )`$ $`{\displaystyle \frac{1}{\sqrt{2}m_2}}\mathrm{cos}\psi \mathrm{sin}\chi (\mathrm{cos}\chi \mathrm{sin}\psi \mathrm{sin}\chi ){\displaystyle \frac{1}{2\sqrt{2}m_3}}\mathrm{sin}(2\psi )`$ $`{\displaystyle \frac{(\frac{m_b}{m_\tau })^2}{m_um_t}}(M^R)_{13}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}m_1}}\mathrm{cos}\psi \mathrm{cos}\chi (\mathrm{cos}\chi \mathrm{sin}\psi \mathrm{sin}\chi )`$ $`+{\displaystyle \frac{1}{\sqrt{2}m_2}}\mathrm{cos}\psi \mathrm{sin}\chi (\mathrm{cos}\chi +\mathrm{sin}\psi \mathrm{sin}\chi ){\displaystyle \frac{1}{2\sqrt{2}m_3}}\mathrm{sin}(2\psi )`$ $`({\displaystyle \frac{m_b}{m_\tau m_c}})^2(M^R)_{22}`$ $`=`$ $`{\displaystyle \frac{1}{2m_1}}(\mathrm{cos}\chi \mathrm{sin}\psi +\mathrm{sin}\chi )^2`$ $`{\displaystyle \frac{1}{2m_2}}(\mathrm{cos}\chi \mathrm{sin}\psi \mathrm{sin}\chi )^2{\displaystyle \frac{1}{2m_3}}\mathrm{cos}^2\psi `$ $`{\displaystyle \frac{(\frac{m_b}{m_\tau })^2}{m_cm_t}}(M^R)_{23}`$ $`=`$ $`{\displaystyle \frac{1}{2m_1}}(\mathrm{sin}^2\chi \mathrm{cos}^2\chi \mathrm{sin}^2\psi )`$ $`+{\displaystyle \frac{1}{2m_2}}(\mathrm{cos}^2\chi \mathrm{sin}^2\psi \mathrm{sin}^2\chi ){\displaystyle \frac{1}{2m_3}}\mathrm{cos}^2\psi `$ $`({\displaystyle \frac{m_b}{m_\tau m_t}})^2(M^R)_{33}`$ $`=`$ $`{\displaystyle \frac{1}{2m_1}}(\mathrm{sin}\chi \mathrm{cos}\chi \mathrm{sin}\psi )^2`$ $`{\displaystyle \frac{1}{2m_2}}(\mathrm{cos}\chi +\mathrm{sin}\chi \mathrm{sin}\psi )^2{\displaystyle \frac{1}{2m_3}}\mathrm{cos}^2\psi `$ From eqs. 3, it is easy to see that $$\|(M^R)_{11}\|<\frac{(\frac{m_\tau m_u}{m_b})^2}{\|m_1\|}$$ (3.3) which implies , in the absence of cancellations for $`(M^R)_{33}`$ <sup>1</sup><sup>1</sup>1Since the ratios of the quark masses have a negligeable dependence on the scale, we can take the values given in : $$\|\frac{(M^R)_{11}}{(M^R)_{33}}\|2\frac{m_u^2}{m_t^2}2\left(\frac{\frac{8}{3}\mathrm{MeV}}{180\mathrm{G}\mathrm{e}\mathrm{V}}\right)^2410^{10}$$ (3.4) To prevent the appearence of such unnatural small factor we impose the vanishing of the r.h.s of last equation in 3: $$\frac{(\mathrm{sin}\chi \mathrm{cos}\chi \mathrm{sin}\psi )^2}{m_1}+\frac{(\mathrm{cos}\chi +\mathrm{sin}\psi \mathrm{sin}\chi )^2}{m_2}+\frac{\mathrm{cos}^2\psi }{m_3}=0$$ (3.5) which requires that not all the $`m_i`$’s have the same sign and $`{\displaystyle \frac{m_b^2}{m_\tau ^2m_cm_t}}(M^R)_{23}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\chi \mathrm{sin}\psi }{2}}({\displaystyle \frac{1}{m_1}}{\displaystyle \frac{1}{m_2}})`$ (3.6) $`(\mathrm{sin}\psi )^2({\displaystyle \frac{(\mathrm{cos}\chi )^2}{m_1}}+{\displaystyle \frac{(\mathrm{sin}\chi )^2}{m_2}}){\displaystyle \frac{(\mathrm{cos}\psi )^2}{m_3}}`$ By neglecting the terms proportional to powers of $`\mathrm{sin}\psi `$, we get $$\|(M^R)_{23}\|710^{11}\mathrm{GeV}$$ (3.7) as in . With $`\mathrm{sin}\psi =0`$, eqs. 3 read: $`({\displaystyle \frac{m_b}{m_\tau m_u}})^2(M^R)_{11}`$ $`=`$ $`{\displaystyle \frac{1}{m_1}}\mathrm{cos}^2\chi `$ $`{\displaystyle \frac{(\frac{m_b}{m_\tau })^2}{m_um_c}}(M^R)_{12}`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}\chi \mathrm{sin}\chi }{\sqrt{2}}}({\displaystyle \frac{1}{m_1}}{\displaystyle \frac{1}{m_2}})`$ $`{\displaystyle \frac{(\frac{m_b}{m_\tau })^2}{m_um_t}}(M^R)_{13}`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}\chi \mathrm{sin}\chi }{\sqrt{2}}}({\displaystyle \frac{1}{m_1}}+{\displaystyle \frac{1}{m_2}})`$ (3.8) $`({\displaystyle \frac{m_b}{m_\tau m_c}})^2(M^R)_{22}`$ $`=`$ $`{\displaystyle \frac{1}{2m_1}}\mathrm{sin}^2\chi {\displaystyle \frac{1}{2m_2}}\mathrm{cos}^2\chi {\displaystyle \frac{1}{2m_3}}`$ $`{\displaystyle \frac{(\frac{m_b}{m_\tau })^2}{m_cm_t}}(M^R)_{23}`$ $`=`$ $`{\displaystyle \frac{1}{2m_1}}\mathrm{sin}^2\chi +{\displaystyle \frac{1}{2m_2}}\mathrm{cos}^2\chi {\displaystyle \frac{1}{2m_3}}`$ $`({\displaystyle \frac{m_b}{m_\tau m_t}})^2(M^R)_{33}`$ $`=`$ $`{\displaystyle \frac{1}{2m_1}}\mathrm{sin}^2\chi {\displaystyle \frac{1}{2m_2}}\mathrm{cos}^2\chi {\displaystyle \frac{1}{2m_3}}`$ In order to get $`(M^R)_{33}=\mathrm{\hspace{0.17em}0}`$, one should have $$\frac{\mathrm{sin}^2\chi }{m_1}+\frac{\mathrm{cos}^2\chi }{m_2}+\frac{1}{m_3}=0$$ (3.9) We look for its solutions with $`m_i`$ ’s constrained by the relations: $`m_2^2`$ $`=`$ $`m_1^2+\mathrm{\Delta }m_s^2`$ $`m_3^2`$ $`=`$ $`m_1^2+\mathrm{\Delta }m_a^2`$ (3.10) With this we do not want to impose that $`(M^R)_{33}`$ vanishes exactly. But, once a solution with no large differences between the matrix elements of $`M^R`$ has been found, we can allow for a value of $`(M^R)_{33}`$ of the same order of the other matrix elements, with a small change for the left-handed neutrino Majorana mass matrix. It is obvious that to obey eq. 3.9 the $`m_i`$ ’s cannot have the same sign and, as $`\|m_3\|>\|m_2\|>\|m_1\|`$, $`m_1`$ and $`m_2`$ should have opposite signs.This fact has been pointed out in . From eq. 3.9 we get: $$\mathrm{tan}^2\chi =\frac{m_1(m_3+m_2)}{m_2(m_3+m_1)}$$ (3.11) which,when $`\|m_2\|<<\|m_3\|`$ simplifies to $$\mathrm{tan}^2\chi =\frac{m_1}{m_2}$$ (3.12) in agreement with the expectation that large mixing angles correspond to almost degenerate masses of the mixed states . One has another approximate solution with almost degenerate lightest neutrino mass eigenstates ($`\|m_2\|m_1=m`$) $$m=\mathrm{cos}(2\chi )m_3$$ (3.13) These approximate solutions correspond to exact solutions of eq. 3.9. In fact, with equal signs for $`m_2`$ and $`m_3`$ we may write eq. 3.9 in the form $$\mathrm{sin}\chi ^2+\mathrm{cos}\chi ^2\frac{m_1}{m_2}+\frac{m_1}{m_3}=0$$ (3.14) which has one and only one solution for negative $`m_1`$ since its l.h.s. is an increasing function of $`m_1`$ in the range $`(\mathrm{},0)`$ varying from $`2\mathrm{cos}^2\chi `$ to $`\mathrm{sin}^2\chi `$ <sup>2</sup><sup>2</sup>2Notice that eqs. 3.10 imply that when $`\left\|m_1\right\|`$ goes to $`\mathrm{}`$ the ratios $`\left\|\frac{m_1}{m_2}\right\|`$ and $`\left\|\frac{m_1}{m_3}\right\|`$ go to 1.. In the case of the same sign for $`m_1`$ and $`m_3`$ it is convenient to rewrite eq. 3.9 in the form $$\mathrm{sin}^2\chi \frac{m_2}{m_1}+\mathrm{cos}^2\chi +\frac{m_2}{m_3}=0$$ (3.15) With its l.h.s., in the range ($`\mathrm{}`$, $`\sqrt{\mathrm{\Delta }m_s^2})`$ for $`m_2`$, going from $`2(\mathrm{sin}^2\chi )`$ to $`\mathrm{}`$. But, if $`\mathrm{cos}^2\chi `$ is sufficiently larger than $`\mathrm{sin}^2\chi `$, can reach positive values, implying two solutions for eq. 3.9, approximately given by eqs. 3.12 and 3.13, respectively. For the MSW solutions one has to take $`\mathrm{cos}^2\chi >\mathrm{sin}^2\chi `$ and therefore one has three solutions. For the vacuum solutions, which allows both signs for $`\mathrm{cos}(2\chi )`$, one has also a solution with $`\mathrm{sin}^2\chi >\mathrm{cos}^2\chi `$. We wish to extend the previous analysis to consider non vanishing values of $`\|\mathrm{sin}\psi \|`$ . In that case, when $$(\mathrm{sin}\chi \mathrm{sin}\psi \mathrm{cos}\chi )^2>(\mathrm{cos}\chi +\mathrm{sin}\psi \mathrm{sin}\chi )^2$$ (3.16) one has only one solution of eq. 3.9; instead, when the r.h.s. of the inequality 3.16 is sufficiently larger than the l.h.s. + 1, there are two solutions. We have performed a numerical study of the solutions of eq. 3.6 with $`\|\mathrm{sin}\psi \|=k(.075)`$, k=0,…,3. To get a general view of our solutions, we give in Table II, for each solar neutrino oscillation scenario, the matrix where the smallest value for the ratios of the moduli of non-vanishing matrix elements of $`M^R`$ (which is, with the only exeption of VACL solution, $`\frac{M_{11}^R}{M_{23}^R}`$) takes the largest value and the ones where the max (ie. highest) matrix element of $`M^R`$ takes the smallest or largest value. From Table II we see that the matrices with smallest excursion for their matrix element correspond to small values for $`m_{1}^{}{}_{}{}^{2}`$ and to $`\psi =0`$, with the only exception of the LMW solution, where $`\|\mathrm{sin}\psi \|=0.075`$. The matrices with the smallest value of $`M_{23}^R`$ are found for the large $`m_{1}^{}{}_{}{}^{2}`$ solutions. Finally the largest values of $`M_{23}^R`$ are found at the boundary of the allowed values for $`\psi `$ and in correspondence of the small $`m_{1}^{}{}_{}{}^{2}`$ solutions. Since our motivation for setting $`M_{33}^R=0`$ is also based on the demand of dealing with $`M^R`$ matrices with not so disparate orders of magnitude for its matrix elements, we favour small $`\psi `$ and $`m_{1}^{}{}_{}{}^{2}`$ choices (in particular, the small MSW solution with the ratio between the smallest and the largest matrix elements of $`M^R`$ given by $`4.310^4`$ ). We never find neutrino masses large enough to be of cosmological relevance. The maximum value found for $`_i\left\|m_i\right\|`$ being .6 eV, corresponding to the boundary of allowed values of the small angle MSW solution. ## 4 Conclusions The hypothesis of a vanishing $`M_{33}^R`$, assumed in and here motivated by the requirement of not having too large fine tunings in the see-saw formula may be implemented in all five solutions to the ”solar neutrino problem”.In general one has solutions with large $`(\mathrm{\Delta }m_a^2`$) or small $`(\mathrm{\Delta }m_s^2)`$ for the square mass of the lightest neutrino mass eigenstate. For large values of $`m_1^2`$ the max matrix element of $`M^R`$ ($`M_{23}^R`$) has a value $`710^{11}`$ GeV, a scale found in a recent work and in agreement with the spontaneous scale of B-L symmetry breaking found in a SO(10) model with SU(4) x SU(2) x SU(2) intermediate symmetry . Moderate values for the matrix elements of $`M^R`$ have been also found in and recently by Akhmedov, Branco and Rebelo . For small values for $`m_1^2`$ one still finds the same order of magnitude for the large angle MSW solution in the entire range of the allowed values for $`\psi `$; the same happens for the other MSW solutions for $`\psi `$ = 0. The other MSW solutions correspond to larger values for the scales, especially for the largest allowed values of $`\|\mathrm{sin}\psi \|`$. The solutions corresponding to small $`m_1^2`$ at $`\psi `$=0 have the appealing feature of having not too different orders of magnitude for the matrix elements of $`M^R`$. Acknowledgements: One of us (F.B.) acknowledges G. Fiorentini for interesting discussions. Table II The matrices $`M_R`$ obeying eq. 3.6 for the five scenarios. For each solution the corresponding values of $`\left\|\mathrm{sin}\psi \right\|`$ and $`\left\|m_1\right\|`$ are given. As $`M_R`$ is symmetric, repeated matrix elements are denoted by #.
warning/0006/astro-ph0006168.html	ar5iv	text	# Exploring the gravitationally lensed system HE 1104-1805: Near-IR SpectroscopyBased on observations collected with the ESO New Technology Telescope (program 61.B-0413) ## 1 Introduction HE 1104$``$1805 is one object in the growing list of gravitationally lensed quasars which might be used to constrain cosmological parameters. It was discovered in the framework of the Hamburg/ESO Quasar Survey and consists of 2 lensed images of a radio-quiet quasar (RQQ) at $`z`$= 2.319, separated by $``$3.15″ (Wisotzki et al. 1993). The lensing galaxy was discovered from ground based near-IR (Courbin, Lidman & Magain, 1998; hereafter C98) and HST optical observations (Remy et al. 1998; hereafter R98). The relatively wide angular separation between the quasar images makes this object suitable for photometric monitoring programs, as conducted at ESO by Wisotzki et al. (1998; hereafter W98). From light curves measured over a period of 6 years, they derived a time delay for HE 1104$``$1805 of $`\mathrm{\Delta }t=0.73`$ years, with a second possible value of 0.3 years. Although we show that the complex lensing potential involved in HE 1104$``$1805 makes it difficult to determine H<sub>0</sub> from the time delay, we also show that HE 1104$``$1805 is probably much more of interest for microlensing studies, provided the lens redshift is better known. The present paper describes an attempt to measure the redshift of the main lensing galaxy from near-IR spectroscopy. Our near-IR observations where motivated by the very red colors measured for the lensing galaxy (C98, R98), and by the better contrast between the lens and the quasar in the near-IR. Although we were unsuccessful in measuring the lens redshift accurately, we did obtain high S/N spectra of the lensed quasar, between 0.95 and 2.5 microns. ## 2 Observations-Reductions The data were taken with SOFI, the near-IR (1 to 2.5 $`\mu `$) imaging spectrograph on the ESO NTT. Two grisms were used to cover the 1 to 2.5 $`\mu `$ wavelength range: a “blue grism”, which covers the region from 0.95 to 1.64 $`\mu `$ and a “red grism” which covers the region 1.53 to 2.52 $`\mu `$. With a 1″ slit, the spectral resolution is around 600. The observations with the blue grism were taken on the night of 1998 June 13, for a total integration time of 4560 seconds and the observations with the red grism were taken on the night of 1999 January 6, with a total integration time of 2400 seconds. Although the seeing for both nights was good, 0.6″ - 0.8″, neither night was photometric. The slit was aligned with the two images of the quasar. As is standard practice in the infrared, the object was observed at two positions along the slit. The strong and highly variable night sky features were effectively removed by subtracting the resulting spectra from each other. The 2-D sky-subtracted spectra were then flat-fielded, registered, and added. The two dimensional combined frames were spatially deconvolved in order to extract the spectrum of the lensing galaxy. For this purpose, we used the method outlined by Courbin et al. (1999, 2000). The algorithm is a spectroscopic extension of the so-called “MCS image deconvolution algorithm” (Magain et al. 1998). It spatially deconvolves 2-D spectra of blended objects, using the spectrum of a reference point source. It also improves their spatial sampling and decomposes them into the individual spectra of point sources (the two quasar images) and extended sources (the lensing galaxy). One also obtains a two-dimensional residual map, i.e., the difference between the data and the deconvolved spectrum (reconvolved by the spectrum of the PSF), in units of the photon noise. The quality of the deconvolution is checked using the residual map, which should be flat with a mean value of 1. The different products of the deconvolution are shown in Fig. 1 for the spectrum taken with the red grism. As the data were obtained before we developed our spectra deconvolution algorithm, we did not observe in an optimal way, in the sense that no reference spectrum was obtained (a spectrum of a star in the field of view). We aligned the slit along the two quasar components, as is usually done for such observations. However, the seeing of the data taken with the red grism was good enough to derive the PSF spectrum from the brighter quasar itself. This was not possible with the data taken with the blue grism. For the observations taken with the red grism, the spectrum of the lensing galaxy, was extracted from the 2-D deconvolved spectrum (extended component only) with standard aperture extraction techniques. The quasar spectra are a product of the deconvolution process and therefore do not show any contamination by the lensing galaxy. For the observations taken with the blue grism, we used wide apertures to extract the quasar spectra from the original data. The lens is therefore contaminating the spectra of the quasar, but by virtue of its very red color, the contamination is negligible. All extracted 1-D spectra were then divided by that of a bright star and multiplied by a blackbody curve that has a temperature that is appropriate for the spectral type of the star. Before the division, spectral features that were visible in the spectra of the bright star, such as the Pachen and Bracket lines of hydrogen, were removed by interpolation. ## 3 Near-IR spectroscopy of the lens ### 3.1 Plausible lens redshift The galaxy spectrum is shown in Fig. 2. The signal-to-noise ratio is very low, so the spectrum has been smoothed with a box car with 200 Å width. Also plotted are the broadband magnitudes of the lens (C98, R98 and Hjorth, private communication). The lens spectrum is scaled to match the $`H`$ and $`K`$ band magnitudes. The spectrum does not lead to a redshift measurement. However, the broadband colours suggest a significant break in the spectrum between the $`I`$ and $`J`$ bands. We have used the publicly available photometric redshift code hyperz (Bolzonella, Miralles & Pelló 2000) to estimate the redshift of the lens. Since there is little evidence for dust in the quasar spectrum (see below), we have fitted dust free models to the data. The best fitting model is a galaxy that was formed in a single burst of star formation. The best fitting redshift is $`z=1`$ with a 1-sigma range of 0.8 to 1.2. The age of the burst is 1.7 Gyrs. The quoted errors on the redshift do not include systematic errors that could be due to the heterogenous nature of the photometric data, which is derived from a mixture of ground and space based observations. The estimated redshift is slightly higher than those estimated from the position of the lens on the fundamental plane ($`z=0.77\pm 0.07`$; Kochanek et al. 2000) or from lens models and the time delay ($`z=0.79`$; W98). Note however that models including a dark component (see next section) can cope with any redshift between 0.7 and 1.3 and reproduce the observed time delay, assuming for example H$`{}_{0}{}^{}=60`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. As the break between 0.7 and 1.0 $`\mu `$ in the model spectra is very strong, a deep spectrum in this region is probably the key to accurately measure its redshift. ### 3.2 Evidence for a high redshift cluster-lens. With only two quasar images available to constrain the lensing potential, the unusual image configuration of HE 1104$``$1805 (R98) is very difficult to model uniquely. The system can not be reproduced with a Singular Isothermal Sphere. Additional shear and convergence, whatever their origin may be (intrinsic ellipticity of the lens and/or intervening lenses), are required to match simultaneously the positions and flux ratio of the quasar images ($`f_A/f_B`$=2.8, see section 4). In a first model, we introduce an ellipticity in the lens model, i.e., we choose an isothermal ellipsoid. Although we can easily obtain a good $`\chi ^2`$ fit, the resulting model has a very large velocity dispersion, over 300 km s<sup>-1</sup>, and an unrealistic ellipticity compared with the ellipticity of the associated light distribution. Finally, such a model predicts time delays of $`470h_{50}`$ days, while the observed value is 265 days, according to W98. The uncertainty on the lens redshift can not explain the discrepancy between the measured and predicted time delays: additional mass is required to describe the image configuration, flux ratio and time delay. We therefore adopt a two component model including (1) the main lensing galaxy, with ellipticity and position angle as constrained by the light distribution of the main galaxy lens, and, (2) a more extended component mimicking an intervening galaxy cluster. For simplicity, we centered the cluster on the main galaxy and assume an elliptical isothermal mass profile with a core radius. The fitted parameters were only the velocity dispersion, the ellipticity and position angle. Both the main lens and cluster components are taken to be at redshift 1.0. Our best fit model is shown in Fig. 3. It involves a cluster with moderate mass, i.e., a velocity dispersion of $`\sigma 575\pm 20`$ km s<sup>-1</sup>. The ellipticity is 0.3 (defined as $`e=[1(a/b)^2]/[1+(a/b)^2`$\]) and PA=10 degrees, which is slightly tilted relative to the axis of the main lens (which has PA=46 degrees) and with the light profile of the galaxy visible in the HST images of Lehar et al. (1999), who gives PA=63$`\pm `$17 degrees). Adding the cluster component also allows one to match better the observed shape parameters of the main lensing galaxy. With the presence of the cluster, the models can accommodate a PA of 46 degrees for the main lensing galaxy and a lower velocity dispersion, $`\sigma `$ 235 km s<sup>-1</sup>. If we assume H$`{}_{0}{}^{}=60`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }=0.3`$ and $`\mathrm{\Lambda }=0.7`$, the galaxy$`+`$cluster model reproduces well the measured time delay, giving a value of $`\mathrm{\Delta }t265`$ days. However, we stress that the lens redshift and the velocity dispersion of the cluster are redundant parameters: increasing the cluster’s mass or decreasing the lens redshift have the same effect on the time delay. This degeneracy between the two parameters will prevent any estimate of H<sub>0</sub> until more observational constraints are available on the cluster component of the lensing matter. The time delay now available in HE 1104$``$1805 can therefore be seen as a new important constraint on the lens model, indicating the presence of a yet undetected cluster, rather than a way to constrain H<sub>0</sub>. The mass within an area of a given radius is shown in Fig. 4 for different model components. If real, the cluster we predict in our model is difficult to detect, with only $`\sigma 575\pm 20`$ km s<sup>-1</sup>. At a redshift of 1, it would be even more difficult to see than the more massive clusters involved in other lenses such as AX J2019+112 (e.g., Benitez et al. 1999) or RX J0911.4+0551 (Burud et al. 1998, Kneib et al. 2000). However the velocity dispersion of such a cluster will change depending on the cluster center position. If it is not aligned with the main galaxy lens, its velocity dispersion will increase significantly. Deep X-ray observations and/or deep IR images of this field would be invaluable in constraining further the models. ## 4 Near-IR spectroscopy of the source at $`z`$ = 2.319 The 0.95 - 2.50 $`\mu `$ spectrum of the quasar pair is shown in Fig. 5. The spectra are on a relative flux scale. Regions of high atmospheric absorption are set to zero. The spectra show clearly the Balmer lines: H$`\alpha `$, H$`\beta `$, H$`\gamma `$ and a partially obscured H$`\delta `$, the \[OIII\] doublet and several broad FeII features (Francis et al. 1991). From the Balmer lines, the redshift is 2.323 for the brighter component (component A) and 2.321 for the fainter (component B). The measurement error is $`\mathrm{\Delta }`$z$`=0.002`$, so the redshift of the two components agree with each other, but are slightly larger than the determination at optical wavelengths ($`z=2.317`$, Smette et al. 1995). As with most quasars (McIntosh et al. 1999b), the \[OIII\] doublet is slightly blue shifted ($`z=2.319`$) with respect to the Balmer lines. Following Wisotzki et al. (1993), we subtracted a scaled version of the fainter component from the brighter one, that is $`f_\lambda (A)c.f_\lambda (B)`$. The scale is set so that the Balmer lines vanish. We find that we require $`c=2.9\pm 0.1`$ for the red spectrum and $`c=3.0\pm 0.1`$ for the blue spectrum. Wisotzki et al. and Smette et al. (1995) have used $`c=2.8`$. The slight difference between Wisotzki’s value and ours may only reflect systematic differences in the way the object was observed and the way the data were reduced rather than anything real. For example, the IR observations were done with a one-arc-second slit, and any small error in the alignment angle could cause such a difference. The difference spectra are plotted in Fig. 6. Here we plot the raw difference spectra as the dotted line, and a smoothed version of this as the continuous line. The spectrum of the brighter component is also displayed. The difference spectra are featureless. The residual after subtracting the strong H$`\alpha `$ line is less than 1%. Not only are the broad hydrogen features removed from the spectra, but the broad iron features and the \[OIII\] doublet are removed as well. As noted by Wisotzki et al. (1993) there appears to be excess continuum in the brighter component. ### 4.1 Extinction The Balmer decrement is around 4 for both components, and this is well within the range expected for unreddened quasars (e.g., Baker et al. 1994). Thus, there is no evidence for absolute reddening. However, the limits we can set on this are weak as the range of values for the Balmer decrement in quasars is rather broad. The limits for differential reddening are considerably stronger. The ratio of the emission lines in the brighter and fainter components is $`2.9\pm 0.1`$. The error brackets the measured variation of this ratio over time (six years of observations) and over wavelength. It is not clear if this variation is real or the result of measurement error. This ratio is remarkably constant over a large wavelength range, from CIV at 1549 Å to H$`\alpha `$, and we can used it to place an upper bound on the amount of differential extinction between the two components. If we assume that the lens is at $`z=1`$ and if we assume that the standard galactic extinction law (Mathis 1990) is applicable, then the differential extinction between the two components is $`\mathrm{\Delta }E(BV)<0.01`$ magnitudes. Recently, Falco et al. (1999) measure a differential extinction of $`\mathrm{\Delta }E(BV)=0.07\pm 0.01`$ for HE 1104$``$1805 in the sense that the B component has a higher extinction. However, their measurements rely on broad band photometry and their results can be mimicked by chromatic amplification of the continuum region by microlensing. If we were to repeat the experiment by comparing the relative strength of the continuum at 1.25 $`\mu `$ and 2.15 $`\mu `$, we would derive a differential extinction of $`\mathrm{\Delta }E(BV)=0.16`$ magnitudes and we would find also that B component was differentially reddened. ### 4.2 Emission line properties of the source. The emission line properties of high redshift quasars have been examined for correlations between line ratios and equivalent widths (McIntosh et al. 1999a,b; Muramaya et al. 1999). As the signal-to-noise ratio and spectral coverage of our IR data are considerably better, we have re-measured the emission line parameters for HE 1104$``$1805. Fitting of the spectrum was done in a similar way to that used in McIntosh et al. (1999a), but with extended spectral coverage. The model spectrum is a sum of Gaussian lines superposed on an exponential continuum to which is added a numerical optical FeII template (4250 Å and 7000 Å). The iron template consists of the optical spectrum of I Zw 1 obtained by Boroson & Green (1992). Before computing the model spectrum the template is smoothed to the resolution of our observations by convolving it with a Gaussian line which has the same FWHM than the broad emission lines of the quasar (CIV, in the present case), i.e., a rest-frame width of 6400 km s<sup>-1</sup>, or 14 Å. A systemic redshift of $`z_{syst}=2.319`$ is determined from the \[OIII\] $`\lambda `$5007 emission line and applied to the data to obtain a rest-frame spectrum (multiplied by $`1+z_{syst}`$ to conserve flux). As the positions of all other emission lines are redshifted relative to the \[OIII\] line by different amounts, their wavelengths are adjusted independently of each other. We measure a mean redshift of $`z_{Balmer}=2.323\pm 0.001`$ from the H$`\gamma `$ $`\lambda `$4340, H$`\beta `$ $`\lambda `$4861, and H$`\alpha `$ $`\lambda `$6562 emission lines. We used a sum of Gaussians to fit each Balmer line. This arbitrary decomposition is certainly not aimed at being representative of any physical model but still allows us to measure fluxes. One single Gaussian was used to fit the H$`\gamma `$ line while two Gaussians are required to fit H$`\beta `$ and three to fit H$`\alpha `$ which shows very wide symmetrical wings. The \[OIII\] doublet is represented by two Gaussians with a fixed line ratio of three (between \[OIII\] $`\lambda `$4959 and \[OIII\] $`\lambda `$5007). All line widths are fixed during the fit and all intensities are adjusted simultaneously (with the conjugate gradient algorithm) with the strength of the iron template and exponential continuum. The results of the fit are reported in Table 1 and Fig. 7. The best fit spectrum has a power law continuum of the type $`F_\lambda =\lambda ^\alpha `$, with $`\alpha =3.6`$. One-sigma errors were estimated by running the fit with different line widths. In addition to these errors, one should consider the error introduced by the continuum determination. Changing the index of the exponential continuum by 10% can affect iron flux measurement by up to 20%. The other, much narrower emission lines, are less affected, but we stress the need for continuum fitting over a very wide wavelength range in order to minimize such systematic errors. This was pointed out by Murayama et al. (1999). It is now obvious on our better data. The quality of the fit is overall very good; however, there are some regions where significant differences exist, as indicated by the residuals shown in the bottom left panel of Fig. 7. Most notably the FeII complex red-wards of \[OIII\] is relatively stronger that the FeII complex blue-wards of H$`\beta `$. The ratio of the EWs of \[FeII\] to H$`\beta `$ is 0.20. This is slightly lower than that measured by McIntosh et al. (1999a), who report 0.29$`{}_{0.09}{}^{}{}_{}{}^{+0.07}`$. The difference is probably not significant but there are two systematic biases that make a direct comparison difficult. Firstly, the continuum in this fit is well determined, whereas the small spectral coverage of the previous work may mean that EWs are underestimated (Muramaya et al. 1999). Secondly, and more fundamentally, EWs measured in lensed quasars are susceptible to microlensing which preferentially amplifies the continuum rather than the larger emission line regions. In fact, the continuum for unlensed quasars also varies. This means that EWs are a poor measure to use. Line fluxes are not susceptible to continuum variations, no matter how the continuum varies, whether it is intrinsic to the AGN or caused by microlensing. From our fit, we measure F(FeII(4810-5090)) / F(H$`\beta `$) = 0.18 $`\pm `$ 0.04 and F(FeII(4434-4685)) / F(H$`\beta `$) = 0.32 $`\pm `$ 0.04 which, according to Lipari et al. (1993) makes of HE 1104$``$1805 a rather weak FeII emitter. ## 5 Is microlensing detected in HE 1104$``$1805 ? The possibility of microlensing can be judged from a comparison between the size of the Einstein ring and the size of a quasar continuum emitting region. The latter is thought to be produced by an accretion disk and is of the order of $`10^{14}`$ to $`3\times 10^{15}`$ cm (Wambsganss et al. 1990; Krolik 1999 and references therein). The former depends on the mass of the microlenses, $`M`$, and is $`3.0\times 10^{16}\sqrt{(}M/M_{})`$ cm. In this calculation and those that follow, we have assumed that the lens is at $`z=1`$, and we have assumed a cosmology where $`H_0=60`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, and $`(\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda })=(0.3,0.0)`$. Thus, given that there is a suitable alignment, microlensing of the continuum is possible. Furthermore, the spectrum of the A component is considerably harder than that of the B component. This chromatic effect supports the idea that the continuum is microlensed since higher energy photons come from the inner part of the accretion disk, and are hence more susceptible to high amplification microlensing than lower energy photons. In fact, with suitable modeling of the lens and additional spectroscopic data, it may be possible to place constraints on accretion disc models (e.g., Agol & Krolik, 1999). The likelihood of microlensing then depends on the density of micro-lenses. If we model the mass distribution of the lensing galaxy as in section $`3.2`$, we can use the distance between the two macro images to determine the mass density at each image position. This gravitational convergence or optical depth, $`\kappa `$, is quite high for both components. For the A component, it is $`\kappa =0.73`$; for the B component, it is 0.53. If this is made entirely of stars then microlensing of either component is highly likely. In detail, however, only the main lensing galaxy is contributing to microlensing. The actual microlensing optical depth at the two quasar image positions is then lower, but still high. The typical time-scale between two consecutive microlensing events depends on the transverse velocity of the source and the velocity dispersion of the microlenses. The velocity dispersion for the lensing galaxy is high (using a single galaxy model one derives $`\sigma 300`$ km s<sup>-1</sup>, or about 235 km s<sup>-1</sup> if a cluster is also involved) and it is probably larger than the transverse bulk velocity. Dividing this velocity directly by the diameter of the Einstein ring, one derives a time scale of 3 years. This is quite long; however, it has been shown that stellar proper motions produce a higher microlensing rate than the one produced by a bulk velocity of the same magnitude (Wambsganss and Kundic 1995, Wyithe et al. 2000). Furthermore, the typical duration of a microlensing event is the time for the continuum emitting region ($`10^{15}\mathrm{cm}`$) to cross a caustic with velocity $`\sigma 300`$ km s<sup>-1</sup>. This is of the order of a few months and much shorter that the time between consecutive microlensing events. Thus, it is likely that microlensing affects the A component. As the stellar density of the lens near the B component is approximately half that of the A component, it is quite likely that microlensing affects the B component as well. We should expect that the continuum of the A component should be preferentially amplified relative to that of the B component for the majority of the time, but we should also expect that the B component should be preferentially amplified for a fraction of the time. HE 1104$``$1805 has now been monitored spectroscopically for six years (Wisotzki et al. 1998). During that time, the continuum of both components have been observed to vary; however, the continuum of the A component has always been harder (Wisotzki, private communication). As the time delay between the two components is of the order of 0.73 year (W98), the hardness of the continuum in the A component cannot be attributed to time delay effects. The most natural explanation is microlensing. Additionally, the relative level of the continuum of the A component is more variable than that of the B component (see Fig. 2 in W98). This cannot be attributed to photometric errors, because the A component is a factor of 3 brighter than B, both components are well separated and the lensing galaxy is much fainter than either component. Conversely, the BLR does not appear to be affected by microlensing. From 1993 to 1999, the ratio of the broad lines between the two components has varied little, with $`2.9\pm 0.1`$ (W98 and this paper). The lines of the BLR in the IR spectra presented here subtract very cleanly, better than 1% of the original line flux. Naively, one may then expect that any substructure in the BLR needs to be considerably larger than the microlensing caustics, i.e., $`3\times 10^{16}`$ cm. However, a more secure estimate requires better modeling of how microlensing in this particular lens can affect the profile of lines from the BLR (e.g., Schneider & Wambsganss, 1990). ## 6 Conclusions We have obtained 1 $`\mu `$ \- 2.5 $`\mu `$ spectra of the gravitational lens HE 1104$``$1805. Although we were not successful in measuring a precise redshift for the lens, the lens is probably an early type galaxy with a plausible redshift of $`0.8<z<1.2`$. This is slightly larger than estimates based on the measured time delay and estimates based from the position of the lens on the fundamental plane. We show however, that we can reconcile time delay and lens redshift by adding a cluster component to the lens models. We find that the continuum in the A component is harder than the continuum in the B component. The most probable explanation is that the A component is microlensed by compact objects in the lens galaxy. The ratio of the emission lines between the two components is $`2.9\pm 0.1`$. This is consistent with that measured at optical wavelengths. The constancy of this ratio over a large wavelength range limits strongly the amount of differential extinction between the two components. We find that the differential extinction is $`\mathrm{\Delta }E(BV)<0.01`$ magnitudes. We find that broad and narrow emission lines can be removed very well by subtracting a scaled version of the spectrum of component B from the spectrum of component A. The residual near the $`H\alpha `$ line is less than $`1\%`$ of the original line flux. It may be possible to use this near perfect subtraction to limit models of the BLR. This possibility should be investigated further. Finally, we note that the time delay measured in HE 1104$``$1805 allows us to demonstrate that the lensing potential is composed of a main lensing galaxy and a more extended “cluster” component. Without the time delay a single galaxy lens would also have been a viable model. With the rapidly increasing number of lenses with known time delay, we can therefore expect to constrain the content in dark matter of lens galaxies and it may be found that intervening clusters are a lot more frequent than first thought. ###### Acknowledgements. We would like to thank Daniel Mc Intosh and Todd Boroson for providing us with the FeII template used for the line fitting. F. Courbin acknowledges financial support through Chilean grant FONDECYT/3990024. Additional support from the European Southern Observatory and through a CNRS/CONICYT grant is also gratefully acknowledged. Jean-Paul Kneib acknowledges support from CNRS.
warning/0006/astro-ph0006278.html	ar5iv	text	# Giant repeated ejections from GRS 1915+105 ## 1 Introduction The ‘microquasar’ GRS 1915+105 is one of the most celebrated and widely-studied astrophysical objects of recent years. The system is extremely luminous and variable in both hard and soft X-rays (e.g. Foster et al. 1996; Morgan, Remillard & Greiner 1997; Belloni et al. 2000) and is a source of relativistic jets observed on arcsec and milliarcsec angular scales (Mirabel & Rodríguez 1994, hereafter MR94; Fender et al. 1999, hereafter F99; Rodríguez & Mirabel 1999, hereafter RM99; Dhawan, Mirabel & Rodríguez 2000). Sams, Eckart & Sunyaev (1996) have reported extended infrared emission from GRS 1915+105, but its relation to the radio ejections is at present unclear. X-ray dips on timescales of minutes have been interpreted by Belloni et al. (1997a,b) as the repeated disappearance and refill of the inner accretion disc, possibly due to extremely rapid transitions between ‘canonical’ black hole accretion states (Belloni et al. 2000). Pooley & Fender (1997; hereafter PF97) reported radio oscillations associated with such dips, and Fender et al. (1997; hereafter F97) discovered infrared analogs of these oscillations. The flat spectrum and correlated radio : infrared behaviour suggested that nonthermal synchrotron emission extended from the radio to the infrared regimes, the first time such high-frequency synchrotron emission had been observed from an X-ray binary (F97). Combined with the unstable accretion disc model of Belloni et al. (1997a,b) we suggested that a fraction of the inner disc was being repeatedly accelerated and ejected from the system (F97; PF97). Eikenberry et al. (1998a) confirmed the association between X-ray and infrared events, and Mirabel et al. (1998, hereafter M98) clearly observed the correlation between X-ray, infrared and radio behaviour in the source. Additional simultaneous observations (Fender & Pooley 1998, hereafter FP98) showed a very clear correlation between sequences of oscillations at radio and infrared wavelengths. Delays between different radio bands (PF97; M98) and between the infrared and radio bands (M98; FP98) clearly indicate that optical depth effects play an important role in the observed emission from these ejections. Eikenberry et al. (1998b) showed that infrared emission line strengths vary in proportion to the continuum during oscillations. A single X-ray dip, spectrally associated with the temporary disappearance of the inner accretion disc, was also found to coincide with a small radio flare (Feroci et al. 1999). More recently Eikenberry et al. (2000) report faint infrared flares whose association with the X-ray behaviour is uncertain, and Ogley et al. (2000) report significant flux from GRS 1915+105 at sub-millimetre wavelengths. ## 2 Observations GRS 1915+105 was observed simultaneously on 1999 May 20 with the United Kingdom Infrared Telescope (UKIRT) and the James Clerk Maxwell Telescope (JCMT), both located on Mauna Kea, Hawaii. ### 2.1 UKIRT GRS 1915+105 was observed with IRCAM3 in the infrared K-band ($`2.2\mu `$m) on 1999 May 20, simultaneously with the longer duration of JCMT SCUBA observations (see below). Data reduction and calibration were performed with iraf, along the lines described in F97. Five clear oscillation events were detected. The undereddened infrared flux densities reached 25 mJy at the peak of the oscillations, the largest amplitude oscillations reported to date in the infrared. The data are plotted in Figs 1 & 2, dereddened by $`A_\mathrm{K}=3.3`$ mag (F97; this value is still rather uncertain). ### 2.2 JCMT The 1350 $`\mu `$m detector of the SCUBA instrument (Holland et al. 1999) on JCMT was used in the photometry. Each integration lasted approximately 4 min. Calibration of the flux-density scale used observations of Mars and Uranus. The airmass ranged from 1.01 to 1.92 during the observations, and the optical depth at 1350 $`\mu `$m was less than 0.2 throughout. The data are plotted in Figs 1 & 2. ### 2.3 Radio In order to piece together the composite radio–mm–infrared spectrum of the source at the epoch of our observations, we have utilised radio data from two different monitoring programs. Firstly we have used public data at 2.3 & 8.3 GHz from the Green Bank Interferometer (GBI) monitoring program (e.g. Waltman et al. 1994). Observations at 15 GHz with the Ryle Telscope (RT, e.g. PF97) reveal strong oscillations for at least 9 days before and 2 days after the simultaneous UKIRT/JCMT observations (Fig 3), with a slowly rising trend in mean level and amplitude. ## 3 Temporal and spectral behaviour The JCMT observations reveal a sequence of 15 millimetre-wavelength oscillations with a quasi-period of $`23`$ min. Four of the oscillations have been observed simultaneously in the near-infrared, and both the (dereddened) near infrared and mm oscillations have an amplitude of 300–350 mJy. These are by far the largest amplitude oscillations ever observed from GRS 1915+105, including radio wavelengths (see e.g. PF97 for ‘typical’ radio oscillations). There is no detectable delay between the emission at the two wavelengths, to the time resolution of the JCMT data, $`4`$ min. F97, M98 and FP98 all report infrared oscillations of comparable amplitude (when dereddened by $`A_\mathrm{K}=3.3`$ mag) to radio oscillations observed around the same time. Although radio observations were not made strictly simultaneously with these mm/infrared observations, monitoring with the RT clearly reveals radio oscillations, with a comparable quasi-period of $`20`$ min, to have been occurring for at least two days before and afterwards (Fig 3). However, the amplitude of the mm and infrared oscillations is about five times greater than that observed at radio wavelengths. The mean radio – infrared spectrum for the four-day period illustrated in Fig 3 is shown in Fig 4. A striking feature of the light curves in Figs 1 & 2 is the infrared oscillation which starts around MJD 51318.59 which does not appear to have a mm counterpart, unlike the other four simultaneously observed events. We have carefully checked the data reduction techniques to see if this was due to human error, but found no evidence of this. We have no clear physical intepretation of this phenomena, except to suggest that it was due to strong and variable absorption which only signficantly affected the lower-frequency emission (for example the optical depth to free-free absorption, $`\tau \nu ^{2.1}`$). As noted above, the overall radio–mm–infrared spectrum at the time of these observations (Fig 4) was steeper than previously observed, perhaps also indicative of some absorption. If this is a correct interpretation of the ‘failed’ mm event, then we would have expected the radio emission to have been completely absorbed at this time also. ## 4 Energetics and mass outflow rate The radiative luminosity of these oscillations is large – for a flat spectrum of amplitude 300 mJy from 1 GHz to $`1.4\times 10^5`$ GHz ($`2.2\mu `$m), at a distance of 11 kpc, it is $`3\times 10^{37}`$ erg s<sup>-1</sup> (the time-averaged radiative luminosity is around half this value). As is the case for all synchrotron emitting plasmas for which adiabatic expansion losses dominate, this is likely to be a significant underestimate of the power being supplied to the jet. Furthermore we assume the emission arises in a partially self-absorbed jet which retains the same power-law distribution of electrons (i.e. $`p=2.6`$ where $`N(E)dEE^pdE`$) as observed in optically thin ejections (F99). In this situation the flat spectrum is produced by a conical, partially self-absorbed jet (e.g. Blandford & Königl 1979; Reynolds 1982). In addition several factors which can further affect the energy budget are uncertain, in particular whether or not the small ejections share the same bulk relativistic motions as the larger ejections, whether each radiating electron has an associated cold proton, and what the filling factor (ie. the effective volume) of the ejecta is. The procedure for calculating the energy and mass of the ejections is as follows: * Transform observed flux densities and frequencies back to their rest frame (identical if no bulk relativistic motion). * Integrate rest-frame luminosity. * Calculate maximum emitting volume – in this case based on the five-minute rise time this is $`3\times 10^{39}`$ cm<sup>3</sup>. The effective emitting volume is this volume multiplied by a ‘filling factor’, $`f`$. * From the volume, spectrum and luminosity, calculate equipartition magnetic field, and corresponding minimum internal energy. * For baryonic case, add one proton for each electron. * For cases with bulk relativistic motion, add in kinetic energy and multiply by two, under the assumption that observed emission was dominated by one (approaching) component only. * Divide by repetition quasi-period of oscillations to obtain time averaged energy and mass outflow rate. We have tabulated results for different cases in table 1; for bulk relativistic motion we have used the Doppler factors corresponding to $`\beta =0.98`$, $`\theta =66^{}`$ from F99. A significant constraint is that the lack of evidence for synchrotron losses (based on the similarity of decay rates at widely different wavelengths) at 2.2$`\mu `$m on a timescale of $`10`$ min, implies that $`B_{\mathrm{max}}30`$ G. For bulk motions with Doppler factor $`\delta `$, this limit is shifted slightly to $`B_{\mathrm{max}}\delta ^{1/3}`$ which, for $`\delta =0.34`$ in this case means the limiting field is $`40`$G, instead of $`30`$G, not a major difference. Thus while reducing the emitting volume via the filling factor (see Table 1) decreases the minimum energy, the stronger derived equipartition magnetic field is irreconcilable with the observed minimum lifetimes. As already noted in F97, a field of order 10 G will cause a cut-off in the spectrum of the oscillations around the optical band. Reducing the magnetic field below the equipartition value results in the internal energy being dominated by the electrons, for which total energy $`EB^{3/2}`$. Because of this constraint, the realistic minimum energy cannot be reduced much below the value for a non-relativistic non-baryonic ejection with filling factor $`f=1.0`$, which is $`4\times 10^{41}`$ erg, with a corresponding time-averaged power requirement of $`3\times 10^{38}`$ erg s<sup>-1</sup>. This corresponds to a radiative efficiency for the outflow of $`5`$%. This minimum power requirement is not strongly affected by our assumption of the spectral form of the electron distribution (e.g. for $`p=2.0`$ the minimum power is reduced by only a factor of three). Only in the case of baryonic ejections at high velocities does the large number of lower energy electrons become significant, as each has an associated proton. As a result it does not matter whether the emission arises in a discrete ‘plasmon’ or an internal shock in a quasi-steady flow. It is interesting to compare the power for baryonic ejections with $`\mathrm{\Gamma }=5`$, $`2\times 10^{43}`$ erg s<sup>-1</sup>, with that calculated for the same criteria for the ‘major’ ejections in F99, $`2\times 10^{39}`$ erg s<sup>-1</sup>. This is due to the observed optically thin cm spectrum, $`S_\nu \nu ^{0.8}`$, being assumed in F99 to have no high-frequency excess, and therefore a much lower integrated luminosity than the flat spectrum oscillations reported here. Further prompt mm and infrared observations during ‘major’ outbursts are required to investigate this. ## 5 Discussion M98 have shown that the wavelength-dependent time delays (radio–radio and radio–infrared) observed from GRS 1915+105 can be approximated by a ‘van der Laan’ (1996) model for an expanding plasmon. However, as discussed in FP98 such a model does not well describe the observed flat spectrum, which seems instead to be better modelled by a partially self-absorbed conical jet of the type developed for AGN (e.g. Blandford & Königl 1979; Reynolds 1982). More recently Kaiser, Sunyaev & Spruit (2000) have further applied a internal shock model to the major 1994 radio outburst of the source reported in MR94. In their model they require approximately the same amount of energy to be associated with the events as calculated for a plasmon model in MR94, but the power requirement is less as they spread the energy input over a much longer period. However, with repeated oscillations as observed here, this cannot be the case, where an entire accretion – ejection cycle is repeated on the timescale of $`20`$ min which we have used to calculate $`P`$ and $`\dot{M}_{\mathrm{jet}}`$ in table 1. Therefore this model cannot be used to evade the enormous amount of continuous power required to generate the observed repeated ejection events (this is not an argument against their model, but one against using it to evade the huge power requirements). Importantly, unless (a) there is a bright mm–infrared contribution from the large ejections which has not to date been observed, and (b) it is only the spatially resolved ejections (F99; RM99) which have bulk relativistic motion and a baryonic content, then GRS 1915+105 injects more energy and matter into the outflow during periods of repeated small events than it does during the large ejections. Note also that Belloni, Migliari & Fender (2000) have found that jet power, calculated as above, appears to be anticorrelated with accretion rate as inferred from X-ray spectral fits, for a small sample of observations with quasi-simultaneous infrared and X-ray coverage. ## 6 Conclusions We have reported giant repeated oscillation events from the black hole system GRS 1915+105 observed simultaneously at mm and infrared wavelengths. Contemporaneous radio observations indicate that these observations were near the end of a sequence of $`10`$ days of oscillations with $`20`$-min quasi-periods (i.e. $`700`$ discrete ejection events). We have investigated in depth the energy and mass flow associated with such events, seeking to minimize the very large required power. However the magnetic field has an upper limit imposed by the lack of observed radiation losses in the infrared band, so that reducing the effective volume by means of a small filling factor cannot significantly reduce the required power. Given the repeated nature of the events, and their almost certain association to an accretion cycle, we cannot spread the minimum required energy over a longer timescale than the repetition quasi-period. As a result we find that at least $`3\times 10^{38}`$ erg s<sup>-1</sup> is required to be channelled into the formation of the ejecta; a very significant fraction of the accretion energy unless the black hole is very massive indeed. The enormous energy budget associated with baryonic ejections at bulk relativistic velocities may rule out this possibility, although at present we do not know which is more unlikely: no baryons or low velocity ejections. ## Acknowledgements We would like to thank John Davies, Graeme Watt, Fred Baas, Ian Robson, Iain Coulson, Andy Adamson, Garret Cotter, Will Grainger, Mark Lacey, Susan Ridgeway, Tim Carroll and Thor Wold for assistance in the realisation of these observations, and Christian Kaiser for stimulating discussions. The JCMT is operated by The Joint Astronomy Centre on behalf of the UK Particle Physics and Astronomy Research Council (PPARC), the Netherlands Organisation for Scientific Research and the National Research Council of Canada. UKIRT is operated by The Observatories on behalf of the PPARC. We thank the staff at MRAO for maintenance and operation of the RT, which is supported by the PPARC.
warning/0006/math0006017.html	ar5iv	text	# Untitled Document The standard deviation effect (or why one should sit first base playing blackjack) E. Muñoz-García , R. Pérez-Marco<sup></sup><sup></sup> UCLA, Dept. of Mathematics, 405, Hilgard Ave., Los Angeles, CA-90095-1555, USA, e-mail: munoz@math.ucla.edu, ricardo@math.ucla.edu. Abstract. For a balanced cardcounting system we study the random variable of the true count after a number of cards are removed from the remaining deck and we prove a close formula for its standard deviation. As expected, the formula shows that the standard deviation increases with the number of cards removed. This creates a ”standard deviation effect” with a two fold consequence: longer long run and presumably larger fluctuations of the bankroll, but a small gain in playing accuracy for the player sitting third base. The opposite happens for the player sitting first base. Thus the optimal position in casino blackjack in terms of shorter long run is first base. Mathematics Subject Classification 2000 : 91A60, 05A19, 60C05 Key Words : Blackjack, true count, standard deviation, long run. Contents 1) Introduction. 2) Standard deviation of the true count. a) Proof of the true count theorem. b) The true count algebra. c) The true count theorem without induction. d) The true count standard deviation formula. 3) Long run. a) Kelly for binomial games. b) Kelly for fuzzy advantage. c) Long run. 4) Practical gambling. a) Comparison of first and third base and head-on play. a1) Full table. In company of ploppys a2) Head-on play. b) Absolute magnitude of the standard deviation. Bibliography. Disclaimer. Cardcounters are asked to forgive us the simplistic presentation of their activity aimed to (non-cardcounter) mathematicians, and mathematicians those elementary computations in the last section aimed to (non-mathematician) cardcounters. Not being gamblers or probability specialists this article is written for fun and with no pretention. 1) Introduction. This article is about casino blackjack and not about tournament blackjack. Cardcounters in the game of blackjack use different count systems to keep track of the ratio of favorable and unfavorable cards that remain on the deck. It is well known now and documented (\[Th\]) that one can get an edge (of the order of 2% ) over the house by increasing the bets in favorable counts and keep them to a minimum in unfavorable ones. The first person who published and analyzed this finding, mathematically and with computer simulations, was E. Thorp (who at the time was a graduate student in the mathematics department in UCLA). One of the main questions that a professional cardcounter player faces is the money management. He plays with a given bankroll. He has to decide the amount to bet at each moment taking into account the limit of the table and the upper bound fixed by his own bankroll. Obviously the maximum bet planned in his strategy should be inferior to the maximum of the table. But how much to bet depending on how favorable is the count ? The goal is to follow a betting pattern that minimizes the risk (that is loosing the whole bankroll) but maximizes the growth rate. The well known sharp strategy is Kelly’s criterion : The bet is a fraction of the total bankroll equal to the advantage you have. This maximizes the expected exponential rate of growth of the bankroll. Kelly’s criterion can be proven to be optimal in a very strong sense: Any other strategy will have a longer expected time to achieve a given amount of the bankroll (L. Breiman \[Br\] <sup></sup><sup></sup>At the time also a UCLA professor). For stablishing this type of results one makes the assumption that there is no minimal unit bet. In section 2 we give a simple derivation of Kelly criterion. For a basic bibliography we refer to \[RT\]. When the advantage is not know precisely but it is a random variable, Kelly criterion also applies (see section 3.c): One should bet according to the expected value of the advantage. The standard deviation of the exponential rate of growth of the bankroll turns out to increase in a significant way with the standard deviation of the advantage. This makes longer the ”long run” (see section 3) and presumably induces larger fluctuations. Thus one should prefer playing conditions that minimize the standard deviation of the advantage. In practice, count systems associate a value to each card. We restrict the discussion here to balanced counts, that is, those for which the total sum of values of a cards in a complete deck is zero. In blackjack only the numerical value of the card matters. A very simple balanced count system called Hi-Lo gives the weight $`1`$ to high cards (10,J,Q,K,A), the weight $`+1`$ to low cards (2,3,4,5,6) and the weight $`0`$ to medium cards (7,8,9). The cardcounter player adds up the weights of the cards as they are revealed. The numerical value obtained is the running count (RC). It tells the player how favorable is the remaining deck (from the Hi-Lo values you can tell that high cards are in favor of the player, low cards are in favor of the dealer). Of course how favorable is the deck depends not directly on the RC but on the true count (TC) that is the ratio of RC by the number of cards remaining (one can also divide by the number of decks remaining, this is just a change of units, and it is what is done in practice). The quantity $`52.TC`$ (i.e. the running count in deck units, $`52`$ is the total number of cards in a deck) gives approximatively the percentage in edge gained by the player due to the unbalanced composition of the deck. The advantage of the house with a game with favorable rules facing a perfect play with no counting from the player (this is called basic strategy) is of the order of $`0.5\%`$ (see for example, \[Gr2\] p.139). Thus it is recommended to beat the minimum when $`52.TC<2`$ and to bet $`52.TC`$ units when $`52.TC>1`$. The blackjack tables are semi-circular. The dealer stands on the flat side and at most seven players (in some tables five players) sit around the circular border. Cards are dealt from the left side to the right side of the dealer. The first position at the left of the dealer is the first base and the last position at the right of the dealer is the third base. The dealer deals the cards from left to right. It seems to be a common belief in the blackjack literature that there is no preferred position for the card counter and that the company of other players at the table has no influence (see \[HC\] p.68, \[Wo\] p. 220, for example). In this article we prove that the positions in the blackjack table are not all equivalent. There is an advantage in terms of shorter long run (and presumably of smaller fluctuations of the bankroll) for the player to be as close as possible to first base. More precisely, a player betting according to Kelly’s criterium will have a larger fluctuations of the exponential rate of growth of his bankroll if he sits on first base than if he sits in third base. The main noticeable effect of this is in the ”long run”. A player sitting third base will need to play about $`2\%`$ more favorable hands than a player sitting first base in order to achieve the same standard deviation for the exponential rate of growth of his bankroll. On the other hand, from the edge point of view, a player sitting closer to third base has a slight playing advantage over the players at his right. The main goal of this article is to explain theoretically these differences in the position of the players. They both have the same origin: the standard deviation effect. In a few words and in a simplistic form, the standard deviation effect appears when one has to take an irreversible decision with respect to a future situation. More the lag is important, more the decision has chances to be incorrect. The game of blackjack is played as follows. First the players decide their bets. A card counter will do this according to the value of the true count, and he will bet following Kelly’s criterium. Then the dealer will deal two cards to each player and two for him. The cards are dealt face up or down, this is not important for our purposes (in all cases there is always a hidden card, the so-called hole card of the dealer). Then the players will make their play decisions (to split, double, hit or stand) in turn from first to third base, and play the hand. When the moment of the play decision comes the true count of the player has changed because he has seen some of the new cards (at least the two he has received). Thus he may not be betting in a favorable situation. The so-called ”True count theorem” justifies the action. The expected value of the true count after some cards have been removed from the remaining deck is the same as the true count before removing the cards. As far as we know, the proof of this theorem in the blackjack context was publicished for the first time by Abdul Jalib M’Hall (\[JM\]) in a message posted to the rec.gambling.blackjack newsgroup. Theorem (True count theorem, A. J. M’Hall). For any balanced count, the expected value of the true count after several cards (but not all) are removed from the remaining deck is the value of the true count before removing the cards. This type of result goes back to the origins of probability theory. There is a well known problem: Each person in a room is asked to take a ticket from a box. One indicates the winner. Should you try to take your ticket as soon as possible or as late as possible? We assume that everyone waits that all tickets have been chosen to look at their ticket (otherwise if they look and reveal the result the problem is not the same). The answer is that it doesn’t matter of course. The expected value of the probability of being a winner (running count) stays the same at all moments. Coming back to the card counting problem, even if the expected value of the true count is independent of the number of cards removed, the standard deviation of the true count will increase as we remove cards from the deck. This results in particular from the exact formula proved in the theorem below. Thus the precise knowledge of the true count ”dilutes” as more cards are played. This makes that often the advantage the player has when his turn of play comes (and also when the dealer’s turn comes) differs from the expected value at the moment of betting. This induces a longer long run. Since the standard deviation increases with the number of cards played, the player sitting in third base will be systematically hurt by a larger standard deviation, thus will experience a longer long run. This is one of the consequences of the standard deviation effect. On the other hand, when the player in third base has to play, he is closer to the dealer’s play. Thus, compared to the player in first base, he knows more accurately the true count at the moment the dealer will play. Thus he can adjust more effectively his play. He has a supplementary bet advantage than the player in first base because he can deviate in a more efficient way from basic strategy. This is the second consequence of the standard deviation effect. The gain from deviating from basic strategy is small. Now we present the ”True count standard deviation formula” which proves the increasing of the standard deviation with the number of cards removed from the remaining deck. Let $`\sigma _n`$ be the standard deviation of the true count after having removed $`n1`$ cards. Theorem (True count standard deviation formula). Let $`N`$ be the number of cards remaining in the deck. After removing $`1n<N`$ cards, the standard deviation of the true count is $$\sigma _n=\sqrt{\frac{N1}{Nn}}\sqrt{n}\sigma _1.$$ The combinatorics in the proof of the true count standard deviation formula reveal a beautiful set of identities that we call the true count algebra. Each one of these identities has a probabilistic interpretation. We only develop in this article the minimum set of identities necessary in the proof of the main theorem. The distribution of the true count does depend on the composition of the deck and the weights of the counting system. This family of distributions contains the hypergeometric distribution. We don’t know a reference for the true count distribution in the classical literature. The authors will be grateful if any reader can provide one. We have the following corollaries (using also the results in section 2). Corollary 1. The standard deviation $`\sigma _n`$ is strictly increasing with $`n`$, even faster when the deck contains fewer cards. Corollary 2. For a given player, at the moment of taking his betting decision, the standard deviation of the true count at the moment of his play decision will be larger if he sits further away from the first base. Corollary 3. The theoretical long run of players sitting further away from the first base are larger. Corollary 4. In terms of shorter theoretical long run, the optimal seat is first base. Corollary 5. For a given player, at the moment of taking his playing decision, the standard deviation of the true count at the moment of the play of the dealer will be smaller if he sits further away from first base. Corollary 6. The playing advantage of players sitting further away from the first base is increased. Corollary 7. In terms of playing advantage the optimal seat is third base. We mainly discuss the effects on long run that is the most relevant one. There is also an effect on larger fluctuations of the bankroll and risk of ruin that is more difficult to analyze with precision. Classically the trade-off between betting efficiency and playing efficiency is something the card counter has to consider carefully when choosing his system (these quantities differ for different systems). Empirical ”proper balances” have been proposed (see \[Gr1\] p.40-49), but of course the proper balance between this two different quantities is up to each player and his goal. Practical gambling. In the last section we present a closer study of how the playing conditions do influence the standard deviation of the true count, thus the long run. The true count standard deviation formula only gives the relation of $`\sigma _n`$ with $`\sigma _1`$. The actual value of $`\sigma _1`$ does depend on the count system and on the remaining distribution of weights in the remaining cards. One has a quite accurate approximate formula only depending on the count system and the number of cards on the deck $$\sigma _1\frac{\mathrm{\Sigma }_0}{N}$$ where $`\mathrm{\Sigma }_0`$ is the standard deviation of weights used in the count system. We discuss some relevant consequences for practical play. Some amusing consequence of the formulas is that for a continuous model of the deck the standard deviation goes to $`+\mathrm{}`$ when $`N0`$. Fortunately (!) the casinos do not use to practice a $`100\%`$ penetration on the decks. Of course, in normal conditions an important penetration will increase in a substantial amount the playing advantage that will become the predominant effect. Nevertheless it is not excluded that casinos could device a set of rules that look advantageous to the players (according to the classical literature) but induce large fluctuations that will wipe out the bankrolls of card counters in the long run. The methods exposed in this article provide also information about higher moments of the true count. The combinatorics one faces is more involved. We plan to study this in the future. Acknowledgements. The authors are grateful to T. Ferguson, T. Liggett, R. Schonmann and B. Rothschild for instructive discussions in game theory, probability and combinatorics (domains in which the authors are amateurs). The authors thank also to F. Cheah and Sonia for their numerical experiments and gambling experience respectively, and also for motivating the present article. 2) Standard deviation of the true count. a) Proof of the true count theorem. To set up the notations we recall first Abdul Jalib M’Hall’s proof \[JM\]. It is enough to prove the result when we remove one card. We denote by $`N`$ the number of cards remaining in the deck. We denote by $`w`$ the different possible weights of the cards given by the count system. Let $`s_w`$ be the total number of cards in the full deck with weight $`w`$. Since the system count is balanced, we have $$\underset{w}{}ws_w=0.$$ Let $`l_w`$ be the total number of cards in the remaining deck with weight $`w`$ (so $`1=_wl_w/N`$). Then if we denote by $`k_w=s_wl_w`$ we have that $$R=\underset{w}{}wk_w=\underset{w}{}wl_w$$ is the running count. The true count is $`T=R/N`$. We compute the expected value $`T_1`$ of the true count after removing $`1`$ card. It is given by $$\begin{array}{cc}\hfill T_1& =\underset{w}{}\frac{R+w}{N1}\frac{l_w}{N}\hfill \\ & =\underset{w}{}\frac{R}{N1}\frac{l_w}{N}+\underset{w}{}\frac{w}{N1}\frac{l_w}{N}\hfill \\ & =\frac{R}{N1}\frac{R}{(N1)N}\hfill \\ & =\frac{R}{N}\hfill \\ & =T\hfill \end{array}$$ and the result follows. b) The true count algebra. For $`1pN`$ we define $$l_w^{w_1\mathrm{}w_p}=l_w\|\{1ip;w_i=w\}\|,$$ that is the number of cards of weight $`w`$ remaining once $`p`$ cards of weight $`w_1,\mathrm{},w_p`$ have been removed. Observe that we have $$\underset{w}{}\frac{l_w^{w_1\mathrm{}w_p}}{Np}=1.$$ These coefficients have a beautiful and rich combinatorics. We call this system of identities the true count algebra. The name is chosen because most relevant formulas have a ”true count” probabilistic interpretation. We concentrate here into the relevant identities in order to compute the standard deviation. The full algebra will be studied in a forthcoming article. Lemma 1. For $`0pN2`$ we have $$S_0=\underset{w}{}\frac{l_w^{w_1\mathrm{}w_p}}{Np}\frac{l_{v_0}^{w_1\mathrm{}w_pw}}{Np1}=\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}.$$ One can translate this relation in a probabilistic statement. Picking at random one card on the remaining deck, the probability that this card has a weight $`v_0`$ is the same than the probability of this same occurrence using the same deck depleted from one card at random. Or, in card counters terms, the running count for weight $`v_0`$ does not change when we remove one card at random. Proof. We split the sum into a sum over all $`wv_0`$ and the term for $`w=v_0`$. Observe that for $`wv_0`$ we have $$l_{v_0}^{w_1\mathrm{}w_pw}=l_{v_0}^{w_1\mathrm{}w_p}.$$ Also for $`w=v_0`$ we have $$l_{v_0}^{w_1\mathrm{}w_pv_0}=l_{v_0}^{w_1\mathrm{}w_p}1.$$ Thus we have $$\begin{array}{cc}\hfill S_0& =\left(\underset{wv_0}{}\frac{l_w^{w_1\mathrm{}w_p}}{Np}\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np1}\right)+\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}\frac{l_{v_0}^{w_1\mathrm{}w_pv_0}}{Np1}\hfill \\ & =\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np1}\left(\underset{wv_0}{}\frac{l_w^{w_1\mathrm{}w_p}}{Np}\right)+\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}\frac{l_{v_0}^{w_1\mathrm{}w_p}1}{Np1}\hfill \\ & =\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np1}\left(1\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}\right)+\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}\frac{l_{v_0}^{w_1\mathrm{}w_p}1}{Np1}\hfill \\ & =l_{v_0}^{w_1\mathrm{}w_p}\left(\frac{1}{Np1}\frac{1}{(Np)(Np1)}\right)\hfill \\ & =\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}\hfill \end{array}$$ q.e.d.$`\mathrm{}`$ Lemma 2. Let $`p0`$, $`q0`$ such that $`p+qN2`$. We have the formula $$\begin{array}{cc}\hfill S_q& =\underset{w}{}\frac{l_w^{w_1\mathrm{}w_p}}{Np}\frac{l_{v_0}^{w_1\mathrm{}w_pw}}{Np1}\frac{l_{v_1}^{w_1\mathrm{}w_pv_0w}}{Np2}\mathrm{}\frac{l_{v_q}^{w_1\mathrm{}w_pv_0\mathrm{}v_{q1}w}}{Npq1}\hfill \\ & =\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}\frac{l_{v_1}^{w_1\mathrm{}w_pv_0}}{Np1}\mathrm{}\frac{l_{v_q}^{w_1\mathrm{}w_pv_0\mathrm{}v_{q1}}}{Npq}\hfill \end{array}$$ Clearing denominators one can also write the less cumbersome formula $$\begin{array}{cc}& \underset{w}{}l_w^{w_1\mathrm{}w_p}l_{v_0}^{w_1\mathrm{}w_pw}l_{v_1}^{w_1\mathrm{}w_pv_0w}\mathrm{}l_{v_q}^{w_1\mathrm{}w_pv_0\mathrm{}v_{q1}w}\hfill \\ & =(Npq1)l_{v_0}^{w_1\mathrm{}w_p}l_{v_1}^{w_1\mathrm{}w_pv_0}\mathrm{}l_{v_q}^{w_1\mathrm{}w_pv_0\mathrm{}v_{q1}}\hfill \end{array}$$ One can translate the identity of the lemma in a probabilistic statement. Picking at random $`q`$ cards on the remaining deck one after the other, the probability that the cards have weights $`v_0`$, $`v_1`$ …and $`v_q`$ in this order is the same than the probability of this same occurrence using the same deck depleted from one card at random. Or, in card counters terms, the running count for an ordered clump of $`q+1`$ cards to have respective weights $`v_0`$, $`v_1`$, … and $`v_q`$ does not change when we remove one card at random. Proof. The proof is done by induction on $`q0`$. The result holds for $`q=0`$ because of lemma 1. We assume that the result holds for $`q0`$ and we prove it for $`q+1`$. We split the sum $$\underset{w}{}l_w^{w_1\mathrm{}w_p}l_{v_0}^{w_1\mathrm{}w_pw}l_{v_1}^{w_1\mathrm{}w_pv_0w}\mathrm{}l_{v_{q+1}}^{w_1\mathrm{}w_pv_0\mathrm{}v_qw}$$ into the sum for $`wv_0`$, $$A=\underset{wv_0}{}l_w^{w_1\mathrm{}w_p}l_{v_0}^{w_1\mathrm{}w_pw}l_{v_1}^{w_1\mathrm{}w_pv_0w}\mathrm{}l_{v_{q+1}}^{w_1\mathrm{}w_pv_0\mathrm{}v_qw}$$ and the term of the sum for $`w=v_0`$ $$B=l_{v_0}^{w_1\mathrm{}w_p}l_{v_0}^{w_1\mathrm{}w_pv_0}l_{v_1}^{w_1\mathrm{}w_pv_0v_0}l_{v_2}^{w_1\mathrm{}w_pv_0v_1v_0}\mathrm{}l_{v_{q+1}}^{w_1\mathrm{}w_pv_0\mathrm{}v_qv_0}$$ For $`wv_0`$ we have $`l_{v_0}^{w_1\mathrm{}w_pw}=l_{v_0}^{w_1\mathrm{}w_p}`$, so $$A=l_{v_0}^{w_1\mathrm{}w_p}\underset{wv_0}{}l_w^{w_1\mathrm{}w_p}l_{v_1}^{w_1\mathrm{}w_pv_0w}\mathrm{}l_{v_{q+1}}^{w_1\mathrm{}w_pv_0\mathrm{}v_qw}$$ Since for $`wv_0`$ we have $`l_w^{w_1\mathrm{}w_p}=l_w^{w_1\mathrm{}w_pv_0}`$, we have $$A=l_{v_0}^{w_1\mathrm{}w_p}\underset{wv_0}{}l_w^{w_1\mathrm{}w_pv_0}l_{v_1}^{w_1\mathrm{}w_pv_0w}\mathrm{}l_{v_{q+1}}^{w_1\mathrm{}w_pv_0\mathrm{}v_qw}$$ Using the induction hypothesis for $`q`$ we conclude that $$\begin{array}{cc}\hfill A=& l_{v_0}^{w_1\mathrm{}w_p}[(N(p+1)q)l_{v_1}^{w_1\mathrm{}w_pv_0}\mathrm{}l_{v_{q+1}}^{w_1\mathrm{}w_pv_0\mathrm{}v_q}\hfill \\ & l_{v_0}^{w_1\mathrm{}w_pv_0}l_{v_1}^{w_1\mathrm{}w_pv_0v_0}l_{v_2}^{w_1\mathrm{}w_pv_0v_1v_0}\mathrm{}l_{v_{q+1}}^{w_1\mathrm{}w_pv_0\mathrm{}v_qv_0}]\hfill \\ & =(Np(q+1))l_{v_0}^{w_1\mathrm{}w_p}l_{v_1}^{w_1\mathrm{}w_pv_0}\mathrm{}l_{v_{q+1}}^{w_1\mathrm{}w_pv_0\mathrm{}v_q}B\hfill \end{array}$$ And the result follows.$`\mathrm{}`$ From the probabilistic interpretation we obtain that there are corresponding formulas for the removal of $`k`$ cards instead of only one. We prove these formulas by induction on $`k`$ using lemma 1 and lemma 2. Lemma 3. Let $`k1`$ and $`0pNk1`$. We have $$\underset{i_1,\mathrm{},i_k}{}\frac{l_{i_1}^{w_1\mathrm{}w_p}}{Np}\frac{l_{i_2}^{w_1\mathrm{}w_pi_1}}{Np1}\frac{l_{i_3}^{w_1\mathrm{}w_pi_1i_2}}{Np2}\mathrm{}\frac{l_{i_k}^{w_1\mathrm{}w_pi_1\mathrm{}i_{k1}}}{Npk+1}\frac{l_{v_0}^{w_1\mathrm{}w_pi_1\mathrm{}i_k}}{Npk}=\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}$$ Lemma 4. Let $`p0`$, $`q0`$ and $`k1`$ such that $`p+k+qN1`$. We have $$\begin{array}{cc}& \underset{i_1,\mathrm{},i_k}{}\frac{l_{i_1}^{w_1\mathrm{}w_p}}{Np}\frac{l_{i_2}^{w_1\mathrm{}w_pi_1}}{Np1}\mathrm{}\frac{l_{i_k}^{w_1\mathrm{}w_pi_1\mathrm{}i_{k1}}}{Npk1}\frac{l_{v_0}^{w_1\mathrm{}w_pi_1\mathrm{}i_{k1}i_k}}{Npk}\hfill \\ & \frac{l_{v_1}^{w_1\mathrm{}w_pv_0i_1\mathrm{}i_k}}{Npk1}\mathrm{}\frac{l_{v_q}^{w_1\mathrm{}w_pv_0\mathrm{}v_{q1}i_1\mathrm{}i_k}}{Npkq}\hfill \\ & =\frac{l_{v_0}^{w_1\mathrm{}w_p}}{Np}\frac{l_{v_1}^{w_1\mathrm{}w_pv_0}}{Np1}\mathrm{}\frac{l_{v_q}^{w_1\mathrm{}w_pv_0\mathrm{}v_{q1}}}{Npq}\hfill \end{array}$$ There are more general formulas, but we don’t need them for the purpose of this article. Next lemma is an application of the formulas. Lemma 5. Let $`W=\{w\}`$ be the set of weights and $`f:W𝐑`$ a real valued function. The expected value of $`f(w)`$, for any $`1nN`$ and $`1jn`$, is given by $$\overline{f}=\underset{w}{}f(w)\frac{l_w}{N}=\underset{w_1\mathrm{}w_n}{}f(w_j)\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}.$$ Proof. We have $$\begin{array}{cc}& \underset{w_1\mathrm{}w_n}{}f(w_j)\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}\hfill \\ & =\underset{w_1\mathrm{}w_j}{}f(w_j)\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_j}^{w_1\mathrm{}w_{j1}}}{Nj+1}\underset{w_{j+1},\mathrm{}w_n}{}\frac{l_{w_{j+1}}^{w_1\mathrm{}w_j}}{Nj}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}\hfill \\ & =\underset{w_1\mathrm{}w_j}{}f(w_j)\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_j}^{w_1\mathrm{}w_{j1}}}{Nj+1}\hfill \\ & =\underset{w_j}{}f(w_j)\underset{w_1,\mathrm{},w_{j1}}{}\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_j}^{w_1\mathrm{}w_{j1}}}{Nj+1}\hfill \\ & =\underset{w_j}{}f(w_j)\frac{l_{w_j}}{N}\hfill \\ & =\overline{f}\hfill \end{array}$$ where we used lemma 4 that gives $$\underset{w_{j+1},\mathrm{}w_n}{}\frac{l_{w_{j+1}}^{w_1\mathrm{}w_j}}{Nj}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}=\underset{w_n}{}\frac{l_{w_n}^{w_1\mathrm{}w_j}}{Nj}=1$$ and $$\underset{w_1,\mathrm{},w_{j1}}{}\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_j}^{w_1\mathrm{}w_{j1}}}{Nj+1}=\frac{l_{w_j}}{N}.$$ $`\mathrm{}`$ c) The true count theorem without induction. We present here a direct computation of the expected value of the true count when $`n`$ cards have been withdraw. It uses the true count theorem for $`n=1`$ but proves the general theorem without induction. This is certainly useless, but it prepares the field and the notations to prove the standard deviation true count formula. Remember that we denote by $`R`$ the running count. The following lemma is immediate by direct computation. Lemma 6. We have the identity $$\frac{R+w_1+\mathrm{}+w_n}{Nn}\frac{R}{N}=\frac{N1}{Nn}\left[\left(\frac{R+w_1}{N1}\frac{R}{N}\right)+\mathrm{}+\left(\frac{R+w_n}{N1}\frac{R}{N}\right)\right]$$ Now the expected value of the true count after removing $`n`$ cards of the deck is (we use lemma 5 here) $$\begin{array}{cc}\hfill T_n& =\underset{w_1,\mathrm{},w_n}{}\frac{R+w_1+\mathrm{}+w_n}{Nn}\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}\hfill \\ & =\frac{R}{N}+\underset{w_1,\mathrm{},w_n}{}\left(\frac{R+w_1+\mathrm{}+w_n}{Nn}\frac{R}{N}\right)\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}\hfill \end{array}$$ Consider the function $$f(w)=\frac{R+w}{N1}\frac{R}{N}.$$ We have (using the true count theorem for $`n=1`$ (!)) $$\overline{f}=0$$ Thus using lemma 6 and lemma 5, $$T_n=\frac{R}{N}+\frac{N1}{Nn}n\overline{f}=\frac{R}{N}.$$ d) The true count standard deviation formula. We still denote by $`f`$ the function defined in the previous section. The square of the standard deviation is given by $$\sigma _n^2=\underset{w_1,\mathrm{},w_n}{}\left(\frac{R+w_1+\mathrm{}+w_n}{Nn}\frac{R}{N}\right)^2\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}$$ Using lemma 6 and developing the square we have $$\sigma _n^2=\left(\frac{N1}{Nn}\right)^2\underset{1j_1,j_2n}{}\underset{w_1,\mathrm{},w_n}{}f(w_{j_1})f(w_{j_2})\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}$$ Immediately from lemma 5 we get Lemma 7. For $`j_1=j_2=j`$, we have $$\underset{w_1,\mathrm{},w_n}{}f(w_j)^2\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}=\sigma _1^2$$ Using next lemma we can finish the computation: $$\begin{array}{cc}\hfill \sigma _n^2& =\left(\frac{N1}{Nn}\right)^2\left(n\sigma _1^2n(n+1)\frac{1}{N1}\sigma _1^2\right)\hfill \\ & =\left(\frac{N1}{Nn}\right)^2\left(1\frac{n1}{N1}\right)n\sigma _1^2\hfill \\ & =\left(\frac{N1}{Nn}\right)n\sigma _1^2\hfill \end{array}$$ q.e.d. Lemma 8. We have for $`j_1j_2`$, $$\underset{w_1,\mathrm{},w_n}{}f(w_{j_1})f(w_{j_2})\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}=\frac{1}{N1}\sigma _1^2$$ Proof. We assume $`j_1<j_2`$ for example. The proof follows the same ideas than the proof of lemma 5. We have using lemma 4 several times $$\begin{array}{cc}& \underset{w_1,\mathrm{},w_n}{}f(w_{j_1})f(w_{j_2})\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_n}^{w_1\mathrm{}w_{n1}}}{Nn+1}\hfill \\ & =\underset{w_1,\mathrm{},w_{j_2}}{}f(w_{j_1})f(w_{j_2})\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_{j_2}}^{w_1\mathrm{}w_{j_21}}}{Nj_2+1}\hfill \\ & =\underset{w_{j_2}}{}f(w_{j_2})\underset{w_1,\mathrm{},w_{j_21}}{}f(w_{j_1})\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_{j_2}}^{w_1\mathrm{}w_{j_21}}}{Nj_2+1}\hfill \\ & =\underset{w_{j_2}}{}f(w_{j_2})\underset{w_1,\mathrm{},w_{j_1}}{}f(w_{j_1})\underset{w_{j_1+1},\mathrm{},w_{j_21}}{}\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_{j_1+1}}^{w_1\mathrm{}w_{j_1}}}{Nj_1}\mathrm{}\frac{l_{w_{j_2}}^{w_1\mathrm{}w_{j_21}}}{Nj_2+1}\hfill \\ & =\underset{w_{j_2}}{}f(w_{j_2})\underset{w_1,\mathrm{},w_{j_1}}{}f(w_{j_1})\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_{j_1}}^{w_1\mathrm{}w_{j_11}}}{Nj_1+1}\frac{l_{w_{j_2}}^{w_1\mathrm{}w_{j_1}}}{Nj_1}\hfill \\ & =\underset{w_{j_2}}{}f(w_{j_2})\underset{w_{j_1}}{}f(w_{j_1})\underset{w_1,\mathrm{},w_{j_11}}{}\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\mathrm{}\frac{l_{w_{j_1}}^{w_1\mathrm{}w_{j_11}}}{Nj_1+1}\frac{l_{w_{j_2}}^{w_1\mathrm{}w_{j_1}}}{Nj_1}\hfill \\ & =\underset{w_{j_1},w_{j_2}}{}f(w_{j_1})f(w_{j_2})\frac{l_{w_{j_1}}}{N}\frac{l_{w_{j_2}}^{w_{j_1}}}{N1}\hfill \\ & =\frac{1}{N1}\sigma _1^2\hfill \end{array}$$ In the last computation we use the following lemma applied to the function $`g(w_1,w_2)=f(w_1)f(w_2)`$ (observe that $`_wg(w,w)(l_w/N)=_wf(w)^2(l_w/N)=\sigma _1^2`$). $`\mathrm{}`$ Lemma 9. Let $`g:W\times W𝐑`$. We have $$\underset{w_1,w_2}{}g(w_1,w_2)\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}=\frac{N}{N1}\overline{g}\frac{1}{N1}\underset{w}{}g(w,w)\frac{l_w}{N}$$ where $`\overline{g}`$ is the expected value of $`g`$ $$\overline{g}=\underset{w_1,w_2}{}g(w_1,w_2)\frac{l_{w_1}}{N}\frac{l_{w_2}}{N}.$$ Proof. We have $$\begin{array}{cc}& \underset{w_1,w_2}{}g(w_1,w_2)\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}\hfill \\ & =\underset{w_1w_2}{}g(w_1,w_2)\frac{l_{w_1}}{N}\frac{l_{w_2}^{w_1}}{N1}+\underset{w}{}g(w,w)\frac{l_w}{N}\frac{l_w1}{N1}\hfill \\ & =\frac{N}{N1}\left(\overline{g}\underset{w}{}g(w,w)\frac{l_w}{N}\frac{l_w}{N}\right)+\frac{N}{N1}\left(\underset{w}{}g(w,w)\frac{l_w}{N}\frac{l_w1}{N1}\right)\hfill \\ & =\frac{N}{N1}\overline{g}+\frac{N}{N1}\left(\underset{w}{}g(w,w)\frac{l_w}{N}\left(\frac{1}{N}\right)\right)\hfill \\ & =\frac{N}{N1}\overline{g}\frac{1}{N1}\underset{w}{}g(w,w)\frac{l_w}{N}\hfill \end{array}$$ $`\mathrm{}`$ 3) Long run. a) Kelly for binomial games. We review in this section the Kelly criterion in the case of an iterated fixed advantage game. We review also the formulas for the expected value and the variance of the exponential rate of growth and discuss the implications for Blackjack play. We assume that we play a repetitive independent game where we have an advantage $`p>1/2`$. We expect an exponential growth of our initial bankroll $`X_0`$ if we follow a reasonable strategy of betting. We assume that there is no minimal unit of bet. By homogeneity of the problem, the sharp strategy will consists in betting a proportion $`f(p)`$ of the total bankroll. Our bankroll after $`n`$ rounds of the game have been played is $$X_n=X_0\underset{i=1}{\overset{n}{}}(1+\epsilon _if(p))$$ where $`\epsilon _i=+1`$ if we won the $`i`$-th hand, and $`\epsilon _i=1`$ if we lost the $`i`$-th hand. The exponential rate of growth of the bankroll is $$G_n=\frac{1}{n}\mathrm{log}\frac{X_n}{X_0}=\frac{1}{n}\underset{i=0}{\overset{n}{}}\mathrm{log}(1+\epsilon _if(p)).$$ The Kelly criterion maximizes the expected value of the exponential rate of growth (the proof is elementary, see section b for a proof in the more general setting of a ”fuzzy advantage”). Theorem 3.1 (Kelly criterion). The expected value of the exponential rate of growth $`G_n`$ is maximized for $$f(p)=2p1.$$ Observe that the expected value is $$𝐄(G_n)=𝐄(G_1)=p\mathrm{log}(1+f(p))+(1p)\mathrm{log}(1f(p)).$$ Also the random variables $$X_i=\mathrm{log}(1+\epsilon _if(p))$$ are independent, and $$\mathrm{𝐕𝐚𝐫}G_n=\frac{1}{n}\mathrm{𝐕𝐚𝐫}X.$$ We compute $$\begin{array}{cc}\hfill \mathrm{𝐕𝐚𝐫}X& =𝐄(X^2)(𝐄(X))^2\hfill \\ & =p\left(\mathrm{log}(1+f(p))\right)^2+(1p)\left(\mathrm{log}(1f(p))\right)^2\left(p\mathrm{log}(1+f(p))+(1p)\mathrm{log}(1f(p))\right)^2\hfill \\ & =p(1p)\left(\mathrm{log}\left(\frac{1+f(p)}{1f(p)}\right)\right)^2\hfill \end{array}$$ So finally $$\mathrm{𝐕𝐚𝐫}G_n=\frac{1}{n}p(1p)\mathrm{log}\left(\frac{p}{1p}\right)^2.$$ Proposition 3.2. The expected value and the standard deviation of $`G_n`$ are $$\begin{array}{cc}\hfill 𝐄(G_n)& =𝐄(G_1)=p\mathrm{log}(2p)+(1p)\mathrm{log}(22p)\hfill \\ \hfill \sigma (G_n)& =\frac{1}{\sqrt{n}}\sqrt{p(1p)}\mathrm{log}\left(\frac{p}{1p}\right)\hfill \end{array}$$ Thus for a slight advantage $`p=1/2+\epsilon `$ we have $$\begin{array}{cc}\hfill 𝐄(G_n)& =2\epsilon ^2+𝒪(\epsilon ^3)\hfill \\ \hfill \sigma (G_n)& =\frac{2\epsilon }{\sqrt{n}}(1+𝒪(\epsilon ^2))\hfill \end{array}$$ From this data in becomes clear how long is ”the long run”. The long run corresponds to play a number of rounds such that $`\sigma (G_n)`$ is of the order of $`𝐄(G_n)`$ (or a small fraction of it). For a smaller number of rounds played there is a good chance that $`G_n`$ is negative, i.e. the bankroll has decreased. In the case of blackjack, $`\epsilon 10^2`$. Thus the ”long run” corresponds to have played of the order of $`10000`$ favorable hands, thus at least $`20000`$ hands. At $`50`$ hands per hour this adds to $`4000`$ hours of play. To play less hands means that we are gambling. This explains that even at the non-professional level the team play makes perfect sense in order to get into the long run quicker (divide the time by the number of members of the team). The situation one faces when playing blackjack is different than a binomial game with a fixed advantage. First the advantage is not the same at different hands. This involves minor modifications of the above computations. Second, due to the standard deviation effect, the advantage the player has at the moment of playing is a random variable at the moment of betting. So he is betting according to a fuzzy advantage. We study this situation in the next sections. Kelly criterion is still sharp. The main effect of the fuzzy advantage is to introduce a supplementary term in the standard deviations of $`G_n`$. b) Kelly for fuzzy advantage. We consider the situation one faces playing blackjack. At each round one decides the amount to bet in function of the true count. We assume that we play multiple independent rounds of the same game. At each round the advantage we have is a random variable $`p`$, $`0<p<1`$, with a known distribution $`D_x(p)`$ and expected value $`p_0(x)>1/2`$ from a family of distributions $`(D_x)`$. The choice of the distribution $`D_x`$ at each round is random with distribution $`\rho (x)`$, $`x[0,1]`$ ($`x`$ is a mere index, we may just adjust it to have a uniform distribution $`\rho (x)=1`$). We want to maximize the expected value of the exponential rate of growth of our bankroll $`X_0`$. The optimal amount to bet is a proportion $`0f(D_x)1`$ of the bankroll depending only on the distribution $`D_x`$. For a sharp strategy we have $`f(D_x)=0`$ when $`p_0(x)1/2`$ (that is, the game is not favorable). The quantity to maximize is $$\begin{array}{cc}& _0^1_0^1\left(p\mathrm{log}(1+f(D_x))+(1p)\mathrm{log}(1f(D_x))\right)D_x(p)𝑑p\rho (x)𝑑x\hfill \\ & =_0^1\left(p_0(x)\mathrm{log}(1+f(D_x))+(1p_0(x))\mathrm{log}(1f(D_x))\right)\rho (x)𝑑x\hfill \end{array}$$ Theorem 3.3 (Kelly criterion). The optimal strategy is obtained for $$f_K(D_x)=2p_0(x)1$$ for $`1/2p_0(x)1`$, and $`f_K(D_x)=0`$ for $`0p_0(x)1/2`$. Proof. We want to maximize the functional $$\begin{array}{cc}\hfill 𝒢(f)& =_0^1\left(p_0(x)\mathrm{log}(1+f(D_x))+(1p_0(x))\mathrm{log}(1f(D_x))\right)\rho (x)𝑑x\hfill \\ & =_{\{x;p_0(x)>1/2\}}\left(p_0(x)\mathrm{log}(1+f(D_x))+(1p_0(x))\mathrm{log}(1f(D_x))\right)\rho (x)𝑑x\hfill \end{array}$$ where we assume that $`f(D_x)=0`$ for $`0<p_0(x)1/2`$. If $`f_0`$ is an extremum then for any perturbation $`h`$ such that $`f_0+\epsilon h`$ is an allowable strategy ($`f_0+\epsilon h>0`$), then if we consider $$g(\epsilon )=𝒢(f_0+\epsilon h)$$ we must have $$g^{}(0)=0.$$ A direct computation gives $$g^{}(\epsilon )=_{\{x;p_0(x)>1/2\}}\frac{h(D_x)}{1\left(f_0(D_x)+\epsilon h(D_x)\right)^2}[(2p_0(x)1)f_0(D_x)\epsilon h(D_x)]\rho (x)𝑑x>$$ So $$g^{}(0)=_{\{x;p_0(x)>1/2\}}\frac{h(D_x)}{1f_0(D_x)^2}(2p_0(x)1f_0(x))\rho (x)𝑑x.$$ Thus when $`p_0(x)>1/2`$ we must have $`f_0(x)=2p_0(x)1=f_K(x)`$. Also by direct computation we have $$g^{\prime \prime }(\epsilon )=_{\{x;p_0(x)>1/2\}}\frac{h^2}{(1(f_0+\epsilon h)^2)^2}\left[(f_0+\epsilon h)^2+12(2p_0(x)1)(f_0+\epsilon h)\right]\rho (x)𝑑x$$ Now since $`f_0+\epsilon h>0`$ and $`2p_0(x)1>0`$ in the range of integration, we have $$g^{\prime \prime }(\epsilon )_{\{x;p_0(x)>1/2\}}\frac{h^2}{(1(f_0+\epsilon h)^2)^2}(f_0+\epsilon h1)^2\rho (x)𝑑x>0$$ Therefore $`\epsilon g(\epsilon )`$ is concave, $`g^{\prime \prime }(0)<0`$, and $`f_0=f_K`$ is a maximum of the functional.$`\mathrm{}`$ c) Long run. We assume here that the player uses a Kelly betting strategy for a repetitive game with fuzzy advantage. For a given advantage distribution $`D_x`$ with $`p_0(x)>1/2`$, the expected exponential rate of growth after playing $`n`$ favorable hands with distributions $`x_1,x_2,\mathrm{}`$ is $$𝐄(G_n)=\frac{1}{n}\underset{i=1}{\overset{n}{}}(p\mathrm{log}(1+f(p_0(x_i)))+(1p)\mathrm{log}(1f(p_0(x_i))))D_{x_i}(p)𝑑p.$$ Thus the expected exponential rate of growth is $$\begin{array}{cc}\hfill 𝐄(G_n)& =(p\mathrm{log}(1+f(p_0(x)))+(1p)\mathrm{log}(1f(p_0(x))))D_x(p)𝑑p\rho (x)𝑑x\hfill \\ & =(p_0(x)\mathrm{log}(1+f(p_0(x)))+(1p)\mathrm{log}(1f(p_0(x))))\rho (x)𝑑x\hfill \\ & =𝐄_x(G_1)\rho (x)𝑑x\hfill \\ & =𝐄(G_1)\hfill \end{array}$$ Observe also that by independence of the hands we have $$\mathrm{𝐕𝐚𝐫}(G_n)=\frac{1}{n}\mathrm{𝐕𝐚𝐫}(G_1).$$ We extimate the standard deviation of $`G_n`$. For this we make the assumption that the distribution $`D_x(p)`$ is a distribution $`D(p)`$ with a small ”noise”, that is $$D_x(p)=D(p)+Z_xd(p)$$ where $`Z_x`$ is a random variable with $`0`$ expectation. We have $$\begin{array}{cc}\hfill Z_x\rho (x)𝑑x=0& \\ \hfill d(p)𝑑p=0& \end{array}$$ the second equation coming from $`D_x(p)𝑑p=1`$. According to the previous notation, the expected advantage is $$p_0(x)=pD_x(p)𝑑p$$ thus $$p_0(x)=p_0+AZ_x$$ where $`p_0=pD(p)𝑑p`$ and we denote $$A=pd(p)𝑑p.$$ Observe also that the quantity $`A`$ is directly related to the variance of $`p_0(x)`$ by $$\mathrm{𝐕𝐚𝐫}(p_0(x))=p_0(x)^2\rho (x)𝑑xp_0^2=A^2Z_x^2\rho (x)𝑑x.$$ Our purpose now is to compute the first order perturbation of the standard deviation of $`G_n`$ introduced by the noise $`Z_x`$. We have to the second order on $`Z_x`$ (we do not care about the coefficient of the first order), $$\begin{array}{cc}\hfill 𝐄_x(G_1)& =(p_0+AZ_x)\mathrm{log}(2p_0+2AZ_x)+(1p_0AZ_x)\mathrm{log}(22p_02AZ_x)\hfill \\ & =(p_0+AZ_x)(\mathrm{log}(2p_0)+\mathrm{log}(1+AZ_x/p_0))+\hfill \\ & +(1p_0AZ_x)(\mathrm{log}(22p_0)+\mathrm{log}(1AZ_x/(1p_0)))\hfill \\ & =p_0\mathrm{log}(2p_0)+(1p_0)\mathrm{log}(22p_0)+AZ_x+\frac{1}{2p_0(1p_0)}A^2Z_x^2+\mathrm{}\hfill \end{array}$$ Thus taking expected values with respect to $`x`$, $$𝐄(G_1)=p_0\mathrm{log}(2p_0)+(1p_0)\mathrm{log}(22p_0)+\frac{1}{2}\frac{\mathrm{𝐕𝐚𝐫}(p_0(x))}{p_0(1p_0)}+\mathrm{}$$ And, on the first order on $`\mathrm{𝐕𝐚𝐫}(p_0(x))`$, we have $$𝐄(G_1)^2=(p_0\mathrm{log}(2p_0)+(1p_0)\mathrm{log}(22p_0))^2+(p_0\mathrm{log}(2p_0)+(1p_0)\mathrm{log}(22p_0))\frac{\mathrm{𝐕𝐚𝐫}(p_0(x))}{p_0(1p_0)}+\mathrm{}$$ Now we compute the expansion of $$𝐄(G_1^2)=𝐄_x(G_1^2)\rho (x)𝑑x.$$ We have (after a long computation) $$\begin{array}{cc}\hfill 𝐄_x(G_1^2)& =p_0(x)(\mathrm{log}(2p_0(x)))^2+(1p_0(x))(\mathrm{log}(22p_0(x)))^2\hfill \\ & =\mathrm{}\hfill \\ & =p_0(\mathrm{log}(2p_0))^2+(1p_0)(\mathrm{log}(22p_0))^2+AZ_x+\hfill \\ & +\frac{1}{p_0(1p_0)}\left(1+(1p_0)\mathrm{log}(2p_0)+p_0\mathrm{log}(22p_0)\right)A^2Z_x^2+\mathrm{}\hfill \end{array}$$ Thus $$\begin{array}{cc}\hfill 𝐄(G_1^2)& =p_0(\mathrm{log}(2p_0))^2+(1p_0)(\mathrm{log}(22p_0(x)))^2+\hfill \\ & +\frac{1}{p_0(1p_0)}\left(1+(1p_0)\mathrm{log}(2p_0)+p_0\mathrm{log}(22p_0)\right)\mathrm{𝐕𝐚𝐫}(p_0(x))+\mathrm{}\hfill \end{array}$$ Finally, putting together the previous formulas, we have Theorem 3.4. The variance of the exponential rate of growth in a Kelly betting strategy with fuzzy advantage is, at the first order in the noise, $$\mathrm{𝐕𝐚𝐫}(G_1)=p_0(1p_0)\left(\mathrm{log}\left(\frac{p_0}{1p_0}\right)\right)^2+\frac{1}{p_0(1p_0)}\left(1(2p_01)\mathrm{log}\left(\frac{p_0}{1p_0}\right)\right)\mathrm{𝐕𝐚𝐫}(p_0(x))+\mathrm{}$$ As an appication to blackjack, where we have $`p_0=1/2+\epsilon `$ with $`\epsilon 10^2`$, and $`\mathrm{𝐕𝐚𝐫}(p_0(x))=\epsilon \sigma _{BET}`$ (where $`\sigma _{BET}`$ is the standard deviation of the true count), $$\mathrm{𝐕𝐚𝐫}(G_1)=4(1+\sigma _{BET}^2)\epsilon ^2+\mathrm{}$$ thus $$\sigma (G_n)=\frac{2\sqrt{1+\sigma _{BET}^2}}{\sqrt{n}}\epsilon +\mathrm{}$$ Thus the long run is increased by a factor $`1+\sigma _{BET}^2`$ caused by the standard deviation effect. In order to carry out precise estimates on the ”long run” we state a precise definition. Note that the factor $`2`$ is somewhat arbitrary. We may want to choose another value in order to have a bound on the probability of loosing after achieving the long run $`N`$ using Tchebichev’s inequality. Definition 3.5 (Long run). The long run $`N`$ is the minimal number of hands in order to have $$\frac{𝐄(G_N)}{\sigma (G_N)}2.$$ Thus from the above computation we have (disregarding integer parts), $$N=2^2(1+\sigma _{BET}^2)\epsilon ^2(1+\sigma _{BET}^2).40000$$ Typically the difference between $`\sigma _{BET}^2`$ between first and third base could be of the order of $`2\%`$ which makes the long run about $`2\%`$ longer, or about $`800`$ more favorables hands to be played, thus about $`2000`$ more hands to be played (that is at least $`40`$ more hours of play). 4) Practical gambling. a) Comparison of first and third base and head-on play. a.1) Full table. The number of cards removed from the deck between the betting decision and the play decision of the first base is exactly $`8\times 2=16`$. The average number of cards for the third base can be computed considering that the average number of cards in a hand of players playing basic strategy is 2.6 (this number is well known). Thus in average the number of cards removed from the deck between the betting decision and the play decision is 19.6. Denoting by $`\mathrm{\Sigma }(1)`$ and $`\mathrm{\Sigma }(2)`$ the corresponding standard deviations of the true count, and considering the approximation $`N>>n`$, we have $$\frac{\mathrm{\Sigma }(2)}{\mathrm{\Sigma }(1)}\frac{4.42}{4}=1.10$$ Thus $`\mathrm{\Sigma }(2)`$ is in average about $`10\%`$ higher than $`\mathrm{\Sigma }(1)`$. In company of ploppys. A recurrent topic in the blackjack literature is about is the influence of inaccurate players (ploppys) in the same table. It is an easy escape gate to blame others for your losses. The systematic answer that one founds all over the blackjack literature is that this is pure superstition. We believe also that there is mostly exaggeration in these complaints, but as we explain next, there is some mathematical foundation based on the standard deviation effect. In the first part of this section, when we were comparing the standard deviation of the first and third base, in order to estimate the number of cards played between first and third base, we made a major assumption: Other players at the table are playing basic strategy, thus the average number of cards in a hand is 2.6. Regular players report that very few players know basic strategy (less than $`5\%`$ we were told) and among these maybe about $`1/4`$ of them are more or less skilled card counters. Playing basic strategy, the casino has an edge of about $`1\%`$ ($`0.5\%`$ with best rules), but the actual win rate of blackjack in a casino is estimated to be between $`3\%`$ and $`7\%`$ (depending of the sources, see \[Gr2\] p.137). This shows how inaccurately the average player plays blackjack. The most common mistake is to think that the game consists in approaching as accurately as possible a total sum of $`21`$. This is totally false. Doing so one risks being busted, thus loosing the whole bet without getting the chance of letting the dealer bust. The real goal is to beat the dealer. And the best way to achieve this is to let him bust. Also ploppys like to split tens for example. This is something at which strategy players look with horror because the count has to be sky high to justify this play (nevertheless it can happen, see \[Us\] p.17 where J. Uston explain how he, on one occasion, he did split eight times tens in one hand…with an initial bet of $ 1000). This mistakes make that the average player uses more cards in his hands than the ideal basic strategy player. We do not have precise statistics on this, but one can guess that having a gang of ploppys spliting tens to your right can make that if you sit in third base the average number of cards played between the betting and the play decision could go as high as $`25`$. In such a table the standard deviation of third base will be $`25\%`$ higher than the one for first base. Thus this will induce longer long run and larger fluctuations, and you will be right to blame the ploppys. a.2) Head-on play. The differences between the standard deviation effect when one plays head-on (that is alone with the dealer) and one plays with a full table is even more important. The exact number of cards played in between decisions for a head-on play is $`4`$. Thus if you sit in first base at a full table your standard deviation will be twice as important (100 % more). And if you sit in third base this figure goes up to $`121\%`$ more. b) Absolute magnitude of standard deviation. The standard deviation true count formula allows us to compute $`\sigma _n`$ from $`\sigma _1`$. We compute now the standard deviation $`\sigma _1`$. It not only depends on $`N`$, and $`R`$, and the system count but also in the actual composition of the remaining deck. Proposition 1. We have $$\sigma _1=\frac{1}{N1}\sqrt{\underset{w}{}w^2\frac{l_w}{N}\left(\frac{R}{N}\right)^2}.$$ Proof. We have $$\begin{array}{cc}\hfill \sigma _1^2& =\underset{w}{}\left(\frac{R+w}{N1}\frac{R}{N}\right)^2\frac{l_w}{N}\hfill \\ & =\frac{1}{N^2(N1)^2}\underset{w}{}(Nw+R)^2\frac{l_w}{N}\hfill \\ & =\frac{1}{N^2(N1)^2}\left[N\left(\underset{w}{}w^2l_w\right)+2R\underset{w}{}wl_w+R^2\right]\hfill \\ & =\frac{1}{N^2(N1)^2}\left[N\left(\underset{w}{}w^2l_w\right)R^2\right]\hfill \\ & =\frac{1}{(N1)^2}\left[N\left(\underset{w}{}w^2l_w\right)\left(\frac{R}{N}\right)^2\right]\hfill \end{array}$$ $`\mathrm{}`$ Let $`𝒮`$ denote the balanced counting system used. A main characteristic of the counting system is the standard deviation $`\mathrm{\Sigma }_0(𝒮)`$ of the weights. Since the system is balanced we have $$\mathrm{\Sigma }_0=\mathrm{\Sigma }_0(𝒮)=\underset{w}{}w^2\frac{S_w}{52}$$ where $`S_w`$ is the total number of cards of weight $`w`$ in a complete deck. There follows the standard deviation of different systems round up to the third digit (for a description the different systems with their list of weights the reader can consult \[Sc\] p.62). $``$Uston ace-five count: $`\mathrm{\Sigma }_0=0.392`$ $``$Hi-Opt I: $`\mathrm{\Sigma }_0=0.784`$ $``$C-R Point count: $`\mathrm{\Sigma }_0=0.855`$ $``$Canfield expert count: $`\mathrm{\Sigma }_0=0.877`$ $``$Hi-Lo: $`\mathrm{\Sigma }_0=0.877`$ $``$Uston advanced plus-minus count: $`\mathrm{\Sigma }_0=0.877`$ $``$Halves count: $`\mathrm{\Sigma }_0=0.920`$ $``$The systematic count: $`\mathrm{\Sigma }_0=0.961`$ $``$Hi-Opt II: $`\mathrm{\Sigma }_0=1.468`$ $``$Canfield master count: $`\mathrm{\Sigma }_0=1.569`$ $``$Zen count: $`\mathrm{\Sigma }_0=1.569`$ $``$Uston advanced point count: $`\mathrm{\Sigma }_0=1.687`$ $``$The Revere point count: $`\mathrm{\Sigma }_0=1.710`$ $``$Uston SS count: $`\mathrm{\Sigma }_0=1.797`$ $``$The Revere advanced point count: $`\mathrm{\Sigma }_0=2.449`$ $``$Griffin seven count: $`\mathrm{\Sigma }_0=3.234`$ $``$Thorp ultimate count: $`\mathrm{\Sigma }_0=5.798`$ We observe, as expected, that $`\mathrm{\Sigma }_0`$ is larger for higher level (i.e. more precise) count systems. From the formula in proposition 1 we can get a good approximation of $`\sigma _1`$ in real gambling situations. Observe that for shoe games $`N`$ is always large due to non complete penetration (typically for shoe games we always have $`N50`$). Thus the square of the standard deviation in the remaining deck, that is $$\underset{w}{}w^2\frac{l_w}{N},$$ is well approximated by $`\mathrm{\Sigma }_0^2`$. Thus the square of the true count $`(R/N)^2`$ can be neglected (remember that the units we use to compute the number of cards is card units and not deck units as is done in practice). We conclude that the quantity $$\sigma _1^{}=\frac{\mathrm{\Sigma }_0}{N}$$ is a good approximation of $`\sigma _1`$ provided that the shoe is not totally depleted (note that in the approximation we replace $`N1`$ by $`N`$ which is acceptable). In practice one measures $`N`$ and the true count in deck units, then using the previous formula one will get $`\sigma _1^{}`$ in deck units that is the standard deviation of the true count in deck units which is the quantity that is used to determine the advantage in percentage. We can also approximate $`\sigma _n`$: $$\sigma _n=\sqrt{\frac{N1}{Nn}}\sqrt{n}\sigma _1\sqrt{n}\frac{\mathrm{\Sigma }_0}{N}$$ where we did approximate $`\sqrt{(N1)/(Nn)}`$ by $`1`$, which is quite accurate (except maybe if just one deck remains, then the value of $`\sigma _n`$ will be larger than the approximation). Using this approximation one can compute the following tables for standard deviation related to the bet decision ($`\sigma _{BET}`$) and the one for the play decision ($`\sigma _{PLAY}`$), for different players and different types of games. \| Position \| $`1`$ \| $`4`$ \| $`7`$ \| \| --- \| --- \| --- \| --- \| \| $`\sigma _{BET}`$ \| $`0.877`$ \| $`0.925`$ \| $`0.971`$ \| \| $`\sigma _{PLAY}`$ \| $`0.449`$ \| $`0.340`$ \| $`0.085`$ \| Hi-Lo, Canfield expert, Uston Advanced plus-minus ($`\mathrm{\Sigma }_0=0.877`$) $`8`$-Deck shoe, $`50\%`$ shoe played \| Position \| $`1`$ \| $`4`$ \| $`7`$ \| \| --- \| --- \| --- \| --- \| \| $`\sigma _{BET}`$ \| $`1.316`$ \| $`1.388`$ \| $`1.457`$ \| \| $`\sigma _{PLAY}`$ \| $`0.674`$ \| $`0.510`$ \| $`0.128`$ \| Hi-Lo, Canfield expert, Uston Advanced plus-minus ($`\mathrm{\Sigma }_0=0.877`$) $`8`$-Deck shoe, $`66.6\%`$ shoe played \| Position \| $`1`$ \| $`4`$ \| $`7`$ \| \| --- \| --- \| --- \| --- \| \| $`\sigma _{BET}`$ \| $`1.754`$ \| $`1.850`$ \| $`1.942`$ \| \| $`\sigma _{PLAY}`$ \| $`0.898`$ \| $`0.680`$ \| $`0.170`$ \| Hi-Lo, Canfield expert, Uston Advanced plus-minus ($`\mathrm{\Sigma }_0=0.877`$) $`8`$-Deck shoe, $`75\%`$ shoe played \| Position \| $`1`$ \| $`4`$ \| $`7`$ \| \| --- \| --- \| --- \| --- \| \| $`\sigma _{BET}`$ \| $`5.780`$ \| $`6.096`$ \| $`6.400`$ \| \| $`\sigma _{PLAY}`$ \| $`2.959`$ \| $`2.241`$ \| $`0.560`$ \| Thorp ultimate ($`\mathrm{\Sigma }_0=5.798`$) $`8`$-Deck shoe, $`50\%`$ shoe played \| Position \| $`1`$ \| $`4`$ \| $`7`$ \| \| --- \| --- \| --- \| --- \| \| $`\sigma _{BET}`$ \| $`8.670`$ \| $`9.144`$ \| $`9.600`$ \| \| $`\sigma _{PLAY}`$ \| $`4.439`$ \| $`3.362`$ \| $`0.840`$ \| Thorp ultimate ($`\mathrm{\Sigma }_0=5.798`$) $`8`$-Deck shoe, $`66.6\%`$ shoe played \| Position \| $`1`$ \| $`4`$ \| $`7`$ \| \| --- \| --- \| --- \| --- \| \| $`\sigma _{BET}`$ \| $`11.56`$ \| $`12.192`$ \| $`12.8`$ \| \| $`\sigma _{PLAY}`$ \| $`5.918`$ \| $`4.482`$ \| $`1.120`$ \| Thorp ultimate ($`\mathrm{\Sigma }_0=5.798`$) $`8`$-Deck shoe, $`75\%`$ shoe played A common feature is that the fluctuations become much larger at the end of the deck. And there is where most of the action takes place because it is at this moment than the true count uses to take its largest values ! This picture is scaring for the card counter and explains why blackjack is a game with high fluctuations for card counters. Let’s see what Stanford Wong and others have to say on this subject in his reference book \[Wo\] (p.199): Comments on risks. Peter Giles said: ”The late Ken Uston was once quoted in the ”Review Journal” as saying, ”It’s really tough to make a living at blackjack. The fluctuations will really wipe out the average guy. If I had to play by myself (instead of on a team), I probably wouldn’t be in it now”. You can quote me (S. Wong) as saying the same (…) I am still trying to determine how many units one is safe with. What is recommended in books is, in my opinion, too risky. This is one area in which it is hard to trust mathematics. Probably not all the mathematics can be found in the classical litterature. It is surprising to note that there is no complete treatment of the long run in the blackjack litterature as the one we carry out in section 3. To limit risks some card counters use to divide their true count by the number of players at the table (as Sonia, a professional gambler, does \[So\]), thus limiting the action, the fluctuation, … but also the advantage. The standard deviation tends to infinite when the deck depletes completely. We can see this formally, making $`N0`$ in the formulas (in a continuous model $`N1`$ is replaced by $`N`$). Of course $`N`$ takes integer values but we can compute for example that if only $`1/4`$ of a deck remains and we assume that $`\mathrm{\Sigma }_01`$ for the count system used, then $$\sigma _14$$ thus in the best case (playing head-on), $$\mathrm{\Sigma }(1)8.$$ And in the case of a full table, third base placement, $$\mathrm{\Sigma }(2)17.7,$$ a true count of the order of 17 is certainly like $`+\mathrm{}`$. When one sees these figures one wonders how it is possible that removing just one card from the deck could have such an effect in the standard deviation. Let’s compute exactly $`\sigma _1`$ in one particular situation. We assume that in a head-on game $`13`$ cards remain (that is $`1/4`$ of a deck). We assume that we use Hi-Lo and the the composing of the remaining cards is $`5`$ high cards, $`5`$ low cards and $`3`$ medium cards (thus $`R=0`$ at this point). The standard deviation of the true count after removing just one card is (in deck units) $$\sigma _1=52.\sqrt{\left(\frac{1}{12}\right)^2\frac{5}{13}+\left(\frac{1}{12}\right)^2\frac{5}{13}+0}=52.\sqrt{\frac{5}{936}}3.80$$ And if still one is not convinced, just look at what happens when we reveal one card, and this card is for example a low card. The new true count (in deck units) will be $$52.\frac{R}{N}=\frac{52}{12}4.33$$ That means that in such situations the true count indicator is meaningless and the standard deviation effect takes over. In an evil scenario casinos could exploit that weakness of card counters. They could offer and advertise very good penetration in order to get tables crowded with card counters. Then deal the shoes almost to bottom to exploit the huge standard deviation effect. Cardcounters will then experience too large fluctuations that will bring them into an important risk of ruin. Thus the casino will wipe out bankrolls of some unlucky card counters. It is unlikely that the casino will make any money in the operation, due to the increased benefits of the survival card counters. Bibliography \[Br\] L. BREIMAN, Optimal gambling systems for favorable games, Fourth Berkeley Symposium on Probability and Statistics, I, p. 65-78, 1961. \[Gr1\] P.A. GRIFFIN, The theory of blackjack, Huntington Press, Las Vegas, 1999. \[Gr2\] P.A. GRIFFIN, Extra stuff, Huntington Press, Las Vegas, 1991. \[HC\] L. HUMBLE, C. COOPER, The world’s greatest blackjack book, Doubleday, New York, 1980. \[JM\] A. JALIB M’HALL, The true count theorem, Message posted in rec.gambling.blackjack, 7/30/1996. \[RT\] L.M. ROTANDO, E.O. THORP, The Kelly criterion and the stock market, Amer. Math. Monthly, 99, p.922-931, 1992. \[Sc\] F. SCOBLETE, Best blackjack, Bonus Books, Chicago, 1996. \[So\] SONIA, Personal communication, 1/2000. \[Th\] E. O. THORP, Beat the dealer, Vintage books, New York, 1962. \[Wo\] S. WONG, Professional blackjack, Pi Yee Press, 1994.
warning/0006/cond-mat0006192.html	ar5iv	text	# Multi-mode dynamics of a coupled ultracold atomic-molecular system ## Abstract We analyze the coherent multi-mode dynamics of a system of coupled atomic and molecular Bose gases. Starting from an atomic Bose-Einstein condensate with a small thermal component, we observe a complete depletion of the atomic and molecular condensate modes on a short time scale due to significant population of excited states. Giant coherent oscillations between the two condensates for typical parameters are almost completely suppressed. Our results cast serious doubts on the common use of the 2-mode model for description of coupled ultracold atomic-molecular systems and should be considered when planning future experiments with ultracold molecules. After experimental achievement of quantum degeneracy in atomic gases of bosons and fermions , several leading groups started the quest for creation of an analogous state in a gas of ultracold molecules . Such systems present a potentially much more challenging experimental task as there are several new degrees of freedom to be controlled. A major step towards this goal has been made by the Austin group of Heinzen , who photoassociated atoms in an atomic condensate into a selected internal molecular state producing for the first time molecules in the nK regime. Cold molecules are believed to be created as a transient state in experiments with Feshbach resonances as well . In this novel area there are only a few theoretical works analyzing the process of making ultracold molecules in atomic condensates and predicting several interesting phenomena in the mixed atomic-molecular systems. Javanainen et al. in a series of papers analyzed the efficiency of photoassociation of an atomic condensate into its molecular counterpart using various theoretical schemes. Others predicted that the ground state of the hybrid system would have a soliton-like nature with liquid-like properties performing Josephson-like oscillations in response to a sudden variation of a magnetic field . Most frequently, coherent oscillations between the atomic and molecular condensates are envisaged . All the approaches employed so far have used a 2-mode approximation (only the condensates involved) to describe the dynamics of a coupled ultracold atomic-molecular system. In this Letter we demonstrate that the 2-mode approach is inadequate as a description of current experiments on stimulated production of cold molecules in atomic Bose-Einstein condensates. For typical parameters the predicted oscillations between the two condensates are strongly damped due to significant population of excited atomic and molecular modes leading to a complete depletion of the initial condensate on a short time scale. The method we use, quantum-optical in spirit, can be regarded as a generalization of the classic Bogolubov approximation . The second-quantized Hamiltionian for the hybrid atomic-molecular system confined in a box with periodic boundary conditions may be written in the following form: $`H`$ $`=`$ $`{\displaystyle \mathrm{d}^3r\left[\mathrm{\Phi }^{}\frac{p^2}{2m}\mathrm{\Phi }+\mathrm{\Psi }^{}\frac{p^2}{4m}\mathrm{\Psi }\right]}`$ (1) $`+`$ $`\sqrt{V}\mathrm{}\mathrm{\Omega }{\displaystyle \mathrm{d}^3r\left[\mathrm{\Psi }^{}\mathrm{\Phi }^2+\mathrm{\Psi }\mathrm{\Phi }^2\right]}`$ (2) $`+`$ $`{\displaystyle \frac{V\mathrm{}g}{2}}{\displaystyle \mathrm{d}^3r\mathrm{\Phi }^{}\mathrm{\Phi }^{}\mathrm{\Phi }\mathrm{\Phi }},`$ (3) where $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are atomic and molecular field operators, respectively, $`V=L^3`$ is a volume of the system ($`L`$ being a size of the box), $`\mathrm{\Omega }`$ parameterizes a coupling between the two fields and $`g=\frac{4\pi \mathrm{}a}{mV}`$ characterizes the atom-atom interactions in the low-energy, s-wave approximation ($`a`$ being the scattering length and $`m`$ – the mass of the atom; in fact, $`a`$ is slightly changed in external, e.g. optical, fields – in present experiments this correction is small though ). A similar Hamiltonian is used to describe a process of second-harmonic generation in nonlinear optics . The atom-molecule and molecule-molecule collisions are not included as their parameters in the low-energy regime are unknown. However, as can be seen from the construction of the method, their incorporation could be easily accomplished. Note that in this model a molecule is created when the positions of two atoms exactly coincide (for an example of a finite-range coupling approach see ). The fields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are expanded in natural modes of the system – the plane waves: $`\mathrm{\Phi }(𝐫)={\displaystyle \frac{1}{\sqrt{V}}}{\displaystyle \underset{𝐤}{}}\mathrm{exp}(i𝐤𝐫)a_𝐤,`$ (4) $`\mathrm{\Psi }(𝐫)={\displaystyle \frac{1}{\sqrt{V}}}{\displaystyle \underset{𝐤}{}}\mathrm{exp}(i𝐤𝐫)b_𝐤,`$ (5) where $`a_𝐤`$ and $`b_𝐤`$ are bosonic annihilation operators for atoms and molecules, respectively, and $`𝐤=\frac{2\pi }{L}𝐧`$ with $`n_i=0,\pm 1,\pm 2,\mathrm{}`$ ($`i=x,y,z`$). With this substitution the Hamiltonian assumes its final form: $`{\displaystyle \frac{H}{\mathrm{}}}`$ $`=`$ $`\xi {\displaystyle \underset{𝐤}{}}n^2(a_𝐤^{}a_𝐤+{\displaystyle \frac{1}{2}}b_𝐤^{}b_𝐤)`$ (6) $`+`$ $`\mathrm{\Omega }{\displaystyle \underset{𝐤,𝐤^{}}{}}b_{𝐤+𝐤^{}}^{}a_𝐤a_𝐤^{}+b_{𝐤+𝐤^{}}a_𝐤^{}a_𝐤^{}^{}+`$ (7) $`+`$ $`{\displaystyle \frac{1}{2}}g{\displaystyle \underset{𝐤,𝐤^{},𝐤^{\prime \prime }}{}}a_{𝐤+𝐤^{}𝐤^{\prime \prime }}^{}a_{𝐤^{\prime \prime }}^{}a_𝐤^{}a_𝐤,`$ (8) where $`\xi =\frac{\mathrm{}}{2m}(\frac{2\pi }{L})^2`$. After elimination of a fast time dependence with the substitution $`\{a,b\}_𝐤=\mathrm{exp}(i\xi n^2t)\{\alpha ,\beta \}_𝐤`$, the Heisenberg equations of motion for the operators $`\alpha _𝐤`$ and $`\beta _𝐤`$ acquire the following form: $`\dot{\alpha _𝐤}=2i\mathrm{\Omega }{\displaystyle \underset{𝐤^{}}{}}\mathrm{exp}({\displaystyle \frac{1}{2}}i\xi \|𝐧𝐧^{}\|^2t)\beta _{𝐤+𝐤^{}}\alpha _𝐤^{}^{}`$ (9) $`ig{\displaystyle \underset{𝐤^{},𝐤^{\prime \prime }}{}}\mathrm{exp}\left[2i\xi (𝐧𝐧^{})(𝐧𝐧^{\prime \prime })t\right]\alpha _{𝐤^{}+𝐤^{\prime \prime }𝐤}^{}\alpha _𝐤^{}\alpha _{𝐤^{\prime \prime }},`$ (10) $`\dot{\beta _𝐤}=i\mathrm{\Omega }{\displaystyle \underset{𝐤^{}}{}}\mathrm{exp}({\displaystyle \frac{1}{2}}i\xi \|𝐧2𝐧^{}\|^2t)\alpha _{𝐤𝐤^{}}\alpha _𝐤^{}.`$ (11) The first hint about limitations of the 2-mode model comes from the following argument. As initially only the atomic condensate (the $`𝐤=0`$ mode) is populated, one might naively suspect that the coupling would primarily lead to an interconversion between atoms and molecules . In such a case only the atomic and molecular condensate ($`𝐤=0`$) modes would be macroscopically populated and therefore we replace the corresponding operators ($`\alpha _0`$ and $`\beta _0`$) by c-numbers (classical complex fields) and set all the others to zero ($`\alpha _{𝐤0}=0`$ and $`\beta _{𝐤0}=0`$). Then, in the absence of atomic collisions, time evolution of the amplitudes can be calculated analytically. Assuming $`\alpha _0=r\mathrm{exp}(i\varphi )`$ and $`\beta _0=\rho \mathrm{exp}(i\theta )`$ with $`\varphi =\varphi _0=const`$ and $`\theta =\theta _0=const`$, the solution is: $`\alpha _0(t)`$ $`=`$ $`{\displaystyle \frac{\sqrt{N}}{\mathrm{cosh}(\sqrt{2N}\mathrm{\Omega }t)}}\mathrm{exp}(i\varphi _0),`$ (12) $`\beta _0(t)`$ $`=`$ $`i\sqrt{{\displaystyle \frac{N}{2}}}\mathrm{tanh}(\sqrt{2N}\mathrm{\Omega }t)\mathrm{exp}(2i\varphi _0),`$ (13) $`\theta _0`$ $`=`$ $`2\varphi _0{\displaystyle \frac{\pi }{2}},`$ (14) where $`N`$, the total number of atoms, is a conserved quantity ($`N=\|\alpha _0(t)\|^2+2\|\beta _0(t)\|^2`$) – see Fig. 1. Note that the inclusion of interactions makes the 2-mode problem analytically insoluble (in particular the condition of constant phases $`\theta `$ and $`\varphi `$ cannot be fulfilled any more). Numerical solutions are of oscillatory character – different from (12). They depict an interconversion between atomic and molecular condensate modes. In the next step we calculate quantum corrections to such a 2-mode model treating the condensates ($`\alpha _0`$ and $`\beta _0`$) as the only sources of particles. This amounts to leaving in Eqs.(9) the terms with at least one $`𝐤=0`$ (in the form of (12)) and neglecting the other ones. In the resulting equations only the zero-momentum as well as the $`𝐤`$ and $`𝐤`$ atomic and molecular modes are present which allows to solve the coupled linear operator equations numerically. Their asymptotics ($`t\mathrm{}`$) may be investigated analytically though, yielding the following expectations values for the number operators: $`\beta _𝐤^{}\beta _𝐤`$ $`=`$ $`const,`$ (15) $`\alpha _𝐤^{}\alpha _𝐤`$ $`=`$ $`{\displaystyle \frac{N\mathrm{\Omega }^2}{2\lambda ^2}}\left[\mathrm{exp}(2\lambda t)+\mathrm{exp}(2\lambda t)2\right],`$ (16) where $`\lambda =\sqrt{2N\mathrm{\Omega }^2\xi ^2n^4}`$. $`\lambda `$ is real for all modes with $`n<\sqrt[4]{2N\mathrm{\Omega }^2/\xi ^2}`$, which sets the number of modes whose population grows in time. As typically both $`\mathrm{\Omega }`$ and $`\xi `$ are of the order of $`1010^2`$ Hz while $`N10^510^6`$ , corrections to the 2-mode model are divergent for many low-lying states. With the parameters used it is only for the 11-th and higher excited atomic modes that the quantum corrections are small and oscillatory (imaginary $`\lambda `$). Therefore, one is not allowed to exclude excited modes ($`𝐤0`$) from a theoretical model . Fig. 1 presents an example of the run-away correction to the 2-mode model obtained numerically together with the classical source terms $`\alpha _0(t)`$ and $`\beta _0(t)`$ (note: in all plots the total populations are normalized by the total number of particles $`N`$). From Fig. 1 one immediately concludes that the 2-mode solution (12) is physically invalid for times larger than $`0.6`$ ms. To cure this problem, from now on, we will use a multi-mode model in the form of Eqs. (9). Solution of the nonlinear operator equations (9) presents an extremely difficult task. A semiclassical approximation, however, is well justified for all except extremely low temperatures. Therefore, we replace all operators by c-number complex amplitudes. From the viewpoint of the Bogolubov method , such an approach is legitimate as indeed many modes are macroscopically populated (i.e. their occupation is greater than quantum fluctuations). This way we are left with a set of nonlinear differential equations which must be solved numerically. The first observation in the multi-mode model is that if one starts from a pure atomic condensate (the $`𝐤=0`$ mode), the 2-mode dynamics is recovered. However, even a very small occupation of excited atomic or molecular modes results in the dynamics beyond the 2-mode approach. Such a behavior resembles an initiation of superfluorescence where quantum, not thermal, fluctuations play a role. In a typical experiment, roughly $`85`$% of the total number of atoms populate the atomic condensate whereas the rest of them is thermally distributed over excited atomic modes and we mimic such a situation in the initial conditions of our model . All the molecular modes are initially unpopulated. Each atomic mode is assigned an initial, randomly chosen, phase. Any subsequent dynamics depends on the initial phases and, in a sense, a single simulation describes a single experimental realization. Values of the parameters in the model are $`\mathrm{\Omega }=`$18.432 Hz, $`\xi =`$71.373 Hz, $`N=10^5`$ and $`g=`$0.018 Hz (the atomic mass and the scattering length are those of <sup>87</sup>Rb and the size of the box is equal to the Thomas-Fermi radius of a condensate of $`N`$ atoms in a trap with frequency of $`\omega _0=2\pi \mathrm{\hspace{0.33em}80}`$ Hz – see ). Sample results for the model with 2622 modes (maximum atomic and molecular excitation $`n_i=5,\mathrm{},5`$, $`i=x,y,z`$) are presented in Fig. 2. In fact the results stabilize for the number of modes exceeding 1000. A striking discovery is that after a relatively short time both the atomic and molecular condensates (the $`𝐤=0`$ modes) are completely depleted and all particles occupy excited modes in roughly equal proportions. For the used (typical) parameters, only one oscillation in the condensate population survives – see Fig. 3. This clearly indicates that if one wants to convert an atomic condensate to its molecular counterpart, it is necessary to precisely tailor the length of the coupling pulse so that the maximum of the molecular $`𝐤=0`$ mode is picked. If one investigates the condensate ($`𝐤=0`$) population and a sum of populations of all excited modes ($`𝐤0`$) solely (like in Fig. 2), their dependence on the initial phases is negligible in a sufficiently big model (we checked this fact by setting different initial phases in our calculations). The effect is due to self-averaging caused by very many random phases present in the system. However, it can be observed in the time evolution of single modes. We emphasize that, contrary to the previous treatments , the effective losses from the condensates are completely due to a Hamiltonian evolution and not because of any phenomenologically introduced loss processes. In other words, due to an external weak coupling of many degrees of freedom, the system is effectively heated and its final equilibrium state certainly does not result from any $`T=0`$ dynamics. The inclusion of atom-molecule and molecule-molecule collisions would randomize the still coherent dynamics leading to an analogous and even more pronounced effect. Remarkably, the character of the dynamics does not depend on the total number of particles – our simulations for $`N=210^3`$ and $`N=510^5`$ still show only one oscillation in the atomic condensate population on a slightly altered (longer for smaller $`N`$) time scale. The amplitude of this oscillation is bigger for smaller $`N`$ indicating a small increase in the molecular condensate production ($`30\%`$ for $`N=210^3`$ and $`23\%`$ for $`N=510^5`$ instead of $`25\%`$ for $`N=10^5`$ as seen in Fig. 2). In order to recover the standard oscillatory dynamics of the 2-mode limit within a multi-mode model, one needs to detune the excited modes. The latter can be achieved by decreasing the size of the box (and so enlarging a spacing between the excited modes) while keeping the total density fixed (in order not to alter the scattering and coupling parameters). The results for the box whose volume is $`15^3`$ times smaller are presented in Fig. 4 (the appropriately decreased number of particles is $`N=10^5/15^330`$). To summarize, we point out the invalidity of a 2-mode model for a description of present experiments with coupled atomic-molecular systems in the Bose-Einstein condensation regime. We find it necessary to employ a multi-mode approach and within it we observe a complete depletion of both the atomic and molecular condensates on a short time scale. For typical parameters, only one giant oscillation between the two condensates is present. Thus, a new destructive mechanism in the system is pointed out which should be taken into account while planning future experiments. An effectively 2-mode dynamics is recovered under special conditions of large detunings. The system presents an interesting example of a quantum coherent dynamical evolution effectively randomized by the coupling of many degrees of freedom. In this Letter we presented the calculations for the box rather than the harmonic oscillator potential. We have chosen the box size such that the first excitation energy is of the order of the level spacing in the experiment of Heinzen . Hence, we may expect that for the harmonic oscillator the time of coherent evolution would be even shorter, because of the quadratic vs. linear dependence of excitation energies in the respective potentials. K.G. and M.G. acknowledge support by Polish KBN grant no 2 P03B 057 15. K.R. and K.G. are supported by the subsidy of the Foundation for Polish Science. Part of the results has been obtained using computers at the Interdisciplinary Centre for Mathematical and Computational Modeling (ICM) at Warsaw University.
warning/0006/math-ph0006016.html	ar5iv	text	# 1 Introduction ## 1 Introduction The von Kármán plate theory is governed by two coupled nonlinear fourth-order partial differential equations in three independent variables (Cartesian coordinates on the plate middle-plane $`x^1,x^2`$ and the time $`x^3`$) and two dependent variables (the transversal displacement function $`w`$ and Airy’s stress function $`\Phi `$), namely \| $`D\mathrm{\Delta }^2w\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }w_{,\alpha \beta }\Phi _{,\mu \nu }+\rho w_{,33}=0,`$ \| \| --- \| \| $`\left(1/Eh\right)\mathrm{\Delta }^2\Phi +\left(1/2\right)\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }w_{,\alpha \beta }w_{,\mu \nu }=0,`$ \| (1) where $`\mathrm{\Delta }`$ is the Laplace operator with respect to $`x^1`$ and $`x^2`$, $`D=Eh^3/12(1\nu ^2)`$ is the bending rigidity, $`E`$ is Young’s modulus, $`\nu `$ is Poisson’s ratio, $`h`$ is the thickness of the plate, $`\rho `$ is the mass per unit area of the plate middle-plane, $`\delta ^{\alpha \beta }`$ is the Kronecker delta symbol and $`\epsilon ^{\alpha \beta }`$ is the alternating symbol. Here and throughout the work: Greek (Latin) indices range over 1, 2 (1, 2, 3), unless explicitly stated otherwise; the usual summation convention over a repeated index is used and subscripts after a comma at a certain function $`f`$ denote its partial derivatives, that is $`f_{,i}=f/x^i,f_{,ij}=f/x^ix^j`$, etc. The von Kármán equations (1) describe entirely the motion of a plate, the membrane stress tensor $`N^{\alpha \beta }`$, moment tensor $`M^{\alpha \beta }`$, shear-force vector $`Q^\alpha `$, strain tensor $`E^{\alpha \beta }`$ and bending tensor $`K_{\alpha \beta }`$ being given in terms of $`w`$ and $`\Phi `$ through the following expressions: \| $`N^{\alpha \beta }=\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }\Phi _{,\mu \nu },M^{\alpha \beta }=D\left\{(1\nu )\delta ^{\alpha \mu }\delta ^{\beta \nu }+\nu \delta ^{\alpha \beta }\delta ^{\mu \nu }\right\}w_{,\mu \nu },`$ \| \| --- \| \| $`Q^\alpha =M_{,\mu }^{\alpha \mu }+N^{\alpha \mu }w_{,\mu },E^{\alpha \beta }=(1/Eh)\left\{(1+\nu )\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }v\delta ^{\alpha \beta }\delta ^{\mu \nu }\right\}\Phi _{,\mu \nu },K_{\alpha \beta }=w_{,\alpha \beta }.`$ \| The theory under consideration allows an exact variational formulation, the von Kármán equations being the Euler-Lagrange equations associated with the action functional $$I[w,\Phi ]=L𝑑x^1𝑑x^2𝑑x^3,L=T\mathrm{\Pi }$$ (2) where \| $`\mathrm{\Pi }`$ \| $`=`$ \| $`(D/2)\left\{\left(\mathrm{\Delta }w\right)^2\left(1\nu \right)\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }w_{,\alpha \beta }w_{,\mu \nu }\right\}`$ \| \| --- \| --- \| --- \| \| \| $``$ \| $`(1/2Eh)\left\{\left(\mathrm{\Delta }\Phi \right)^2\left(1+\nu \right)\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }\Phi _{,\alpha \beta }\Phi _{,\mu \nu }\right\}+(1/2)\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }\Phi _{,\alpha \beta }w_{,\mu }w_{,\nu },`$ \| is the strain energy per unit area of the plate middle-plane and $$T=\left(\rho /2\right)\left(w_{,3}\right)^2,$$ is the kinetic energy per unit area of the plate middle-plane. ## 2 Conservation laws In the recent paper , all Lie point symmetries of system (1) are shown to be variational symmetries of the functional (2), and all corresponding (via Noether’s theorem) conservation laws admitted by the smooth solutions of the von Kármán equations are established. Each such conservation law is a linear combination of the basic linearly independent conservation laws $$\frac{\Psi _{(j)}}{x^3}+\frac{P_{(j)}^\mu }{x^\mu }=0\left(j=1,2,\mathrm{},14\right),$$ whose densities $`\Psi _{(j)}`$ and fluxes $`P_{(j)}^\mu `$ are presented (together with the generators of the respective symmetries) on the Table 1 below in terms of $`Q^\alpha ,M^{\alpha \beta },G^{\alpha \beta }`$ and $`F^\alpha `$, $$G^{\alpha \beta }=\left(1/Eh\right)\left\{(1+\nu )\delta ^{\alpha \mu }\delta ^{\beta \nu }\nu \delta ^{\alpha \beta }\delta ^{\mu \nu }\right\}\Phi _{,\mu \nu }\left(1/2\right)\epsilon ^{\alpha \mu }\epsilon ^{\beta \nu }w_{,\mu }w_{,\nu },F^\alpha =G_{,\nu }^{\alpha \nu }.$$ \| $`w`$ \- translations \| transversal linear momentum (first von Kármán eqn) \| \| --- \| --- \| \| $`X_1=\frac{}{w}`$ \| $`P_{(1)}^\alpha =Q^\alpha ,\Psi _{(1)}=\rho w_{,3}`$ \| \| $`\Phi `$ \- translations \| compatibility condition (second von Kármán eqn) \| \| --- \| --- \| \| $`X_{14}=\frac{}{\Phi }`$ \| $`P_{(14)}^\alpha =F^\alpha ,\Psi _{(14)}=0`$ \| time - translations energy $`X_4=\frac{}{x^3}`$ $`P_{(4)}^\alpha =w_{,3}Q^\alpha \Phi _{,3}F^\alpha +w_{,3\beta }M^{\alpha \beta }+\Phi _{,3\beta }G^{\alpha \beta }`$ $`\Psi _{(4)}=T+\mathrm{\Pi }`$ $`x^1\&x^2`$\- translations wave momentum $`X_2=\frac{}{x^1}`$ $`P_{(2)}^\alpha =\delta ^{\alpha 1}L+w_{,1}Q^\alpha +\Phi _{,1}F^\alpha w_{,1\beta }M^{\alpha \beta }\Phi _{,1\beta }G^{\alpha \beta }`$ $`\Psi _{(2)}=\rho w_{,1}w_{,3}`$ $`X_3=\frac{}{x^2}`$ $`P_{(3)}^\alpha =\delta ^{\alpha 2}L+w_{,2}Q^\alpha +\Phi _{,2}F^\alpha w_{,2\beta }M^{\alpha \beta }\Phi _{,2\beta }G^{\alpha \beta }`$ $`\Psi _{(3)}=\rho w_{,2}w_{,3}`$ rotations moment of the wave momentum $`X_6=x^2\frac{}{x^1}x^1\frac{}{x^2}`$ $`P_{(6)}^\alpha =x^2P_{(2)}^\alpha x^1P_{(3)}^\alpha +\epsilon _\nu ^\mu w_{,\mu }M^{\alpha \nu }+\epsilon _\nu ^\mu \Phi _{,\mu }G^{\alpha \nu }`$ $`\Psi _{(6)}=x^2\Psi _{(2)}x^1\Psi _{(3)}`$ rigid body rotations angular momentum $`X_7=x^1\frac{}{w}`$ $`P_{(7)}^\alpha =M^{\alpha 1}x^1Q^\alpha +w\epsilon ^{\alpha \nu }\Phi _{,\nu 2},\Psi _{(7)}=\rho x^1w_{,3}`$ $`X_8=x^2\frac{}{w}`$ $`P_{(8)}^\alpha =M^{\alpha 2}x^2Q^\alpha +w\epsilon ^{\nu \alpha }\Phi _{,\nu 1},\Psi _{(8)}=\rho x^2w_{,3}`$ scaling $`X_5=x^\mu \frac{}{x^\mu }+2x^3\frac{}{x^3}`$ $`P_{(5)}^\alpha =x^1P_{(2)}^\alpha +x^2P_{(3)}^\alpha 2x^3P_{(4)}^\alpha w_{,\beta }M^{\alpha \beta }\Phi _{,\beta }G^{\alpha \beta }`$ $`\Psi _{(5)}=x^1\Psi _{(2)}+x^2\Psi _{(3)}2x^3\Psi _{(4)}`$ Galilean boost center-of-mass theorem $`X_9=x^3\frac{}{w}`$ $`P_{(9)}^\alpha =x^3Q^\alpha ,\Psi _{(9)}=\rho (x^3w_{,3}w)`$ ## 3 Balance laws Given a region $`\mathrm{\Omega }`$ in the plate middle-plane with sufficiently smooth boundary $`\mathrm{\Sigma }`$ of outward unit normal $`n_\alpha `$, a balance law $$\frac{d}{dt}\underset{\mathrm{\Omega }}{}\Psi _{(j)}𝑑x^1𝑑x^2+\underset{\mathrm{\Sigma }}{}P_{(j)}^\alpha n_\alpha 𝑑\mathrm{\Sigma }=0,$$ (3) corresponds to each of the conservation laws listed in Table 1. It holds, just as the respective conservation law, for every smooth solution of the von Kármán equations. The balance laws are applicable even if $`\mathrm{\Omega }`$ is intersected by a discontinuity (singular) manifold (on which the corresponding densities $`\Psi _{(j)}`$ and fluxes $`P_{(j)}^\alpha `$ may suffer jump discontinuities) provided that the integrals exist. We are ready now to extend the “continuous“ von Kármán plate theory so as to cover situations when some physical quantities suffer jump discontinuities at a certain curve. ## 4 Acceleration waves ###### Definition 1 A discontinuity solution of the von Kármán equations is a couple of functions $`(w,\Phi )`$, defined in a certain region $`\mathrm{\Omega }`$, such that the two balance laws corresponding to the von Kármán equations themselves, namely $$\frac{d}{dt}\underset{\stackrel{~}{\mathrm{\Omega }}}{}\rho w_{,3}𝑑x^1𝑑x^2\underset{\stackrel{~}{\mathrm{\Sigma }}}{}Q^\alpha \stackrel{~}{n}_\alpha 𝑑\mathrm{\Sigma }=0,\underset{\stackrel{~}{\mathrm{\Sigma }}}{}F^\alpha \stackrel{~}{n}_\alpha 𝑑\mathrm{\Sigma }=0,$$ (4) hold $`\stackrel{~}{\mathrm{\Omega }}\mathrm{\Omega }`$ with boundary $`\stackrel{~}{\mathrm{\Sigma }}`$ of outward unit normal $`\stackrel{~}{n}_\alpha `$, and $`(w,\Phi )`$ is a solution of the (local) von Kármán equations (1) almost everywhere in $`\mathrm{\Omega }`$ except for a moving curve $`\mathrm{\Gamma }`$ at which some of the derivatives of $`w`$ or $`\Phi `$ have jumps. ###### Definition 2 A discontinuity solution of the von Kármán equations is an acceleration wave if at the wave front – a smoothly propagating connected singular curve $`\mathrm{\Gamma }`$ $$\mathrm{\Gamma }:\gamma (x^1,x^2,x^3)=0,(x^1,x^2)\mathrm{\Omega }𝐑^2,x^3𝐑^+,\gamma C^1(\mathrm{\Omega }\times 𝐑^+),$$ we have $$[w]=[\Phi ]=[w_{,i}]=[\Phi _{,i}]=0,[w_{,33}]0.$$ (5) (Here and in what follows, the square brackets are used to denote the jump of any field $`f`$ across the curve $`\mathrm{\Gamma }`$, i.e., $`[f]=f_2f_1,`$ where $`f_2`$ and $`f_1`$ are the limit values of $`f`$ behind $`\mathrm{\Gamma }`$ and ahead of $`\mathrm{\Gamma }`$.) The moving curve $`\mathrm{\Gamma }`$ divides the region $`\mathrm{\Omega }`$ into two parts $`\mathrm{\Omega }^+`$ and $`\mathrm{\Omega }^{}`$ and forms the common border between them. It is assumed that ahead of the wave front (in the region $`\mathrm{\Omega }^+`$) we have the known unperturbed fields $`w^+(x^1,x^2,x^3)`$, $`\Phi ^+(x^1,x^2,x^3)`$ and behind it (in the region $`\mathrm{\Omega }^{}`$) – the unknown perturbed fields $`w^{}(x^1,x^2,x^3)`$, $`\Phi ^{}(x^1,x^2,x^3)`$. At the wave front $`\mathrm{\Gamma }`$, we have the jump conditions (5). The jumps of the derivatives of $`w`$ and $`\Phi `$ across $`\mathrm{\Gamma }`$ are permissible if they obey the compatibility conditions following by Hadamard’s lemma . Thus ###### Proposition 1 If $`[w_{,33}]0`$, then $$[w_{,\alpha \beta }]=\lambda n_\alpha n_\beta ,[w_{,\alpha 3}]=\lambda Cn_\alpha ,[w_{,33}]=\lambda C^2,$$ where $`\lambda `$ is an arbitrary factor, $`C`$ and $`n_\alpha `$, $$C=\left\|\gamma \right\|^1\gamma /x^3,n_\alpha =\left\|\gamma \right\|^1\gamma /x^\alpha ,\left\|\gamma \right\|=\sqrt{(\gamma /x^1)^2+(\gamma /x^2)^2},$$ are the speed of displacement and the direction of propagation of the wave front $`\mathrm{\Gamma }`$. ###### Proposition 2 If at least one of the third derivatives of $`w`$ suffers a jump at $`\mathrm{\Gamma }`$, then the compatibility conditions for the jumps of the third derivatives of the displacement field across $`\mathrm{\Gamma }`$ are: \| $`\left[w_{,\alpha \beta \gamma }\right]`$ \| $`=`$ \| $`\lambda ^{}n_\alpha n_\beta n_\gamma +\lambda /s\left(n_\alpha n_\beta t_\gamma +n_\alpha t_\beta n_\gamma +t_\alpha n_\beta n_\gamma \right)`$ \| \| --- \| --- \| --- \| \| \| $`+`$ \| $`\lambda a\left(t_\alpha t_\beta n_\gamma +t_\alpha n_\beta t_\gamma +n_\alpha t_\beta t_\gamma \right),`$ \| where $`t_\alpha `$ is the unit tangent vector to $`\mathrm{\Gamma }`$, $`\lambda ^{}`$ is an arbitrary factor, while $`\lambda =\left[w_{,\alpha \beta }\right]n^\alpha n^\beta `$ and $`a=t_\alpha n^\alpha /s`$, $`s`$ being the natural parameter (arc-length) of the curve $`\mathrm{\Gamma }`$. ###### Proposition 3 If at least one of the second derivatives of $`\Phi `$ suffers a jump at $`\mathrm{\Gamma }`$, then $$[\Phi _{,\alpha \beta }]=\mu n_\alpha n_\beta ,[\Phi _{,\alpha 3}]=\mu Cn_\alpha ,[\Phi _{,33}]=\mu C^2,$$ where $`\mu `$ is an arbitrary factor, are the compatibility conditions for the jumps of the second derivatives of the stress field across $`\mathrm{\Gamma }`$. ###### Proposition 4 If at least one of the third derivatives of $`\Phi `$ suffer a jump at $`\mathrm{\Gamma }`$, then the compatibility conditions for the jumps of the third derivatives of the stress field across $`\mathrm{\Gamma }`$ are: \| $`\left[\Phi _{,\alpha \beta \gamma }\right]`$ \| $`=`$ \| $`\mu ^{}n_\alpha n_\beta n_\gamma +\mu /s\left(n_\alpha n_\beta t_\gamma +n_\alpha t_\beta n_\gamma +t_\alpha n_\beta n_\gamma \right)`$ \| \| --- \| --- \| --- \| \| \| $`+`$ \| $`\mu a\left(t_\alpha t_\beta n_\gamma +t_\alpha n_\beta t_\gamma +n_\alpha t_\beta t_\gamma \right),`$ \| where $`\mu =[\Phi _{,\alpha \beta }]n^\alpha n^\beta `$ and $`\mu ^{}`$ is an arbitrary factor. According to the divergence theorem (see e.g. ), a couple of functions $`(w,\Phi )`$ suffering jump discontinuities at a singular curve $`\mathrm{\Gamma }`$ is a discontinuity solution of the von Kármán equations in the sense of Definition 1 iff the following jump conditions $$C\left[\rho w_{,3}\right]+\left[Q^\alpha \right]n_\alpha =0,\left[F^\alpha \right]n_\alpha =0,$$ (6) hold at $`\mathrm{\Gamma }`$, and a balance law of form (3) holds on this solution iff at $`\mathrm{\Gamma }`$: $$C\left[\Psi _{\left(j\right)}\right]\left[P_{\left(j\right)}^\alpha \right]n_\alpha =0.$$ (7) Definition 2, Propositions 1, 2, 3, 4 and jump conditions (6) imply that: ###### Proposition 5 If an acceleration wave in the von Kármán plate theory is such that $`[w_{,\alpha \beta \gamma }]0`$ ($`[\Phi _{,\alpha \beta \gamma }]0`$) at the curve of discontinuity $`\mathrm{\Gamma }`$, then $`\lambda ^{}=\lambda a`$ ($`\mu ^{}=\mu a`$). Given a discontinuity solution of the von Kármán equations, the two corresponding balance laws (4) being satisfied, the other balance laws do not necessarily hold for this solution. The jump conditions associated with the most important conservation laws from Table 1 are derived using (7) and presented on the Table 2 below, where $`X_j^+(w_{,\beta })`$ and $`X_j^+(\Phi _{,\beta })`$ denote the limit values of $`X_j(w_{,\beta })`$ and $`X_j(\Phi _{,\beta })`$ ahead of $`\mathrm{\Gamma }`$: time - translations energy $`X_4=\frac{}{x^3}`$ $`\frac{C}{2}\left(D\lambda ^2\frac{\mu ^2}{Eh}\right)=\left(D\lambda X_4^+(w_{,\beta })\frac{\mu }{Eh}X_4^+(\Phi _{,\beta })\right)n^\beta `$ $`x^1`$ & $`x^2`$\- translations wave momentum $`X_2=\frac{}{x^1}`$ $`\frac{1}{2}\left(D\lambda ^2\frac{\mu ^2}{Eh}\right)n^1=\left(D\lambda X_2^+(w_{,\beta })\frac{\mu }{Eh}X_2^+(\Phi _{,\beta })\right)n^\beta `$ $`X_3=\frac{}{x^2}`$ $`\frac{1}{2}\left(D\lambda ^2\frac{\mu ^2}{Eh}\right)n^2=\left(D\lambda X_3^+(w_{,\beta })\frac{\mu }{Eh}X_3^+(\Phi _{,\beta })\right)n^\beta `$ rotations moment of the wave momentum $`X_6=x^2\frac{}{x^1}x^1\frac{}{x^2}`$ $`\frac{1}{2}\epsilon _\beta ^\alpha n_\alpha x^\beta \left(D\lambda ^2\frac{\mu ^2}{Eh}\right)=`$ $`\left\{D\lambda \left(\epsilon _\beta ^\alpha w_{,\alpha }+X_6^+(w_{,\beta })\right)\frac{\mu }{Eh}\left(\epsilon _\beta ^\alpha \Phi _{,\alpha }+X_6^+(\Phi _{,\beta })\right)\right\}n^\beta `$ scaling $`X_5=x^\mu \frac{}{x^\mu }+2x^3\frac{}{x^3}`$ $`\frac{1}{2}\left(x^\alpha n_\alpha 2Cx^3\right)\left(D\lambda ^2\frac{\mu ^2}{Eh}\right)=`$ $`\left\{D\lambda \left(w_{,\beta }+X_5^+(w_{,\beta })\right)\frac{\mu }{Eh}\left(\Phi _{,\beta }+X_5^+(\Phi _{,\beta })\right)\right\}n^\beta `$ The center-of-mass theorem holds for any discontinuity solutions of the von Kármán equations. The balance laws, associated with the infinitesimal symmetries $`X_7`$, $`X_8`$, $`X_{10}`$ and $`X_{11}`$ hold iff $`\lambda =0`$, while those associated with $`X_{12}`$ and $`X_{13}`$ – iff $`\mu =0`$. For this reason, there do not exist acceleration waves in the von Kármán plate theory satisfying all balance laws. Obviously, when dealing with discontinuity solutions, from physical point of view it seems reasonably that at least the balance of energy should hold in addition to the balance laws corresponding to the fundamental equations considered. Observing Table 2 it is evident that for acceleration waves propagating into an undisturbed plate the balance of energy implies also the balances of wave momentum, moment of the wave momentum as well as the balance related to the scaling symmetry. ## 5 Examples As an example, we consider acceleration waves such that behind and ahead of the wave front the plate motion is described by solutions of the von Kármán equations invariant under the group generated by $`X_3`$ and $`X_2+(1/c)X_4`$, where $`c`$ is an arbitrary constant. The most general form of such group-invariant solutions is \| $`w=u(\xi )=u_0+u_1\xi +u_2\mathrm{sin}\omega \xi +u_3\mathrm{cos}\omega \xi ,`$ \| \| --- \| \| $`\Phi =\phi (\xi )=\phi _0+\phi _1\xi +\phi _2\xi ^2+\phi _3\xi ^3,`$ \| where $`\xi =x^1cx^3`$, $`u_j`$, $`\phi _j`$ are arbitrary constants, and $`\omega =c\sqrt{\rho /D}`$. Let ($`u^+`$,$`\phi ^+`$) – an arbitrary solution of that kind – describes the plate motion ahead of the wave front. Then, Definition 2 and Propositions 1 to 4 imply that each acceleration wave of the type considered reads $$u=\{\begin{array}{cc}u^++c_1(1\mathrm{cos}\omega \xi ),\hfill & \xi <0,\hfill \\ u^+,\hfill & \xi >0,\hfill \end{array}\phi =\{\begin{array}{cc}\phi ^++c_2\xi ^2,\hfill & \xi <0,\hfill \\ \phi ^+,\hfill & \xi >0,\hfill \end{array}$$ (8) where $`c_1`$ and $`c_2`$ are arbitrary constants, but $`c_10`$; the wave front in this case is the moving straight line $`\mathrm{\Gamma }:\xi =0`$. In general, however, an acceleration wave of form (8) does not satisfy the balance laws other than (4). Indeed, after a little manipulation, the jump conditions from Table 2 simplify to $`DEh\omega ^4c_1(c_12u_3^+)`$ $`=`$ $`4c_2(c_2+2\phi _2^+),`$ (9) $`DEh\omega ^2c_1(u_1^++\omega u_2^+)`$ $`=`$ $`2c_2\phi _1^+,`$ (10) where $`u_j^+`$, $`\phi _j^+`$ are the constants in $`u^+`$ and $`\phi ^+`$, respectively. The jump condition (9) is necessary and sufficient for the balances of energy, wave momentum and moment of wave momentum to hold, while the balance related to $`X_5`$ requires both (9) and (10). The second relation is treated in a different manner according to the acceleration wave under consideration. If the wave is such that $`c_2=0`$, then (10) holds for this wave only if $`u_1^+=\omega u_2^+`$. On the other hand, if we consider waves with $`c_20`$, then choosing the coefficient $`\phi _1^+`$ in a suitable manner, we could satisfy (10) identically (note that adding a linear function of the independent variables to Airy’s stress function does not change the membrane stress tensor $`N^{\alpha \beta }`$). Hence, the balance associated with the scaling symmetry $`(X_5)`$ holds only for acceleration waves of form (8) satisfying (9) and which are such that either $`c_20`$ or $`u_1^+=\omega u_2^+`$. Another example, discussed in details in , is an axisymmetrically expanding acceleration wave composed by solutions of the von Kármán equations that are joined invariants of the rotation ($`X_6`$) and scaling ($`X_5`$) symmetries. Acknowledgments. This work was supported by Contract No. MM 517/1995 with NSF, Bulgaria.
warning/0006/cond-mat0006065.html	ar5iv	text	# Variety and Volatility in Financial Markets ## I INTRODUCTION In recent years physicists started to interact with economists to concur to the modeling of financial markets as model complex systems . This triggered the interest of a group of physicists into the analysis and modeling of price dynamics in financial markets performed by using paradigms and tools of statistical and theoretical physics . One target of these researches is to implement a stochastic model of price dynamics in financial markets which reproduces the statistical properties observed in the time evolution of stock prices. In the last few years physicists interested in financial analysis have performed several empirical researches investigating the statistical properties of stock price and volatility time series of a single stock (or of an index) at different temporal horizons . Such a kind of analysis does not take into account any interaction of the considered financial stock with other stocks traded simultaneously in the same market. It is known that the synchronous price returns time series of different stocks are pair correlated and several researches has been performed also by physicists in order to extract information from the correlation properties . A precise characterization of collective movements in a financial market is of key importance in understanding the market dynamics and in controlling the risk associated to a portfolio of stocks. The present study contributes to the understanding of collective behavior of a portfolio of stocks in normal and extreme days of market activity. Specifically, we address the question: Is the complexity of a financial market essentially limited to the statistical behavior of each financial time series or rather a complexity of the overall market exists? To answer this question, we present the results of an empirical analysis performed adopting the following point of view. We investigate the price returns of an ensemble of $`n`$ stocks simultaneously traded in a financial market at a given day. With this approach we quantify what we call the variety of a financial market at a given trading day . The variety provides statistical information about the amount of different behavior observed in stock return in a given ensemble of stocks at a given trading time horizon (in the present case, one trading day). We observe that the distribution of variety is sensitive to the composition of the portfolio investigated (especially to the capitalization of the considered stocks). The return distribution shows a typical shape for most of the trading days. However, the typical behavior is not observed during crash and rally days. The shape and parameters characterizing the ensemble return distribution are relatively stable during normal phases of the market activity while become time dependent in the periods subsequent to crashes. The variety is characterized by a long-range correlated memory showing that no typical time scale can be expected after a rally or a crash for the expected relaxation to a “normal” market phase. Moreover a simple model such as the single-index model is not able to reproduce the statistical properties empirically observed. The paper is organized as follow. In Section II we illustrate our database and the ensemble of stocks considered. Sect. III is devoted to the investigation of the statistical properties of the time evolution of each single stock. In Section IV, we discuss the statistical properties of ensemble return distribution. Specifically we consider the behavior of the central lowest moments, their distribution and correlation, a comparison of time and portfolio average, and the role of the size and homogeneity of the investigated portfolio. In Section V we compare the statistical properties observed in a real financial market with the prediction of the single-index model. In Section VI we present a discussion of the obtained results. ## II DATABASE AND INVESTIGATED VARIABLES The investigated market is the New York Stock Exchange (NYSE) during the 12-year period from January 1987 to December 1998 which corresponds to 3032 trading days. We consider the ensemble of all stocks traded in the NYSE. The number of stocks traded in the NYSE is increasing in the investigated period and it ranges from $`1128`$ at the beginning of 1987 to $`2788`$ at the end of 1998. The total number of data records exceeds $`6`$ millions. The variable investigated in our analysis is the daily price return, which is defined as $$R_i(t)\frac{Y_i(t+1)Y_i(t)}{Y_i(t)},$$ (1) where $`Y_i(t)`$ is the closure price of $`i`$th stock at day $`t`$ ($`t=1,2,..`$). For each trading day $`t`$, we consider $`n`$ returns, where $`n`$ is depending on the total number of stocks traded in the NYSE at the selected day $`t`$. In our study we use a “market time”. With this choice, we consider only the trading days and we remove the weekends and the holidays from the calendar time. A database of more than 6 millions records unavoidably contains some errors. A direct control of a so large database is not realistic. For this reason, to avoid spurious results we filter the data by not considering price returns which are in absolute values greater than $`50\%`$. The companies traded in the NYSE are quite different the one from the other. Differences among the companies are observed both with respect to the sector of their economic interests and with respect to their size. One measure of the size of a company is its capitalization. The capitalization of a stock is the stock price times the number of outstanding shares. In this study, we discuss the role of the different capitalization in the price dynamics. ## III SINGLE stock PROPERTIES The distribution of returns with different time horizons of a single stock or index has been studied by several authors . The stocks traded in a financial market have different capitalization. An important point is whether the differences in capitalization are reflected in the statistical properties of the price returns of the stocks. To answer this question we investigate the distribution of daily returns of $`2188`$ stocks traded in the NYSE at an arbitrarily chosen day that we select as June 10th, 1996. We compare the statistical properties of daily price return distribution of each stock as a function of its capitalization. We order the $`2188`$ stocks in decreasing order according to their capitalization at June 10th, 1996. Our ordering procedure gives to the most capitalized stock (the General Electric Co., GE) the rank $`i=1`$, to the second one (the Coca Cola Company) the rank $`i=2`$, and so on. An analysis of the return probability density function (pdf) for the $`2188`$ stocks shows that the distributions are different. This is due in general to: (i) different scale and (ii) different shape of the return pdfs. In order to eliminate one source of difference we analyze the pdf of the normalized returns $`(R_i(t)\mu _i)/\sigma _i`$ $`(i=1,2,\mathrm{},2188)`$, where $`\mu _i`$ and $`\sigma _i`$ are the first two central moments of the time series $`R_i(t)`$ defined as $`\mu _i={\displaystyle \frac{1}{T_i}}{\displaystyle \underset{t=1}{\overset{T_i}{}}}R_i(t),`$ (2) $`\sigma _i=\sqrt{{\displaystyle \frac{1}{T_i}}\left({\displaystyle \underset{t=1}{\overset{T_i}{}}}(R_i(t)\mu _i)^2\right)},`$ (3) where $`T_i`$ is the number of trading days of the stock $`i`$ during the investigated period. The quantity $`\mu _i`$ gives a measure of the overall performance of stock $`i`$ in the period. The standard deviation $`\sigma _i`$ is called historical volatility in the financial literature and quantifies the risk associated with the $`i`$-th stock. This quantity is of primary importance in risk management and in option pricing. The pdf of normalized daily returns of all the stocks ordered by capitalization is shown in Fig. 1. The central part of the distribution of the most capitalized stocks has a bell-shaped profile. Moving towards less capitalized stocks the central part of the distribution becomes more peaked and the tails of the distribution become fatter. The pdf of the less capitalized stocks is therefore more leptokurtic than the pdf of the more capitalized ones. The typical estimation of the degree of leptokurtosis of a pdf is done by considering its kurtosis. The evaluation of the kurtosis of the pdf is in general difficult for small set of data because the fourth moment and all the moments higher than the second are extremely sensible to the highest absolute returns. This implies that the kurtosis calculated from a relatively small set of records is dominated by the highest absolute returns rather than by the shape of the pdf and therefore it is not a good statistical estimation. To avoid this problem, we quantify the distance between the empirically calculated pdf of daily returns of $`i`$th stock and the Gaussian distribution by considering the quantity $$h\frac{<\|x\|>}{\sqrt{<x^2><x>^2}}.$$ (4) The quantity $`h`$ is nondimensional and depends on the first two moments. For the Gaussian distribution $$P_G(x)=\frac{1}{\sqrt{2\pi \sigma _G^2}}\mathrm{exp}\left(\frac{(x\mu _G)^2}{2\sigma _G^2}\right),$$ (5) the parameter $`h`$ is equal to $$h_G=\sqrt{\frac{2}{\pi }}\left(\mathrm{exp}(\frac{\mu _G^2}{2\sigma _G^2})+\sqrt{\frac{\pi }{2}}\frac{\mu _G}{\sigma _G}Erf\left(\frac{\mu _G}{\sqrt{2}\sigma _G}\right)\right).$$ (6) The parameter $`h_G`$ is a function of the ratio $`\mu _G/\sigma _G`$ ranging from the lower bound $`\sqrt{2/\pi }`$ when $`\mu _G/\sigma _G=0`$ to infinity. For a leptokurtic pdf, as for example a Laplace distribution or a Student’s t-distribution with finite variance, $`h`$ is always smaller than $`h_G`$. The distance of $`h`$ from $`h_G`$ is able to quantify the degree of leptokurtosis of the considered pdf. Figure 2 shows the parameter $`h`$ for the stocks traded in the NYSE as a function of their capitalization. In the figure, we show also the lower bound of $`h_G`$ for comparison. The empirically calculated parameter $`h`$ is systematically smaller than $`h_G`$. The mean value $`<h>`$ of the overall market is $`<h>=0.67`$ and its standard deviation is $`\sigma _h=0.06`$. Hence this result suggests that as a first approximation one can assume that the large majority of stocks are characterized by a roughly similar pdf. However we wish to point out that this conclusion is only valid as a first approximation because a trend of $`h`$ is clearly detected in Fig. 2. Specifically $`h`$ increases as the capitalization increases. Therefore the less capitalized stocks have a more leptokurtic daily return pdf than the more capitalized ones. The second moment of return distribution has been found finite in recent research . In order to verify the convergence of the pdf towards a Gaussian pdf at large temporal horizons, we evaluate the $`h`$ parameter for weekly $`<h_w>`$ and monthly $`<h_m>`$ return pdfs. We obtain from our analysis $`<h_w>=0.70`$ and $`<h_m>=0.74`$. These results show that the values of $`h`$ moves towards $`h_G=\sqrt{2/\pi }0.80`$ when the time horizon of returns is increased, supporting the conclusion of finite second moment. ## IV ensemble return distribution In the previous section we focused on statistical properties of time evolution of price returns for each single stocks traded in the NYSE. In this section we perform a synchronous analysis on the return of all the stock traded in the NYSE. To this aim we extract the $`n`$ returns of the $`n`$ stocks for each trading day $`t`$. The distribution of these returns $`P_t(R)`$ provides information about the kind of activity occurring in the market at the selected trading day $`t`$. Figure 3 shows the logarithm of the pdf as a function of the return and of the trading day. In this figure we show the interval of daily returns from $`25\%`$ to $`25\%`$. The central part of the distribution is roughly triangular in a logarithmic scale and this shape and its scale are conserved for long time periods. Sometimes the shape and scale of the ensemble return pdf changes abruptly either in the presence of large average positive returns or large average negative returns. Figure 4 shows the same data of Fig. 3 in a contour plot. The contour lines describe equiprobability regions. In order to point out the properties of the central part of the distribution, in Fig. 4 we plot only the returns which are less than $`15\%`$ in absolute value. Only a few points of the contour lines fall behind this limit during the 1987 and 1998 crises. In Fig. 4 there are long time periods in which the central part of the distribution maintains his shape and the equiprobability contour lines are approximately parallel one to each other. As an example, one can consider the three-year period 1993-1995. On the other hand there are time periods in which the shape of the distribution changes drastically. In general these periods corresponds to financial turmoil in the market. For example a dramatic change of the shape and of the scale of the pdf is observed in Fig. 4 during and after the 19 Oct. 1987 crash, at the beginning of 1991 and at the end of 1998. A systematic analysis of the change of the shape and scale of the ensemble return distribution during extreme events of the market has been discussed elsewhere One key aspect of the ensemble return distribution concerns its shape during the normal periods of activity of the market. Is the distribution approximately Gaussian or systematic deviation from a Gaussian shape are quantitatively observed? We already cited that a direct inspection of Fig. 3 suggests that the central part of the empirical return distribution is roughly Laplacian (triangular in a logarithmic scale) and not Gaussian. To make this analysis more quantitative, we show in Fig. 5 the ratio between the value of $`h`$ determined for each trading day from the ensemble return distribution and the quantities $`h_G`$ calculated by determining the mean and the standard deviation of $`P_t(R)`$ and hypothesizing a Gaussian shape by using Eq. (6). The ratio $`h/h_G`$ is systematically smaller than one and this implies that the Gaussian hypothesis for the shape of the distribution is not verified by the empirical analysis. In other words the Gaussian distribution is not a good approximation both for the central part and for the tails of the distribution and the deviation from the Gaussian behavior is systematically observed for all the trading days of the 12 years time period analyzed in our study. In summary the ensemble return distribution well characterizes the market activity. It has a typical shape and scale during long periods of “normal” activity of the market characterized by moderately low average daily return. During extreme events the shape and scale are dramatically changed in a systematic way. Specifically during crises the ensemble return distribution becomes negatively skewed whereas during rallies a positive skewness is observed . Figure 4 clearly shows that extreme events (such as for example October 87 crash) triggers an “aftershock” period, in the ensemble return pdf, that can last for a period of time of several months. ### A Central moments In order to characterize more quantitatively the ensemble return distribution at day $`t`$, we extract the first two central moments at each of the $`3032`$ trading days. Specifically, we consider the average and the standard deviation defined as $`\mu (t)={\displaystyle \frac{1}{n_t}}{\displaystyle \underset{i=1}{\overset{n_t}{}}}R_i(t),`$ (7) $`\sigma (t)=\sqrt{{\displaystyle \frac{1}{n_t}}\left({\displaystyle \underset{i=1}{\overset{n_t}{}}}(R_i(t)\mu (t))^2\right)},`$ (8) where $`n_t`$ indicates the number of stocks traded at day $`t`$. The mean of price returns $`\mu (t)`$ quantifies the general trend of the market at day $`t`$. The standard deviation $`\sigma (t)`$ gives a measure of the width of the ensemble return distribution. We call this quantity variety of the ensemble because it gives a measure of the variety of behavior observed in a financial market at a given day. A large value of $`\sigma (t)`$ indicates that different companies are characterized by rather different returns at day $`t`$. In fact in days of high variety some companies perform great gains whereas others have great losses. The mean and the standard deviation of price returns are not constant and fluctuate in time. We study the temporal series of $`\mu (t)`$ and $`\sigma (t)`$ in order to characterize the temporal evolution of the ensemble return distribution quantitatively. We investigate these fluctuating parameters by investigating their time correlation properties and their pdfs. ### B Probability distributions of the central moments The empirical pdf of the mean $`\mu (t)`$ for the 3032 trading days investigated is shown in Fig. 6. The central part of this distribution is non-Gaussian and is roughly described by a Laplace distribution. The mean $`\mu (t)`$ is proportional to the sum of $`n`$ random variables $`R_i(t)`$ $`(i=1,2,\mathrm{},n)`$. The Central Limit Theorem prescribes that the sum of $`n`$ independent random variables with finite variance converges to a Gaussian pdf. By assuming a finite value for the volatility of stocks, the observation that the pdf of the mean return $`\mu (t)`$ is non-Gaussian can be therefore attributed to the presence of correlation between the stocks. Figure 7 shows the pdf of the variety $`\sigma (t)`$. The central part of this distribution is approximated by a lognormal distribution. A deviation from the lognormal behavior is observed in the tail of higher values of variety. This deviation is depending on the size of the portfolio and will be discussed in subsection IV E. ### C Correlations in the central moments Another important statistical property of $`\mu (t)`$ and $`\sigma (t)`$ concerns their correlation properties . For the considered portfolio, we calculate the autocorrelation function of a variable $`x(t)`$ which is defined as $$R(\tau )\frac{<x(t)x(t+\tau )><x(t)><x(t+\tau )>}{<x(t)^2><x(t)>^2}.$$ (9) In agreement with previous results , we find that the mean $`\mu (t)`$ is approximately delta correlated, whereas the autocorrelation function of $`\sigma (t)`$ is long-range correlated. The empirical autocorrelation function of $`\sigma (t)`$ is well approximated by a power-law function $`R(\tau )\tau ^\delta `$ . By performing a best fit with a maximum time lag of $`50`$ trading days, we determine the exponent $`\delta =0.230\pm 0.006`$. This result indicates that the variety $`\sigma (t)`$ has a long-time memory in the market. We recall that the historical volatility is characterized by long time memory of the same nature . Another way to investigate the long-range correlation is to determine the power spectrum of the investigated variable. We evaluate the power spectrum of $`\sigma (t)`$ and we perform a best fit of the power spectrum with a functional form of the kind $$S(f)\frac{1}{f^\eta }.$$ (10) Our best fit for the power spectrum of $`\sigma (t)`$ gives for the exponent $`\eta 1.1`$. This result confirms that the variety $`\sigma (t)`$ is a long-range time correlated random variable. ### D Time and portfolio average Figure 6 shows two curves. In fact in Fig. 6 we also show the pdf of the mean $`\mu _i`$. The quantity $`\mu _i`$ (see Eq. (2)) is the mean return of stock $`i`$ averaged over the investigated time interval. The pdf of $`\mu _i`$ is non-Gaussian and it is much more peaked than the pdf of $`\mu (t)`$. Hence the statistical behavior observed by investigating a large portfolio in a market day is not representative of the statistical behavior observed by investigating the time evolution of single stocks. This comparison can be performed also for the second moment of the distributions. In Fig. 7 we compare the pdf of the volatility $`\sigma _i`$ and the pdf of the variety $`\sigma (t)`$. Also in this case, the statistical properties of $`\sigma _i`$ and $`\sigma (t)`$ are different. Specifically, the pdf of $`\sigma (t)`$ is more peaked than the pdf of $`\sigma _i`$. In order to understand the different behavior of the time-averaged and the portfolio-averaged quantities, for the sake of simplicity, we consider a portfolio composed by $`N`$ stocks which are traded in a period of $`T`$ trading days. We first study the properties of the two means, $`\mu _i`$ and $`\mu (t)`$. It is straightforward to verify that $$<\mu _i>_i=<\mu (t)>_t\mu ,$$ (11) where $`<..>_t`$ indicates temporal average and $`<..>_i`$ indicates ensemble average. The variances of $`\mu _i`$ and $`\mu (t)`$ are in general different. We obtain for the variance of $`\mu (t)`$ the expression $$Var[\mu (t)]_t\frac{1}{T}\underset{t=1}{\overset{T}{}}(\mu (t)\mu )^2=\frac{1}{N^2}\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}\sigma _{ij}^2,$$ (12) where $`\sigma _{ij}^2`$ is the return covariance between stock $`i`$ and $`j`$ defined as $$\sigma _{ij}^2=<R_i(t)R_j(t)>_t<R_i(t)>_t<R_j(t)>_t.$$ (13) The width of the pdf of $`\mu (t)`$ (shown in Fig. 6) is the square root of $`Var[\mu (t)]_t`$. Equations (12) and (13) indicate that this quantity depends both on the ensemble averaged square volatility (terms with $`i=j`$ in Eq. (12)) and on the mean of the synchronous cross-covariances between pairs of stocks (terms with $`ij`$ in Eq. (12)). With similar methods we show that the variance of $`\mu _i`$ can be written as $$Var[\mu _i]_i\frac{1}{N}\underset{i=1}{\overset{N}{}}(\mu _i\mu )^2=\frac{1}{T^2}\underset{t=1}{\overset{T}{}}\underset{t^{}=1}{\overset{T}{}}\sigma _{tt^{}}^2,$$ (14) where we define the return covariance between trading day $`t`$ and $`t^{}`$ as $$\sigma _{tt^{}}^2=<R_i(t)R_i(t^{})>_i<R_i(t)>_i<R_i(t^{})>_i.$$ (15) This quantity gives an estimate of the correlation present in the whole portfolio at trading day $`t`$ and $`t^{}`$. The double sum in Eq. (14) can be split in a term depending on the average square variety ($`t=t^{}`$) and in a term depending on the correlation between different trading days ($`tt^{}`$). We verify that the average square variance and volatility satisfy the sum rule $$Var[\mu _i]_i+<\sigma _i^2>_i=Var[\mu (t)]_t+<\sigma ^2(t)>_t.$$ (16) Combining Eq.s (12), (14) and (16) we show that $`{\displaystyle \frac{T1}{T}}<\sigma ^2(t)>_t+{\displaystyle \frac{2}{N^2}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \underset{i<j}{}}\sigma _{ij}^2=`$ (17) $`={\displaystyle \frac{N1}{N}}<\sigma _i^2>_i+{\displaystyle \frac{2}{T^2}}{\displaystyle \underset{t=1}{\overset{T}{}}}{\displaystyle \underset{t^{}<t}{}}\sigma _{tt^{}}^2.`$ (18) Since $`N,T>>1`$, we approximate $`(N1)/N(T1)/T1`$ and Eq. (17) becomes $$<\sigma _i^2>_i<\sigma ^2(t)>_t<\sigma _{ij}^2>_{ij}<\sigma _{tt^{}}^2>_{tt^{}},$$ (19) or equivalently $$Var[\mu (t)]_tVar[\mu _i]_i<\sigma _{ij}^2>_{ij}<\sigma _{tt^{}}^2>_{tt^{}}.$$ (20) Figure 6 shows that $`Var[\mu (t)]_t>Var[\mu _i]_i`$. This empirical observation together with the last relation tell us that the synchronous cross-correlations between the stocks are on average stronger than the single stock correlation present in the whole portfolio at two different trading day. This result is consistent with previous observations that synchronous returns of different stocks are significantly cross-correlated , whereas single price returns are poorly autocorrelated in time. This conclusion is also verified by our empirical observation that $`<\sigma _i^2>_i><\sigma ^2(t)>_t`$. ### E Portfolio size One key aspect of the previous results concerns the degree of generality of the observed stylized facts. In other words, are the empirical properties of the variety depending on the considered portfolio? In Section II we have shown that all the stocks are not equivalent with respect to their statistical properties (see the spread of points observed in Fig. 2). In fact a trend is observed in the degree of non-Gaussian shape of the return distribution as a function of the stock capitalization. To test the degree of sensitivity of our results to the average capitalization of the selected portfolio, we repeat the analysis presented in subsection III.B for three other portfolios of stocks traded in the NYSE. Specifically we investigate: (a) the set of 30 stocks used to compute the Dow Jones Industrial Average index; (b) the set of stocks traded in the NYSE and used to compute the Standard & Poor’s 100 index; and (c) the set of stocks traded in the NYSE and used to compute the Standard & Poor’s 500 index. The results obtained for all the stocks traded in the NYSE are also considered for reference. The four sets are different with respect to two aspects. They differ for the number of stocks present in the set and for the average capitalization of the considered stocks. The empirical pdfs of $`\mu (t)`$ for the four considered sets are roughly the same. An evident different behavior is observed for the variety. In Fig. 8 we show the pdf of the variety of the considered portfolios of stocks. Specifically panels (a), (b), (c) and (d) of Fig. 8 are the results obtained for the Dow Jones 30, Standard & Poor’s 100, Standard & Poor’s 500 and NYSE sets of stocks, respectively. By moving from the smallest to the largest portfolio of stocks two effects take place. The pdf of the variety becomes progressively sharper and deviates more from a lognormal profile. The fact that the pdf of the variety becomes progressively sharper is probably due to the fact the number of elements in the considered set increases whereas we interpret the progressive deviation from the lognormal profile as a direct manifestation of the progressive increases of the degree of inhomogeneity of the portfolio of stocks. In summary the presence of inhomogeneity in capitalization in the portfolio of stocks affects the statistical properties of the variety of the portfolio. This fact should be kept in mind when results about the variety such as results about other statistical properties included return distribution are obtained by considering the statistical properties of a set of inhomogeneous stocks. ## V Single-index model In this section we compare the results of our empirical analysis obtained for the NYSE portfolio of stocks with the results obtained by modeling the stock price dynamics with the single-index model. The single-index model is a basic model of price dynamics in financial markets. It assumes that the returns of all stocks are controlled by one factor, usually called the “market”. In this model, for any stock $`i`$ we have $$R_i(t)=\alpha _i+\beta _iR_M(t)+ϵ_i(t),$$ (21) where $`R_i(t)`$ and $`R_M(t)`$ are the return of the stock $`i`$ and of the “market” at day $`t`$, respectively, $`\alpha _i`$ and $`\beta _i`$ are two real parameters and $`ϵ_i(t)`$ is a zero mean noise term characterized by a variance equal to $`\sigma _{ϵ_i}^2`$. The noise terms of different stocks are assumed to be uncorrelated, $`<ϵ_i(t)ϵ_j(t)>_t=0`$ for $`ij`$. Moreover the covariance between $`R_M(t)`$ and $`ϵ_i(t)`$ is set to zero for any $`i`$. Each stock is correlated with the market and the presence of such a correlation induces a correlation between any pair of stocks. It is customary to adopt a broad-based stock index for the market $`R_M(t)`$. Our choice for the “market” time series is the Standard and Poor’s 500 index. The best estimation of the model parameters $`\alpha _i`$, $`\beta _i`$ and $`\sigma _{ϵ_i}^2`$ is done with the ordinary least squares method . In order to compare our empirical results with those predicted by the single-index model we build up an artificial market according to Eq. (20). To this end we first evaluate the model parameters for all the stocks traded in the NYSE and then we generate a set of $`n`$ of surrogate time series according to Eq. (20). To make the simulation as realistic as possible, in the generation of our surrogate data set we use as “market” time series the true time series of the Standard and Poor’s 500 index. We evaluate the central moments $`\mu (t)`$ and $`\sigma (t)`$ defined in Eqs (7-8) for the surrogate data. In Fig. 9(a) we show the time series of $`\mu (t)`$ of the real data and in Fig. 9(b) we show the same quantity for the surrogate market data generated according to the single-index model. The agreement between the two time series is pretty high and therefore the single-index model describes quite well the mean returns of the market at time $`t`$ provided that the behavior of the “market” $`R_M(t)`$ is known . This result is also confirmed by Fig. 10 where the pdf of $`\mu (t)`$ for real and surrogate data are shown. Also the time correlation properties of surrogate $`\mu (t)`$ are pretty similar to the real ones. In fact, a fast decaying autocorrelation function of $`\mu (t)`$ is observed in surrogate data. A good agreement is also observed when one investigates the statistical properties of $`\mu _i`$ and $`\sigma _i`$. The single-index model approximates quite well the empirical distribution of $`\mu _i`$ and $`\sigma _i`$. A different behavior is observed for the variety $`\sigma (t)`$. Figure 9(c) and 9(d) show the time series of $`\sigma (t)`$ for real and surrogate data, respectively. The real time series of the variety is non stationary and shows several bursts of activity. On the contrary the surrogate time series is quite stationary with the exception of the 1987 crash. Figure 11 shows the pdfs of $`\sigma (t)`$ for real and surrogate data. The model fails in describing the distribution of $`\sigma (t)`$. In summary, the single-index model gives a good approximation of the statistical behavior of $`\mu (t)`$, $`\mu _i`$ and $`\sigma _i`$ whereas it describes poorly the statistical behavior of the variety of a portfolio of stocks traded in a financial market. This conclusion is also supported by the observation that the autocorrelation function of the surrogate variety decays in $`23`$ trading days to the value $`0.1`$ and the power spectrum is very similar to a white noise spectrum, whereas long-range correlation is observed in real data. A more refined analysis shows that the artificial ensemble return distribution is systematically less leptokurtic than the real one. Moreover, in Ref. we show that the single-index model is unable to predict the change in the symmetry properties of the ensemble return distribution in crash and rally days. The differences observed between the behavior of real data and the behavior of surrogate data suggest that the correlations among the stocks can be explained by the single-index model only for “normal” periods in first approximation whereas the model miss completely to reproduce the correlation behavior during extreme events. ## VI Conclusions The present study shows that one needs to consider not only the statistical properties characterizing the time evolution of price for each stock traded but also the synchronous collective behavior of the portfolio considered to reveal the overall complexity of a financial market. We show that such a collective behavior of a portfolio of stock is efficiently monitored by the variety of the ensemble return distribution. This variable is directly observable for each portfolio and presents interesting statistical properties. It is non-Gaussian distributed and long-range correlated. The detailed statistical properties depends on the considered portfolio of stocks. We verify that for a portfolio of stocks characterized by comparable capitalization the distribution of the variety is approximately lognormal. Deviation from the lognormal behavior are observed for less homogeneous (in capitalization) portfolios. The shape of the distribution and the long-term memory of the variety are not reproduced by considering surrogated data simulated by using a single-index model with a realistic time series for the “market”. This implies that the complexity detected by the performed empirical analysis cannot be modeled with a similar simple stock price model. The correlations present in the market are more complex than the ones hypothesized by the single-index model. The correct modeling of the statistical properties of the variety can be then used as a benchmark for stock price models more sophisticated than the single-index model. The ensemble return distribution shows a qualitatively and quantitatively different behavior in “normal” and extreme trading days. The variety of a portfolio is then able to detect quite clearly shocks and aftershocks occurring in the market. Hence, it is a promising direct observable able to measure how much a portfolio is under pressure and how distant is from the typical market activity in a specific trading day. A theoretical challenge is to relate this empirical ensemble observation directly with the correlations active between pairs of stocks of a correlation. In summary, we believe that the overall complexity of a financial market can be detected and modeled only by considering simultaneously – (i) the statistical properties of the time evolution of stock prices of the considered portfolio and (ii) the statics and dynamics of the correlations existing between stocks. ## VII Acknowledgements The authors thank INFM and MURST for financial support. This work is part of the FRA-INFM project Volatility in financial markets. F. Lillo acknowledges FSE-INFM for his fellowships. We wish to thank Giovanni Bonanno for help in numerical calculations.
warning/0006/math0006157.html	ar5iv	text	# 1 Introduction ## 1 Introduction Since 1930, with the renowned papers by Douglas and Radó on minimal surfaces, the study of parametric two-dimensional surfaces with prescribed mean curvature, satisfying different kinds of geometrical or topological side conditions, has constituted a very challenging problem and has played a prominent role in the history of the Calculus of Variations. Surfaces with prescribed constant mean curvature are usually known as “soap films” or “soap bubbles”. This case has been successfully and deeply investigated by several authors, and nowadays a quite wide description of the problem is available in the literature (see the survey book by Struwe ). The phenomenon of the formation of an electrified drop is closely related to soap film and soap bubbles. As experimentally observed (see for example , , ), an external electric field may affect the shape of the drop, and its surface curvature turns out to be nonconstant, in general. However, as regards the mathematical treatment of the case of nonconstant prescribed mean curvature, only few existence results of variational type are known. Apart from few papers on the existence of a “small” solution for the Plateau problem (we quote, for instance, , and , see also ), all the other variational-type results hold true in a perturbative setting, namely, for curvatures of the form $`H(u)=H_0+H_1(u)`$ with $`H_0\{0\}`$ and $`H_1C^1(^3)L^{\mathrm{}}`$ having $`H_1_{\mathrm{}}`$ small. In particular, let us mention the papers , , , and , which deal with the Plateau problem, or the corresponding Dirichlet problem. In this paper we are interested in the existence of $`𝕊^2`$-type parametric surfaces in $`^3`$ having prescribed mean curvature $`H`$, briefly, $`H`$-bubbles. More precisely, for $`HC^1(^3)`$, an $`H`$-bubble is a nonconstant conformal function $`\omega :^2^3`$, smooth as a map on $`𝕊^2`$, satisfying the following problem: $$\{\begin{array}{cc}\mathrm{\Delta }\omega =2H(\omega )\omega _x\omega _y\hfill & \text{in }^2\hfill \\ _^2\|\omega \|^2<+\mathrm{}.\hfill & \end{array}$$ (1.1) Here $`\omega _x=(\frac{\omega _1}{x},\frac{\omega _2}{x},\frac{\omega _3}{x})`$, $`\omega _y=(\frac{\omega _1}{y},\frac{\omega _2}{y},\frac{\omega _3}{y})`$, $`\mathrm{\Delta }\omega =\omega _{xx}+\omega _{yy}`$, $`\omega =(\omega _x,\omega _y)`$, and $``$ denotes the exterior product in $`^3`$. In case of nonzero constant mean curvature $`H(u)H_0`$, Brezis and Coron proved that the only nonconstant solutions to (1.1) are spheres of radius $`\|H_0\|^1`$. In the present paper, we study the existence of $`H`$-bubbles with minimal energy in case $`H:^3`$ is a smooth function satisfying: $`sup_{u^3}\|H(u+\xi )uu\|<1`$, for some $`\xi ^3`$, $`H(u)H_{\mathrm{}}`$ as $`\|u\|\mathrm{}`$, for some $`H_{\mathrm{}}`$. The assumption $`(𝐡_\mathrm{𝟏})`$ is a global condition on the radial component of $`H(+\xi )`$ that, roughly speaking, measures how far $`H`$ differs from a constant. In addition, we also need that $`H`$ is nonzero on some sufficiently large set. This condition will be made clear in the following. In order to state our result we need some preliminaries. Let us point out that problem (1.1) has a natural variational structure, since solutions to (1.1) are formally the critical points of the functional $$_H(u)=\frac{1}{2}_^2\|u\|^2+2_^2Q(u)u_xu_y,$$ where $`Q:^3^3`$ is any vector field such that $`\mathrm{div}Q=H`$. Roughly speaking, the functional $`_^2Q(u)u_xu_y`$ has the meaning of a volume, for $`u`$ in a suitable space of functions. This is clear when $`H(u)H_0`$. Indeed in this case, taking $`Q(u)=\frac{H_0}{3}u`$, one deals with the standard volume functional $`_^2uu_xu_y`$ which is a determinant homogeneous in $`u`$ and, for $`u`$ constant far out, measures the algebraic volume enclosed by the surface parametrized by $`u`$. Moreover, it turns out to be bounded with respect to the Dirichlet integral by the Bononcini-Wente isoperimetric inequality. These facts hold true more generally when $`H`$ is a bounded nonzero function on $`^3`$ (see ). In particular, the functional $`_^2Q(u)u_xu_y`$ is essentially cubic in $`u`$ and it satisfies a generalized isoperimetric inequality. For this reason, we expect that $`_H`$ has a mountain pass structure, and this gives an indication for the existence of a nontrivial critical point. The natural space in order to look for $`𝕊^2`$-type solutions seems to be the Sobolev space $`H^1(𝕊^2,^3)`$, modulo stereographic projection. However, working with this space gives some technical difficulties due to the fact that $`H`$ may be nonconstant. In any case, we can define a mountain pass level for $`_H`$ restricted to some class of smooth functions. In addition, thanks to the assumption $`(𝐡_\mathrm{𝟏})`$, we can restrict ourselves to radial paths spanned by functions in $`𝒮_\xi =\{\xi +C_c^{\mathrm{}}(^2,^3):u\xi \}`$, where $`\xi ^3`$ is the same as in $`(𝐡_\mathrm{𝟏})`$. Thus we are lead to introduce the value $$c_H=\underset{u𝒮_\xi }{inf}\underset{s>0}{sup}_H(su).$$ The assumption $`(𝐡_{\mathrm{}})`$ guarantees that $$0<c_H\frac{4\pi }{3H_{\mathrm{}}^2}.$$ Note that if $`H_{\mathrm{}}0`$, the value $`\frac{4\pi }{3H_{\mathrm{}}^2}`$ equals the mountain pass level for the energy functional $`_H_{\mathrm{}}`$ corresponding to the constant mean curvature $`H_{\mathrm{}}`$. Moreover by the results proved by Brezis and Coron in , this value is the least critical value for $`_H_{\mathrm{}}`$ in $`H^1(𝕊^2,^3)`$, and it is attained by the spheres (with degree 1) of radius $`\|H_{\mathrm{}}\|^1`$. Now, our result can be stated as follows: ###### Theorem 1.1 Let $`HC^1(^3)`$ satisfy $`(𝐡_\mathrm{𝟏})`$ and $`(𝐡_{\mathrm{}})`$. If $$c_H<\frac{4\pi }{3H_{\mathrm{}}^2}$$ $`()`$ holds, then there exists an $`H`$-bubble $`\omega `$ such that $`_H(\omega )=c_H`$. Moreover, called $`_H`$ the set of $`H`$-bubbles, it holds that $`c_H=inf_{\omega _H}_H(\omega )`$. We point out that, thanks to $`(𝐡_\mathrm{𝟏})`$, the condition $`()`$ requires that $`_H(\overline{u})<0`$ for some $`\overline{u}𝒮_\xi `$ and then excluded the case $`H0`$. Clearly, when $`_H(\overline{u})<0`$ somewhere and $`H_{\mathrm{}}=0`$, then $`()`$ is automatically satisfied. Moreover, when $`H_{\mathrm{}}>0`$, the condition $`()`$ turns out to be true if $`H(u)>H_{\mathrm{}}`$ for $`\|u\|`$ large. Note that, in general, even if $`H(u)=H_{\mathrm{}}`$ for $`\|u\|R`$, Theorem 1.1 ensures that the $`H`$-bubble we find is different from the $`H_{\mathrm{}}`$-bubble located in the region $`\|u\|R`$. We also notice that in general we have no information about the position of the $`H`$-bubble given by Theorem 1.1. In particular, we can exhibit examples of radial curvatures $`H`$ for which $`H`$-bubbles with minimal energy exist but cannot be radial. The main difficulties in approaching problem (1.1) with variational methods concern the study of the Palais-Smale sequences. In particular, we emphasize the following problems: boundedness of a Palais-Smale sequence with respect to the Dirichlet norm, and in $`L^{\mathrm{}}`$; blow up analysis for a (bounded) Palais-Smale sequence. Concerning the first problem, the assumption $`(𝐡_\mathrm{𝟏})`$ can be useful in order to guarantee the boundedness with respect to the gradient $`L^2`$-norm. However the boundedness in $`L^{\mathrm{}}`$ in general cannot be deduced a priori and it is not just a technical difficulty. In fact, one can exhibit examples of Palais-Smale sequences which are bounded with respect to the Dirichlet norm, but not in $`L^{\mathrm{}}`$, and the lack of boundedness in $`L^{\mathrm{}}`$ cannot be eliminated in any way. Hence, because of these difficulties, we tackle the problem by using an approximation method in the spirit of a celebrated paper by Sacks and Uhlenbeck . More precisely, we construct a family of approximating solutions on which global and local estimates can be proved. In particular, assuming that $`H`$ is constant far out, we can obtain boundedness both with respect to the Dirichlet norm, and in $`L^{\mathrm{}}`$. Then, a limit procedure, involving a (partial) blow up analysis, is carried out, in order to show the existence of an $`H`$-bubble with minimal energy. In the last step, we remove the assumption that $`H`$ is constant far out, by an approximation argument on the curvature function, and we recover the full result stated in Theorem 1.1. We point out that for a curvature $`HC^1(^3)`$ satisfying $`(𝐡_\mathrm{𝟏})`$ and such that $`H(u)H_{\mathrm{}}0`$ for $`\|u\|`$ large, the set $`_H`$ of $`H`$-bubbles is nonempty a priori, and the existence of a minimal $`H`$-bubble can be obtained with a direct argument, just minimizing the energy functional $`_H`$ over $`_H`$, without using the above mentioned approximation method. In fact, the hard step lies in removing the condition that $`H`$ is constant far out, just asking to $`H`$ the asymptotic behaviour stated in $`(𝐡_{\mathrm{}})`$. To this goal, it is important to know that the energy of the minimal $`H`$-bubble is exactly $`c_H`$, and proving this needs either a sharp study of the behaviour of the Palais Smale sequences, or an (almost equivalent) approximation argument as, for instance, the Sacks-Uhlenbeck type argument that we develop. This step requires much more work and constitutes the largest part of this paper. We finally mention a result by Bethuel and Rey that states the existence of an $`H`$-bubble passing through an arbitrarily prescribed point in $`^3`$ in case $`H`$ is a perturbation of a nonzero constant. This result expresses the fact that the bubbles with constant curvature $`H_00`$ are stable with respect to small $`L^{\mathrm{}}`$ perturbations of $`H_0`$. Actually, in our opinion, the proof of this result is not completely clear and we are not able to recover it with our method. In fact, we think that the problem of existence of $`H`$-bubbles for a prescribed bounded curvature function $`H`$ has some similarities with a semilinear elliptic problem on $`^N`$ of the form $$\{\begin{array}{cc}\mathrm{\Delta }u+u=a(x)u^p\hfill & \text{on }^N\hfill \\ u>0\hfill & \text{on }^N\hfill \\ uH^1(^N)\hfill & \end{array}$$ (1.2) where $`1<p<\frac{N+2}{N2}`$ and $`a`$ is a bounded positive function on $`^N`$. It is known that the existence of solutions to (1.2) is strongly affected by the behaviour of the coefficient $`a(x)`$, and in some cases problem (1.2) has no solution. In particular, this may happen also when $`a(x)`$ is a small $`L^{\mathrm{}}`$ perturbation of a positive constant. In our opinion, similar considerations hold also for the problem of $`H`$-bubbles, and the behaviour of $`H(u)`$ plays a similar role of the coefficient $`a(x)`$ in (1.2). Hence, as well as for problem (1.2), we suspect that the existence of $`H`$-bubbles with minimal energy may depend in a very sensitive way on the function $`H`$. ## 2 The variational approach This Section is structured as follows. In the first part we introduce some notation in view of setting up a variational framework to study problem (1.1). In particular we define the $`H`$-volume functional, the energy functional associated to problem (1.1), and we recall some generalized isoperimetric inequality. In the second part we define a mountain pass level $`c_H`$ for the energy functional $`_H`$ and we discuss some properties related to the value $`c_H`$ strongly depending on the assumption $`(𝐡_\mathrm{𝟏})`$. ### 2.1 Notation and isoperimetric inequality First, let us introduce the space $$X=\{v\varphi :vH^1(𝕊^2,^3)\}$$ where $`\varphi :^2𝕊^2`$ is the (inverse of the) standard stereographic projection and it is given by $$\varphi (z)=(\mu x,\mu y,1\mu ),\mu =\mu (z)=\frac{2}{1+\|z\|^2},$$ (2.3) being $`z=(x,y)`$ and $`\|z\|^2=x^2+y^2`$. Notice that $`uX`$ if and only if $`u,\widehat{u}H_{loc}^1(^2,^3)`$ and $`_^2\|u\|^2<+\mathrm{}`$, where $`\widehat{u}(z)=u\left(\frac{z}{\|z\|^2}\right)`$. Let us also set $`H_0^1=H_0^1(D,^3)`$, where $`D`$ is the open unit disc in $`^2`$. Clearly, $`H_0^1X`$. For every $`uX`$ we denote the Dirichlet integral by $$𝒟(u)=\frac{1}{2}_^2\|u\|^2.$$ Now, given $`HC^1(^3)`$, we construct the $`H`$-volume functional as follows. Set $$m_H(u)=_0^1H(su)s^2𝑑s.$$ Thus, for every $`u^3`$ one has $$\mathrm{div}(m_H(u)u)=H(u).$$ (2.4) Then, let $`𝒱_H:XL^{\mathrm{}}`$ be defined by $$𝒱_H(u)=_^2m_H(u)uu_xu_y.$$ In case $`H(u)1`$, one has $`m_H(u)\frac{1}{3}`$, and the functional $`𝒱_H`$ reduces to the classical volume functional which satisfies the standard isoperimetric inequality. In fact the following generalization holds, as proved by Steffen in . ###### Lemma 2.1 If $`HC^1(^3)`$ is bounded on $`^3`$ then there exists $`S_H>0`$ such that $$S_H\|𝒱_H(u)\|^{2/3}𝒟(u)\mathrm{𝑓𝑜𝑟}\mathrm{𝑒𝑣𝑒𝑟𝑦}uXL^{\mathrm{}}.$$ (2.5) ###### Remark 2.2 In fact Steffen in proves that the functional $`𝒱_H`$ admits a continuous extension on $`H_0^1`$ and (2.5) holds true also for every $`uH_0^1`$. Finally we introduce the energy functional $`_H:XL^{\mathrm{}}`$, defined for every $`uXL^{\mathrm{}}`$ by $$_H(u)=𝒟(u)+2𝒱_H(u).$$ In the following result we state some properties of the functional $`_H`$. ###### Lemma 2.3 Let $`HC^1(^3)`$. Then: for every $`uXL^{\mathrm{}}`$ one has $`_H(su)=s^2𝒟(u)+o(s^2)`$ as $`s0`$, for every $`uXL^{\mathrm{}}`$ and for $`hC_c^{\mathrm{}}(^2,^3)`$ the directional derivative of $`_H`$ at $`u`$ along $`h`$ exists, and it is given by $$d_H(u)h=_^2uh+2_^2H(u)hu_xu_y,$$ for every bounded solution $`\omega `$ to (1.1) one has $$𝒟(\omega )+_^2H(\omega )\omega \omega _x\omega _y=0.$$ (2.6) ###### Remark 2.4 If $`\omega XL^{\mathrm{}}`$ is a weak solution to (1.1), i.e., $`d_H(\omega )h=0`$ for every $`hC_c^{\mathrm{}}(^2,^3)`$, then, since $`HC^1(^3)`$, by a Heinz regularity result , $`\omega C^3(^2,^3)`$, it is conformal, and smooth as a map on $`𝕊^2`$. In particular there exists $`lim_{\|z\|\mathrm{}}\omega (z)=\omega _{\mathrm{}}^3`$. Proof. Part $`(i)`$ is a consequence of Lemma 2.1. Part $`(ii)`$ follows by the results in , using (2.4). Finally, (2.6) can be proved multiplying the system $`\mathrm{\Delta }\omega =2H(\omega )\omega _x\omega _y`$ by $`\omega `$, integrating on $`D_R`$, and passing to the limit as $`R+\mathrm{}`$. To conclude this Subsection, we point out a consequence of assumption $`(𝐡_{\mathrm{}})`$. Actually, the following result holds true under a much weaker condition. ###### Lemma 2.5 Let $`HC^1(^3,)`$ satisfy $$\|H(su)\|H_0>0\mathrm{𝑓𝑜𝑟}ss_0\mathrm{𝑎𝑛𝑑}u\mathrm{\Sigma },$$ (2.7) being $`\mathrm{\Sigma }`$ a nonempty open set in $`𝕊^2`$. Then there exists $`\overline{u}H_0^1L^{\mathrm{}}`$ such that $`_H(s\overline{u})\mathrm{}`$ as $`s+\mathrm{}`$. Proof. Thanks to the rotational invariance of the problem we may assume that $`\mathrm{\Sigma }`$ is an open neighborhood of the point $`e_3=(0,0,1)`$. Furthermore, let us suppose that $`H(su)H_0>0`$ for $`s>s_0`$ and $`u\mathrm{\Sigma }`$. For $`\delta (0,1)`$ let us define $$u^\delta (z)=\{\begin{array}{cc}\varphi (z)\hfill & \text{as }\|z\|<\delta \hfill \\ \frac{1\|z\|}{1\delta }\varphi \left(\frac{\delta }{\|z\|}z\right)\hfill & \text{as }\delta \|z\|1,\hfill \end{array}$$ where $`\varphi :^2𝕊^2`$ is the function introduced in (2.3). It holds that $`u^\delta H_0^1L^{\mathrm{}}`$, and $`u^\delta `$ parametrizes the boundary of the sector of cone defined by $$A_\delta =\{\xi ^3:\|\xi \|\mathrm{cos}\theta _\delta >\xi e_3,\|\xi \|<1\}$$ where $`\theta _\delta =\mathrm{arccos}\frac{1\delta ^2}{1+\delta ^2}`$. In addition one has that $$u^\delta (z)u_x^\delta (z)u_y^\delta (z)=\{\begin{array}{cc}\mu (z)^2\hfill & \text{as }\|z\|<\delta \hfill \\ 0\hfill & \text{as }\|z\|>\delta \hfill \end{array}$$ and for every $`s>0`$, by the divergence theorem, $$𝒱_H(su^\delta )=_{sA_\delta }H(\xi )𝑑\xi .$$ Since $`\varphi (0)=e_3`$ and $`\varphi `$ is continuous, we can find $`\delta _0(0,1)`$ such that $`\varphi (z)\mathrm{\Sigma }`$ as $`\|z\|<\delta _0`$. Set $`\overline{u}=u^{\delta _0}`$ and $`A=A_{\delta _0}`$. Therefore, by the hypothesis, for $`s>s_0`$ one has $`𝒱_H(s\overline{u})`$ $`=`$ $`{\displaystyle _{s_0A}}H(\xi )𝑑\xi {\displaystyle _{sAs_0A}}H(\xi )𝑑\xi `$ $``$ $`𝒱_H(s_0\overline{u}){\displaystyle _{sAs_0A}}H_0𝑑\xi `$ $`=`$ $`𝒱_H(s_0\overline{u})𝒱_{H_0}(s_0\overline{u})s^3\|𝒱_{H_0}(\overline{u})\|.`$ Then $$_H(s\overline{u})s^2𝒟(\overline{u})+2(𝒱_H(s_0\overline{u})𝒱_{H_0}(s_0\overline{u}))2s^3\|𝒱_{H_0}(\overline{u})\|.$$ Passing to the limit as $`s+\mathrm{}`$ we obtain the thesis. Finally, we observe that in case $`H(su)H_0<0`$ for $`s>s_0`$ and $`u\mathrm{\Sigma }`$, one can repeat the same argument taking $`v(x,y)=u(y,x)`$. ### 2.2 The mountain pass level Assume that $`HC^1(^3)L^{\mathrm{}}`$ is such that there exists $`\overline{u}C_c^{\mathrm{}}(^2,^3)`$ with $`_H(\overline{u})<0`$. In particular, this excludes the case $`H0`$. Then, let $$c_H=\underset{\genfrac{}{}{0pt}{}{uC_c^{\mathrm{}}(^2,^3)}{u0}}{inf}\underset{s>0}{sup}_H(su).$$ (2.8) Note that $`c_H`$ is well defined and, thanks to Lemma 2.1, it is positive and finite. In particular, by (2.5), one can estimate $$c_H\left(\frac{S_H}{3}\right)^3,$$ where $`S_H`$ is the isoperimetric constant associated to $`H`$. ###### Remark 2.6 When $`H(u)H_0\{0\}`$, the volume functional is purely cubic and one can easily prove that $$c_{H_0}=\left(\frac{S_{H_0}}{3}\right)^3=\frac{4\pi }{3H_0^2}=_{H_0}(\omega ^0)=\underset{s>0}{sup}_{H_0}(s\omega ^0)$$ where $`\omega ^0=\frac{1}{H_0}\varphi `$ and $`\varphi `$ is defined in (2.3). Notice that $`\omega ^0`$ is a conformal parametrization of the sphere of radius $`\|H_0\|^1`$ centered at the origin, it satisfies $`\mathrm{\Delta }\omega ^0=2H_0\omega _x^0\omega _y^0`$ on $`^2`$, $`𝒟(\omega ^0)=\frac{4\pi }{H_0^2}`$, and $`𝒱_{H_0}(\omega ^0)=\frac{4\pi }{3H_0^2}`$. The results that follow better explain the role of the condition $`(𝐡_\mathrm{𝟏})`$ with respect to the definition of $`c_H`$. To this extent, we point out that, since problem (1.1) is invariant under translations, in the assumption $`(𝐡_\mathrm{𝟏})`$ we may suppose that $`\xi =0`$. Hence, setting: $$M_H=\underset{u^3}{sup}\|H(u)uu\|$$ (2.9) the hypothesis $`(𝐡_\mathrm{𝟏})`$ reads: $`M_H<1`$. It is convenient to introduce also the value $$\overline{M}_H=2\underset{u^3}{sup}\|(H(u)3m_H(u))u\|.$$ (2.10) In fact, several estimates in the sequel need a bound just on $`\overline{M}_H`$. ###### Remark 2.7 $`(i)`$ By (2.4) and by the definition of $`m_H`$, it turns out that $`\overline{M}_HM_H`$, but the strict inequality may also occur. Indeed one can construct functions $`HC^1(^3)`$ such that $`M_H=+\mathrm{}`$, while $`\overline{M}_H<+\mathrm{}`$. $`(ii)`$ If $`HC^1(^3)`$ satisfies $`\overline{M}_H<+\mathrm{}`$ then it turns out that $`HL^{\mathrm{}}(^3)`$. Furthermore, for every $`u𝕊^2`$ there exists $`lim_{s+\mathrm{}}H(su)=\widehat{H}(u)`$ and $`\widehat{H}C^0(𝕊^2)`$. Thus, if $`\overline{M}_H<+\mathrm{}`$, then the condition (2.7) used in Lemma 2.5 is verified whenever $`lim\; sup_{s+\mathrm{}}\|H(su)\|>0`$ for some $`u𝕊^2`$. First, we give a positive lower bound on the energy of any $`H`$-bubble. ###### Proposition 2.8 Let $`HC^1(^3)`$ satisfy $`(𝐡_\mathrm{𝟏})`$. If $`\omega `$ is an $`H`$-bubble, then $`_H(\omega )c_H`$. The proof of Proposition 2.8 is based on the following Lemma. ###### Lemma 2.9 Let $`HC^1(^3)`$ satisfy $`\overline{M}_H<1`$ and let $`uH_0^1L^{\mathrm{}}\{0\}`$. If $`sup_{s>0}_H(su)<+\mathrm{}`$ then $`_H(su)as^2bs^3`$ for every $`s>0`$, with $`a,b>0`$ depending on $`u`$, if $`_H(s_0u)<0`$ for some $`s_0>0`$ then $`sup_{s>0}_H(su)=\mathrm{max}_{s[0,s_0]}_H(su)`$, if $`sup_{s>0}_H(su)=_H(\overline{s}u)`$, then $`𝒱_H(\overline{s}u)<0`$. Proof. Fix $`uH_0^1L^{\mathrm{}}\{0\}`$ and set $`f(s)=_H(su)`$ for every $`s0`$. Notice that $`f`$ is differentiable and $$f^{}(s)=s_D\|u\|^2+2s^2_DH(su)uu_xu_y.$$ Using (2.10), one has that $$f^{}(s)(1\overline{M}_H)𝒟(u)s+\frac{3}{s}f(s).$$ (2.11) If $`sup_{s>0}f(s)<+\mathrm{}`$, since $`\overline{M}_H<1`$, (2.11) implies that $`lim_{s+\mathrm{}}f^{}(s)=\mathrm{}`$ and then there exists $`s_0>0`$ such that $`f(s)<0`$ for $`ss_0`$. Setting $`\overline{a}=(1\overline{M}_H)𝒟(u)`$ and integrating (2.11) over $`[s_0,s]`$ one obtains $$f(s)\left(\frac{f(s_0)}{s_0^3}\frac{\overline{a}}{s_0}\right)s^3+\overline{a}s^2$$ (2.12) for every $`ss_0`$. Keeping into account that $`f(s)=s^2𝒟(u)+o(s^2)`$ as $`s0^+`$, one can find $`a\overline{a}`$ such that $$f(s)\left(\frac{f(s_0)}{s_0^3}\frac{\overline{a}}{s_0}\right)s^3+as^2$$ for every $`s0`$, namely $`(i)`$. Now, let us prove $`(ii)`$. If $`f(s_0)<0`$, by (2.12), one infers that $`sup_{s>0}f(s)<+\mathrm{}`$. Moreover (2.11) implies in particular that $`f^{}(s)<0`$ whenever $`f(s)0`$. Hence also $`(ii)`$ holds true. Finally, if $`sup_{s>0}_H(su)=_H(\overline{s}u)`$, then $`f^{}(\overline{s})=0`$, and consequently, by (2.10), $$3𝒱_H(\overline{s}u)=3𝒱_H(\overline{s}u)\overline{s}f^{}(\overline{s})\overline{s}^2\left(1\frac{\overline{M}_H}{2}\right)𝒟(u)<0,$$ that is $`(iii)`$. Proof of Proposition 2.8. By Remark 2.4 an $`H`$-bubble $`\omega `$ is smooth and bounded. Moreover the mapping $`f(s)=_H(s\omega )`$ is well defined, and twice differentiable on $`(0,+\mathrm{})`$, with $$f^{\prime \prime }(s)=_D\|\omega \|^2+4s_DH(s\omega )\omega \omega _x\omega _y+2s^2_DH(s\omega )\omega \omega \omega _x\omega _y.$$ Since $`M_H<1`$, one obtains that $$f^{\prime \prime }(s)2(1M_H)𝒟(u)+\frac{2}{s}f^{}(s).$$ (2.13) In particular, by (2.13), if $`f^{}(\overline{s})=0`$ for some $`\overline{s}>0`$ then $`f^{\prime \prime }(\overline{s})<0`$. This shows that there exists at most one value $`\overline{s}>0`$ where $`f^{}(\overline{s})=0`$. In fact, one knows that $`f^{}(1)=0`$ because of (2.6). Hence $`sup_{s>0}f(s)=f(1)`$ and, arguing as in the proof of Lemma 2.9, $`f(s)\mathrm{}`$ as $`s+\mathrm{}`$. Now, for every $`\delta (0,1)`$ let $`u^\delta :D^3`$ be defined as follows: $$u^\delta (z)=\{\begin{array}{cc}0\hfill & \text{as }\|z\|\delta \hfill \\ \left(\frac{\mathrm{log}\|z\|}{\mathrm{log}\delta }1\right)\omega _{\mathrm{}}\hfill & \text{as }\delta ^2\|z\|<\delta \hfill \\ \left(\frac{\mathrm{log}\|z\|}{2\mathrm{log}\delta }1\right)(\omega ^\delta (z)\omega _{\mathrm{}})+\omega _{\mathrm{}}\hfill & \text{as }\delta ^4\|z\|<\delta ^2\hfill \\ \omega ^\delta (z)\hfill & \text{as }\|z\|<\delta ^4\hfill \end{array}$$ where $`\omega _{\mathrm{}}=lim_{\|z\|\mathrm{}}\omega (z)`$, and $`\omega ^\delta (z)=\omega (\frac{z}{\delta ^5})`$. Note that $`u^\delta H_0^1L^{\mathrm{}}`$ and $`u^\delta _{\mathrm{}}\omega _{\mathrm{}}`$. Let us set $`f_\delta (s)=_H(su^\delta )`$. We claim that for every $`s^{}>0`$ $$\underset{s[0,s^{}]}{sup}\|f_\delta (s)f(s)\|0\mathrm{as}\delta 0.$$ (2.14) Assuming for a moment that (2.14) holds, let us complete the proof. Let $`s_0>1`$ be such that $`f(s_0)<0`$. By (2.14), for $`\delta >0`$ small enough, $`f_\delta (s_0)<0`$ and then, by Lemma 2.9, $`sup_{s>0}f_\delta (s)`$ is attained in $`(0,s_0)`$. Hence, using again (2.14), we have $$c_H\underset{s>0}{sup}f_\delta (s)=\underset{s[0,s_0]}{\mathrm{max}}f_\delta (s)\underset{s[0,s_0]}{\mathrm{max}}f(s)+o(1)=f(1)+o(1).$$ Therefore the thesis follows. Finally, let us prove the claim (2.14). For every $`s0`$ we can write $`f_\delta (s)f(s)`$ $`=`$ $`s^2\left({\displaystyle _{\|z\|>\delta ^4}}\|u^\delta \|^2{\displaystyle _{\|z\|>\delta ^1}}\|\omega \|^2\right)`$ $`+`$ $`2s^3({\displaystyle _{\|z\|>\delta ^4}}m_H(su^\delta )u^\delta u_x^\delta u_y^\delta {\displaystyle _{\|z\|>\delta ^1}}m_H(s\omega )\omega \omega _x\omega _y).`$ We observe that $`2\left\|{\displaystyle _{\|z\|>\delta ^4}}m_H(su^\delta )u^\delta u_x^\delta u_y^\delta \right\|`$ $``$ $`m_H_{\mathrm{}}\omega _{\mathrm{}}{\displaystyle _{\|z\|>\delta ^4}}\|u^\delta \|^2`$ $`2\left\|{\displaystyle _{\|z\|>\delta ^1}}m_H(s\omega )\omega \omega _x\omega _y\right\|`$ $``$ $`m_H_{\mathrm{}}\omega _{\mathrm{}}{\displaystyle _{\|z\|>\delta ^1}}\|\omega \|^2.`$ Moreover, one can check that $$_{\|z\|>\delta ^4}\|u^\delta \|^20\mathrm{as}\delta 0,$$ and, since $`\omega X`$, also $$_{\|z\|>\delta ^1}\|\omega \|^20\mathrm{as}\delta 0.$$ Therefore (2.14) immediately follows and this concludes the proof. Notice that the full condition $`M_H<1`$ enters just in the previous step. Now we are going to prove two technical Lemmata that will be used in the sequel. ###### Lemma 2.10 Let $`HC^1(^3)`$ satisfy $`\overline{M}_H<1`$. Then $`c_Hc_{\lambda H}`$ for every $`\lambda (0,1]`$. Proof. Firstly, notice that for $`\lambda (0,1]`$, the isoperimetric inequality (2.5) holds true also for $`\lambda H`$ (with $`S_{\lambda H}=\lambda ^{\frac{2}{3}}S_H`$), and then the value $`c_{\lambda H}`$ is well defined and positive. Suppose that it is finite and, given $`ϵ>0`$, let $`uC_c^{\mathrm{}}(^2,^3)\{0\}`$ be such that $`sup_{s>0}_{\lambda H}(su)<c_{\lambda H}+ϵ`$. Since $`\overline{M}_{\lambda H}=\lambda \overline{M}_H<1`$, by Lemma 2.9, $`lim_{s+\mathrm{}}_{\lambda H}(su)=\mathrm{}`$. In particular, $`𝒱_{\lambda H}(su)<0`$ for $`s`$ large. Hence $`_H(su)_{\lambda H}(su)<0`$ for $`s`$ large. Using again Lemma 2.9 there exists $`\overline{s}>0`$ such that $`sup_{s>0}_H(su)=_H(\overline{s}u)`$. Furthermore $`𝒱_H(\overline{s}u)<0`$. Therefore $$c_H_H(\overline{s}u)=_{\lambda H}(\overline{s}u)+2(1\lambda )𝒱_H(\overline{s}u)_{\lambda H}(\overline{s}u)\underset{s>0}{sup}_{\lambda H}(su)c_{\lambda H}+ϵ.$$ Then the thesis follows because of the arbitrariness of $`ϵ>0`$. The next result states the upper semicontinuity of $`c_H`$ with respect to $`H`$. ###### Lemma 2.11 Let $`HC^1(^3)`$ satisfy $`\overline{M}_H<1`$. Let $`(H_n)C^1(^3)`$ be a sequence of functions satisfying $`\overline{M}_{H_n}<1`$, and such that $`H_nH`$ uniformly on compact sets of $`^3`$. Then $`lim\; sup_{n+\mathrm{}}c_{H_n}c_H`$. Proof. Suppose that $`c_H`$ is finite and, given $`ϵ>0`$ take $`uC_c^{\mathrm{}}(^2,^3)\{0\}`$ such that $`sup_{s>0}_H(su)<c_H+ϵ`$. One can check that $`lim_{n+\mathrm{}}_{H_n}(su)=_H(su)`$ for every $`s0`$. By Lemma 2.9, $`_H(s_0u)<0`$ for some $`s_0>0`$, and then also $`_{H_n}(s_0u)<0`$ for $`n`$ large enough. Therefore, since $`H_n`$ satisfies $`\overline{M}_{H_n}<1`$, using again Lemma 2.9, $`sup_{s>0}_{H_n}(su)=_{H_n}(\overline{s}_nu)`$ for some $`\overline{s}_n[0,s_0]`$. Then, for a subsequence, $`\overline{s}_n\overline{s}`$ and, since $`H_nH`$ uniformly on compact sets, $`_{H_n}(\overline{s}_nu)_H(\overline{s}u)`$. Consequently one has $$c_{H_n}_{H_n}(\overline{s}_nu)=_H(\overline{s}u)+o(1)\underset{s>0}{sup}_H(su)+o(1)c_H+ϵ+o(1).$$ Passing to the limit as $`n+\mathrm{}`$ and taking into account of the arbitrariness of $`ϵ>0`$, the thesis is proved. Lastly, we give an estimate for $`c_H`$ from above. Here, just the assumption $`(𝐡_{\mathrm{}})`$, and in fact a more general condition, is enough. ###### Lemma 2.12 Let $`HC^1(^3)`$ satisfy (2.7) for some nonempty open set $`\mathrm{\Sigma }𝕊^2`$. Then $`c_H\frac{4\pi }{3H_0^2}`$. Proof. As in the proof of Lemma 2.5, we may assume that $`\mathrm{\Sigma }`$ is an open neighborhood of the point $`e_3=(0,0,1)`$ and that $`H(su)H_0>0`$ for $`s>s_0`$ and $`u\mathrm{\Sigma }`$. Let us consider the function $`\omega ^0:^2^3`$ defined as in Remark 2.6. For every $`r>0`$ set $`\omega ^r=\omega ^0re_3`$. Notice that $`\omega ^r`$ is a conformal parametrization of a sphere of radius $`r_0=\frac{1}{H_0}`$ and center $`re_3`$. Hence, using the divergence theorem, one has that $$𝒱_H(s\omega ^r)=_{B_{sr_0}(se_3)}H(\xi )𝑑\xi =s^3_{B_{r_0}(0)}H(s\xi se_3)𝑑\xi .$$ (2.15) Setting $`s_r=\frac{s_0}{rr_0}`$ and using (2.15), one obtains that for $`s[0,s_r]`$ $$_H(s\omega ^r)4\pi (r_0s_r)^2+\frac{8\pi }{3}H_{\mathrm{}}(r_0s_r)^3=O(s_r^2),$$ while, for $`ss_r`$, by the hypothesis (2.7), one has $$_H(s\omega ^r)4\pi (r_0s)^2\frac{8\pi }{3}H_0(r_0s)^3\frac{4\pi }{3H_0^2}.$$ Then $$\underset{s>0}{sup}_H(s\omega ^r)\mathrm{max}\{\frac{4\pi }{3H_0^2},O(s_r^2)\}=\frac{4\pi }{3H_0^2}$$ (2.16) for $`r>0`$ large enough. Now, as in the proof of Proposition 2.8, one can construct $`u^{r,\delta }H_0^1L^{\mathrm{}}`$ such that $`sup_{s>0}_H(su^{r,\delta })sup_{s>0}_H(s\omega ^r)+o(1)`$, with $`o(1)0`$ as $`\delta 0`$. Hence, by (2.16), one obtains $`c_H\frac{4\pi }{3H_0^2}+o(1)`$, that is, the thesis. From the previous proof, one immediately infers the next estimate. ###### Corollary 2.13 Let $`HC^1(^3)`$ satisfy $`(𝐡_{\mathrm{}})`$. Then $`c_H\frac{4\pi }{3H_{\mathrm{}}^2}`$. If, in addition, $`H(u)>H_{\mathrm{}}>0`$ for $`\|u\|`$ large, then $`c_H<\frac{4\pi }{3H_{\mathrm{}}^2}`$. ## 3 Approximating problems Aim of this Section is to introduce a family of perturbed energy functionals having a mountain pass critical point at a level which approximate the value $`c_H`$ introduced in the previous Section. The advantage in following this procedure (already used in a different framework by Sacks and Uhlenbeck ) is due to the possibility to obtain some uniform global and local estimates on the critical points of the perturbed problems. Thus, for every $`\alpha >1`$ ($`\alpha `$ will be taken close to 1) we consider the Sobolev space $`H_0^{1,2\alpha }=H_0^{1,2\alpha }(D,^2)`$ and the functional $`_H^\alpha :H_0^{1,2\alpha }`$ defined by $$_H^\alpha (u)=\frac{1}{2\alpha }_D\left((1+\|u\|^2)^\alpha 1\right)+2𝒱_H(u).$$ It is convenient to denote $$𝒟^\alpha (u)=\frac{1}{2\alpha }_D\left((1+\|u\|^2)^\alpha 1\right).$$ Since $`H^{1,2\alpha }L^{\mathrm{}}H^1`$, the functional $`_H^\alpha `$ turns out to be well defined and regular on $`H_0^{1,2\alpha }`$, when $`H`$ is any bounded, smooth function. More precisely, $`_H^\alpha `$ is of class $`C^1`$ on $`H_0^{1,2\alpha }`$ and $$d_H^\alpha (u)h=_D(1+\|u\|^2)^{\alpha 1}uh+2_DH(u)hu_xu_y$$ for every $`u,hH_0^{1,2\alpha }`$ (see ). Our first goal is to prove that for every $`\alpha >1`$ sufficiently close to 1 the functional $`_H^\alpha `$ has a mountain pass geometry and a corresponding mountain pass critical point, as stated in the following result. ###### Lemma 3.1 Let $`HC^1(^3)L^{\mathrm{}}`$ be such that there exists $`\overline{u}C_c^{\mathrm{}}(D,^3)`$ with $`_H(\overline{u})<0`$. Then there exists $`\overline{\alpha }>1`$ such that for every $`\alpha (1,\overline{\alpha })`$ the class $`\mathrm{\Gamma }^\alpha =\{\gamma C([0,1],H_0^{1,2\alpha }):\gamma (0)=0,_H^\alpha (\gamma (1))<0\}`$ is nonempty and the value $$\overline{c}_H^\alpha =\underset{\gamma \mathrm{\Gamma }^\alpha }{inf}\underset{s[0,1]}{\mathrm{max}}_H^\alpha (\gamma (s))$$ is positive. If in addition $`\overline{M}_H<+\mathrm{}`$ then for every $`\alpha (1,\overline{\alpha })`$ there exists $`u^\alpha H_0^{1,2\alpha }`$ such that $`_H^\alpha (u^\alpha )=\overline{c}_H^\alpha `$ and $`d_H^\alpha (u^\alpha )=0`$. The second step consists in obtaining some uniform estimates on the mountain pass critical points $`u^\alpha `$ of the perturbed functionals $`_H^\alpha `$. ###### Proposition 3.2 Let $`HC^1(^3)`$ be such that $`\overline{M}_H<1`$ and, for every $`\alpha (1,\overline{\alpha })`$, let $`u^\alpha H^{1,2\alpha }`$ be the critical point of $`_H^\alpha `$ at level $`\overline{c}_H^\alpha `$ given by Lemma 3.1. Then $`\underset{\alpha 1}{lim\; sup}_H^\alpha (u^\alpha )c_H,`$ $`\underset{\alpha (1,\overline{\alpha })}{sup}u^\alpha _2<+\mathrm{},`$ $`\underset{\alpha (1,\overline{\alpha })}{inf}u^\alpha _2>0,`$ where $`c_H`$ is defined by (2.8). If, in addition, $`H(u)=H_0`$ for $`\|u\|R_0`$, for some $`R_0>0`$, then $$\underset{\alpha (1,\overline{\alpha })}{sup}u^\alpha _{\mathrm{}}<+\mathrm{}.$$ The proofs of Lemma 3.1 and Proposition 3.2 will be carried out in Subsections 3.1 and 3.2, respectively. The last result of this Section states the behaviour of the family of the mountain pass critical points $`u^\alpha `$ in the limit as $`\alpha 1`$. This result describes a blow up phenomenon, and it will be proved in the Appendix, in a more general situation. ###### Proposition 3.3 Let $`HC^1(^3)`$ be such that $`\overline{M}_H<1`$ and $`H(u)=H_0`$ for $`\|u\|R_0`$, for some $`R_0>0`$. For every $`\alpha (1,\overline{\alpha })`$, let $`u^\alpha H^{1,2\alpha }`$ be the critical point of $`_H^\alpha `$ at level $`\overline{c}_H^\alpha `$ given by Lemma 3.1. Then, there exist sequences $`(ϵ_\alpha )(0,+\mathrm{})`$, $`(z_\alpha )\overline{D}`$, a number $`\lambda (0,1]`$, and a function $`\omega XL^{\mathrm{}}`$ such that, setting $`v^\alpha (z)=u^\alpha (ϵ_\alpha z+z_\alpha )`$, for a subsequence, one has: $`ϵ_\alpha 0`$ and $`ϵ_\alpha ^{2(\alpha 1)}\lambda `$, $`v^\alpha \omega `$ strongly in $`H_{loc}^1(^2,^3)`$ and uniformly on compact sets of $`^2`$, $`\omega `$ is a nonconstant solution to $`\mathrm{\Delta }\omega =2\lambda H(\omega )\omega _x\omega _y`$ on $`^2`$, $`_{\lambda H}(\omega )\lambda lim\; inf_{\alpha 1}_H^\alpha (u^\alpha )`$. ### 3.1 Proof of Lemma 3.1 ###### Lemma 3.4 Let $`\rho (0,(\frac{S_H}{2})^{3/2}]`$ being $`S_H`$ given by (2.5). Then, for every $`uH_0^{1,2\alpha }`$ such that $`u_2\rho `$ one has $`_H^\alpha (u)\frac{1}{2}𝒟(u)`$. Proof. Using the inequality $`(1+s^2)^\alpha 1+\alpha s^2`$, one infers that $`_H^\alpha (u)_H(u)`$ for every $`uH_0^{1,2\alpha }`$. In addition, by the isoperimetric inequality (2.5) one has $`_H(u)𝒟(u)2S_H^{3/2}𝒟(u)^{3/2}`$. Therefore $`u_2\rho `$ implies $`_H^\alpha (u)(1\sqrt{2}S_H^{3/2}\rho )𝒟(u)`$ and the thesis follows since $`\rho (\frac{S_H}{2})^{3/2}`$. ###### Lemma 3.5 If $`uH_0^{1,2\overline{\alpha }}`$ for some $`\overline{\alpha }>1`$, then $`_H^\alpha (su)_H(su)`$ as $`\alpha 1`$, uniformly with respect to $`s[0,\overline{s}]`$ for every $`\overline{s}>0`$. Proof. The thesis follows by the estimate $`0`$ $``$ $`_H^\alpha (su)_H(su)={\displaystyle \frac{1}{2\alpha }}{\displaystyle _D}\left((1+s^2\|u\|^2)^\alpha 1\alpha s^2\|u\|^2\right)`$ $``$ $`{\displaystyle \frac{1}{2\alpha }}{\displaystyle _D}\left(2^{\alpha 1}1+2^{\alpha 1}s^{2\alpha }\|u\|^{2\alpha }\alpha s^2\|u\|^2\right),`$ and by standard techniques. ###### Lemma 3.6 If $`\overline{M}_H<+\mathrm{}`$ then for $`\alpha (1,\frac{3}{2})`$ the functional $`_H^\alpha `$ satisfies the Palais-Smale condition on $`H_0^{1,2\alpha }`$. Proof. First, note that for every $`uH_0^{1,2\alpha }`$, using (2.10) one has $`3_H^\alpha (u)d_H^\alpha (u)u`$ $``$ $`\left({\displaystyle \frac{3}{2\alpha }}1\right)𝒟^\alpha (u)+2{\displaystyle _D}(3m_H(u)H(u))uu_xu_y`$ $``$ $`\left({\displaystyle \frac{3}{2\alpha }}1\right)u_{2\alpha }^{2\alpha }{\displaystyle \frac{\overline{M}_H}{2}}u_2^2.`$ Hence, $$\left(\frac{3}{2\alpha }1\right)u_{2\alpha }^{2\alpha }\overline{M}_HC_\alpha u_{2\alpha }^2+d_H^\alpha (u)u_{2\alpha }+3_H^\alpha (u).$$ (3.17) Now, let $`(u^n)H_0^{1,2\alpha }`$ be a Palais-Smale sequence for $`_H^\alpha `$. By (3.17) the sequence $`(u^n)`$ is bounded in $`H_0^{1,2\alpha }`$. Then, there exists $`\overline{u}H_0^{1,2\alpha }`$ such that (for a subsequence) $`u^n\overline{u}`$ weakly in $`H^{1,2\alpha }`$ and uniformly on $`\overline{D}`$ (by Rellich Theorem). We need the following auxiliary result (see , for a proof): ###### Lemma 3.7 Let $`(u^n),(v^n)H_0^1L^{\mathrm{}}`$ be such that $`u^nu`$ weakly in $`H^1`$ and $`v^nv`$ uniformly. Then $$_Dv^nu_x^nu_y^n_Dvu_xu_y.$$ Since for every $`hH_0^{1,2\alpha }`$ $$_D(1+\|u^n\|^2)^{\alpha 1}u^nh+2_DH(u^n)hu_x^nu_y^n0$$ as $`n+\mathrm{}`$, thanks to Lemma 3.7 we obtain that $`d_H^\alpha (\overline{u})=0`$. In particular $`0=d_H^\alpha (\overline{u})(u^n\overline{u})=d𝒟^\alpha (u^n)(u^n\overline{u})+o(1)`$. On the other hand, we can use again Lemma 3.7 to get $`o(1)=d_H^\alpha (u^n)(u^n\overline{u})=d𝒟^\alpha (u^n)(u^n\overline{u})+o(1)`$. Therefore, $`(d𝒟^\alpha (u^n)d𝒟^\alpha (\overline{u}))(u^n\overline{u})=o(1)`$. Finally we note that $`𝒟^\alpha `$ is strictly convex on $`H_0^{1,2\alpha }`$, and hence $`d𝒟^\alpha `$ is strictly monotone. This readily leads to the conclusion. In conclusion, we notice that the first part of Lemma 3.1 is an immediate consequence of Lemmata 3.4 and 3.5. The existence of the critical point $`u^\alpha `$ is obtained as an application of the mountain pass theorem, and by Lemma 3.6. ### 3.2 Proof of Proposition 3.2 In order to show the first estimate, it is useful to introduce, for every $`\alpha (1,\overline{\alpha })`$, the value $$c_H^\alpha =\underset{\genfrac{}{}{0pt}{}{uH_0^{1,2\alpha }}{u0}}{inf}\underset{s>0}{sup}_H^\alpha (su).$$ ###### Lemma 3.8 Let $`HC^1(^3)`$ satisfy $`\overline{M}_H<1`$. Then $`lim\; sup_{\alpha 1}c_H^\alpha c_H`$. Proof. Fix $`ϵ>0`$ and take $`uC_c^{\mathrm{}}(D,^3)`$ such that $`sup_{s>0}_H(su)<c_H+ϵ`$. For every $`s0`$, using Lemma 2.9, one has $`_H^\alpha (su)`$ $`=`$ $`𝒟^\alpha (su)𝒟(su)+_H(su)`$ (3.18) $``$ $`C_0(s^{2\alpha }+1)C_1s^3`$ with $`C_0,C_1>0`$ depending just on $`u`$ (and not on $`\alpha `$). Therefore, for $`\alpha (1,\frac{3}{2})`$ there exists $`\overline{s}_\alpha >0`$ such that $`_H^\alpha (\overline{s}_\alpha u)=sup_{s>0}_H^\alpha (su)`$. From (3.18) it follows that $`\overline{s}_\alpha `$ is uniformly bounded. Then, by Lemma 3.5, $`lim_{\alpha 1}_H^\alpha (\overline{s}_\alpha u)=_H(\overline{s}u)`$ for some $`\overline{s}>0`$. Hence, $`lim\; sup_{\alpha 1}c_H^\alpha c_H+ϵ`$ and the thesis follows by the arbitrariness of $`ϵ>0`$. Concerning the $`H_0^1`$ bounds we have the following result. ###### Lemma 3.9 Let $`HC^1(^3)`$ satisfy $`\overline{M}_H<1`$. If $`uH_0^{1,2\alpha }`$ is a nonzero critical point of $`_H^\alpha `$, then $$\frac{1}{2}\left(\frac{2\overline{M}_H}{3}\right)^2S_H^3_D\|u\|^2\left(\frac{1}{2\alpha }\frac{1}{3}\frac{\overline{M}_H}{6}\right)^1_H^\alpha (u),$$ where $`S_H`$ is given by (2.5). Proof. Using (2.10) one has $`3_H^\alpha (u)`$ $`=`$ $`3𝒟^\alpha (u)d𝒟^\alpha (u)u+2{\displaystyle _D}(3m_H(u)H(u))uu_xu_y`$ $``$ $`\left({\displaystyle \frac{3}{2\alpha }}1\right){\displaystyle _D}\left(1+\|u\|^2\right)^{\alpha 1}\|u\|^2+{\displaystyle \frac{3}{2\alpha }}{\displaystyle _D}\left(\left(1+\|u\|^2\right)^{\alpha 1}1\right)`$ $`{\displaystyle \frac{\overline{M}_H}{2}}{\displaystyle _D}\|u\|^2`$ $``$ $`\left({\displaystyle \frac{3}{2\alpha }}1{\displaystyle \frac{\overline{M}_H}{2}}\right){\displaystyle _D}\|u\|^2.`$ Moreover, by (2.5) and (2.10) again, one has $`2𝒟(u)`$ $``$ $`{\displaystyle _D}\left(1+\|u\|^2\right)^{\alpha 1}\|u\|^2`$ $`=`$ $`6𝒱_H(u)+2{\displaystyle _D}\left(3m_H(u)H(u)\right)uu_xu_y`$ $``$ $`6S_H^{\frac{3}{2}}𝒟(u)^{\frac{3}{2}}+\overline{M}_H𝒟(u).`$ Since $`u0`$ one gets the thesis. Finally, to show the $`L^{\mathrm{}}`$ bound, $`H`$ is asked to be constant far out and the following estimate holds. ###### Lemma 3.10 Let $`HC^1(^3)`$ be such that $`H(u)=H_0`$ for $`\|u\|R_0`$, where $`R_0>0`$ is given. If $`uH_0^{1,2\alpha }`$ is a critical point of $`_H^\alpha `$, then $$u_{\mathrm{}}C\|H_0\|u_2^2+R_0$$ where $`C`$ is a universal positive constant (independent of $`\alpha ,R_0,H_0`$ and $`u`$). Proof. If $`uH_0^{1,2\alpha }`$ is a critical point of $`_H^\alpha `$, then $`u`$ is a weak solution to problem $$\{\begin{array}{cc}\mathrm{div}(a_\alpha (z)u)=2H(u)u_xu_y\hfill & \text{in }D\hfill \\ u=0\hfill & \text{on }D\hfill \end{array}$$ where $`a_\alpha (z)=(1+\|u(z)\|^2)^{\alpha 1}`$. Fix $`R>R_0`$ and let $`\mathrm{\Omega }_0`$ be a component of $`\{zD:\|u(z)\|>R\}`$, if there exists. Since $`u`$ is continuous, the set $`\mathrm{\Omega }_0`$ is nonempty, bounded, open and connected, and $`\|u\|=R`$ on $`\mathrm{\Omega }_0`$. Taking $`\delta (0,RR_0)`$ one can find a bounded, smooth domain $`\mathrm{\Omega }=\mathrm{\Omega }_\delta `$ close to $`\mathrm{\Omega }_0`$ such that $`\|u(z)\|>R_0`$ for $`z\mathrm{\Omega }`$ and $`\|u(z)\|R+\delta `$ for $`z\mathrm{\Omega }`$. Hence $`u`$ satisfies $$\mathrm{div}(a_\alpha (z)u)=2H_0u_xu_y\mathrm{on}\mathrm{\Omega }.$$ (3.19) For every $`k`$ let $`a_\alpha ^k=\mathrm{min}\{a_\alpha ,k\}`$ and let $`\phi ^k`$ be the solution to problem $$\{\begin{array}{cc}\mathrm{div}(a_\alpha ^k(z)\phi )=g\hfill & \text{in }\mathrm{\Omega }\hfill \\ \phi =0\hfill & \text{on }\mathrm{\Omega }\hfill \end{array}$$ (3.20) where $`g=2H_0u_xu_y`$. Since $`a_\alpha ^k`$ is a continuous bounded function on $`\mathrm{\Omega }`$ and $`a_\alpha ^k1`$, by a result of Bethuel and Ghidaglia, Theorem 1.3 in , there exists a constant $`C>0`$ such that $$\phi ^k_{\mathrm{}}+\phi ^k_2C\|H_0\|u_2^2$$ (3.21) and $`C`$ is independent of $`k,\alpha ,\mathrm{\Omega }`$ and $`u`$. Hence the sequence $`(\phi ^k)`$ is bounded in $`H_0^1(\mathrm{\Omega })`$ and thus, there exists $`\phi H_0^1(\mathrm{\Omega })`$ such that, for a subsequence, $`\phi ^k\phi `$ weakly in $`H_0^1(\mathrm{\Omega })`$ and pointwise a.e. We remark that $`a_\alpha L^{\frac{\alpha }{\alpha 1}}(\mathrm{\Omega })`$ since $`uH_0^{1,2\alpha }`$. In particular $`a_\alpha L^2(\mathrm{\Omega })`$ for $`\alpha <2`$ and $`a_\alpha ^ka_\alpha `$ strongly in $`L^2(\mathrm{\Omega })`$. By (3.20) for every $`hC_c^{\mathrm{}}(\mathrm{\Omega })`$ $$_\mathrm{\Omega }a_\alpha ^k(z)\phi ^kh=_\mathrm{\Omega }gh.$$ Hence, by a standard limit procedure, we obtain that for every $`hC_c^{\mathrm{}}(\mathrm{\Omega })`$ $$_\mathrm{\Omega }a_\alpha (z)\phi h=_\mathrm{\Omega }gh$$ that is, $`\phi `$ is a weak solution to $$\{\begin{array}{cc}\mathrm{div}(a_\alpha (z)\phi )=g\hfill & \text{in }\mathrm{\Omega }\hfill \\ \phi =0\hfill & \text{on }\mathrm{\Omega }\text{ .}\hfill \end{array}$$ (3.22) Moreover, by (3.21) we also get $$\phi _{\mathrm{}}+\phi _2C\|H_0\|u_2^2.$$ (3.23) Now, we observe that, thanks to (3.19) and (3.22), the function $`\psi =u\phi `$ is the solution to problem $$\{\begin{array}{cc}\mathrm{div}(a_\alpha (z)\psi )=0\hfill & \text{in }\mathrm{\Omega }\hfill \\ \psi =u\hfill & \text{on }\mathrm{\Omega }\hfill \end{array}$$ (3.24) and it can be characterized as the minimum for the problem $$inf\{_\mathrm{\Omega }a_\alpha (z)\|\psi \|^2:\psi u+H_0^1(\mathrm{\Omega })\}.$$ (3.25) Hence $`\psi _{\mathrm{}}R+\delta `$. Otherwise, if $`P`$ denotes the projection on the disc $`D_{R+\delta }`$, that is $$P(z)=\{\begin{array}{cc}z\hfill & \text{if }\|z\|R+\delta \hfill \\ (R+\delta )\frac{z}{\|z\|}\hfill & \text{if }\|z\|>R+\delta ,\hfill \end{array}$$ then $`\overline{\psi }=P\psi `$ will be a solution to (3.25) and then to (3.24). In conclusion, using (3.23), $$u_{\mathrm{}}\phi _{\mathrm{}}+\psi _{\mathrm{}}C\|H_0\|u_2^2+R+\delta ,$$ and by the arbitrariness of $`R>R_0`$ and $`\delta >0`$ one gets the thesis. Finally, Proposition 3.2 follows by Lemmata 3.8, 3.9 and 3.10, noting that $`c_H^\alpha \overline{c}_H^\alpha =_H^\alpha (u^\alpha )`$. ## 4 Proof of the main theorem Here we give the proof of Theorem 1.1. First, as a preliminary result we consider the case in which $`H`$ is constant outside a ball (Subsection 4.1). Then, in Subsection 4.2, we remove this condition, just asking $`H`$ to be asymptotic to a constant at infinity, according to the assumption $`(𝐡_{\mathrm{}})`$. ### 4.1 Case $`H`$ constant far out The results obtained in the previous Sections allow us to deduce the existence of an $`H`$-bubble when the prescribed curvature $`H`$ satisfies $`(𝐡_\mathrm{𝟏})`$ and is constant far out (this last condition enters in order to guarantee an $`L^{\mathrm{}}`$ bound on the approximating solutions). More precisely, the following result holds. ###### Theorem 4.1 Let $`HC^1(^3)`$ verify $`(𝐡_\mathrm{𝟏})`$ and the following conditions: there exists $`\overline{u}C_c^{\mathrm{}}(^2,^3)`$ such that $`_H(\overline{u})<0`$, there exists $`R_0>0`$ and $`H_0`$ such that $`H(u)=H_0`$ as $`\|u\|R_0`$. Then there exists an $`H`$-bubble $`\omega `$ such that $`_H(\omega )=c_H`$, where $`c_H`$ is defined by (2.8). ###### Remark 4.2 Suppose that in the assumption $`(ii)`$ $`H_00`$. Then, by Lemma 2.5, the condition $`(i)`$ is automatically fulfilled. Moreover, in this case problem (1.1) admits the (trivial) solution $`\omega ^0`$ which parametrizes a sphere of radius $`\|H_0\|^1`$ placed in the region $`\|u\|>R_0`$. However, the additional information on the energy of the $`H`$-bubble $`\omega `$ makes meaningful the above result, since if $`c_H<\frac{4\pi }{3H_0^2}`$ then $`\omega `$ is geometrically different from $`\omega ^0`$. Proof. From the assumptions $`(i)`$ and $`(ii)`$, and since $`\overline{M}_H<1`$, thanks to Propositions 3.2 and 3.3, there exists a function $`\omega XL^{\mathrm{}}`$ which is a $`\lambda H`$-bubble with $`\lambda (0,1]`$ and $`_{\lambda H}(\omega )\lambda c_H`$. Since $`M_H<1`$, by Proposition 2.8 (applied with $`\lambda H`$ instead of $`H`$), $`_{\lambda H}(\omega )c_{\lambda H}`$. Finally, Lemma 2.10 implies $`_{\lambda H}(\omega )c_H`$. Then $`\lambda =1`$ and the Theorem is proved. ### 4.2 General case Now we want to remove the hypothesis that $`H`$ is constant far out, by requiring just an asymptotic behaviour at infinity as stated by $`(𝐡_{\mathrm{}})`$. To this aim, we will use the condition $`()`$. Our argument consists in approximating $`H`$ with a sequence of functions $`(H_n)C^1(^3)`$ satisfying the hypotheses of Theorem 4.1 and then passing to the limit on the sequence $`(\omega ^n)`$ of the corresponding $`H_n`$-bubbles. The information on the energies $`_{H_n}(\omega ^n)`$ together with the condition $`()`$ will permit us to obtain some $`L^{\mathrm{}}`$ bound on the sequence $`(\omega ^n)`$, and then to get the result. Thus, let us start with the construction of the sequence $`(H_n)`$. ###### Lemma 4.3 Let $`HC^1(^3)`$ satisfying $`(𝐡_{\mathrm{}})`$ and let $`M_H`$ be defined by (2.9). Then there exists a sequence $`(H_n)C^1(^3)`$ such that: $`H_nH`$ uniformly on $`^3`$, for every $`n`$ there exists $`R_n>0`$ such that $`H_n(u)=H_{\mathrm{}}`$ as $`\|u\|R_n`$, $`sup_{u^3}\|H_n(u)uu\|:=M_{H_n}M_H`$. Proof. It is not restrictive to suppose $`H_{\mathrm{}}=0`$. Hence, for every $`u^3\{0\}`$ one has $$H(u)=_1^+\mathrm{}H(su)u𝑑s.$$ Let $`\chi C^{\mathrm{}}(,[0,1])`$ be such that $`\chi (r)=1`$ as $`r0`$, $`\chi (r)=0`$ as $`r1`$ and $`\|\chi ^{}\|2`$. Given any sequence $`r_n+\mathrm{}`$ set $`\chi _n(r)=\chi (rr_n)`$ and $$H_n(u)=_1^+\mathrm{}\chi _n(s\|u\|)H(su)u𝑑s$$ for every $`u^3\{0\}`$. By continuity, $`H_n`$ is well defined and continuous on $`^3`$. In fact $`H_nC^1(^3)`$ and for each $`u^3\{0\}`$ $$H_n(u)u=\frac{d}{ds}H_n(su)\|_{s=1}=\chi _n(\|u\|)H(u)u.$$ (4.26) Therefore $`(iii)`$ holds true. By the definition of $`H_n`$, one has $`H_n(u)=0`$ as $`\|u\|>r_n+1`$. Thus $`(ii)`$ follows, with $`R_n=r_n+1`$. Moreover (4.26) implies $`(iii)`$. Now, notice that $$H_n(u)=\chi _n(\|u\|)H(u)+_{r_n}^{r_n+1}\chi _n^{}(t)H\left(t\frac{u}{\|u\|}\right)𝑑t.$$ (4.27) Setting $`ϵ_n=sup_{\|u\|r_n}\|H(u)\|`$, one has that $`\left\|{\displaystyle _{r_n}^{r_n+1}}\chi _n^{}(t)H\left(t{\displaystyle \frac{u}{\|u\|}}\right)𝑑t\right\|`$ $``$ $`2ϵ_n`$ $`\left\|(\chi _n(\|u\|)1)H(u)\right\|`$ $``$ $`2ϵ_n.`$ Hence, (4.27) implies that $`\|H_n(u)H(u)\|4ϵ_n`$ for every $`u^3`$ and then, since $`ϵ_n0`$, also $`(i)`$ is proved. As a further tool, we also need the following result. ###### Lemma 4.4 Let $`(\stackrel{~}{H}_n)C^1(^3)`$, $`H_{\mathrm{}}`$ and $`(\stackrel{~}{\omega }_n)XL^{\mathrm{}}`$ be such that: $`\stackrel{~}{H}_nH_{\mathrm{}}`$ uniformly on compact sets, $`sup_n\left(\stackrel{~}{\omega }^n_2+\stackrel{~}{\omega }^n_{\mathrm{}}\right)<+\mathrm{}`$, for every $`n`$ the function $`\stackrel{~}{\omega }^n`$ solves $`\mathrm{\Delta }\stackrel{~}{\omega }^n=2\stackrel{~}{H}_n(\stackrel{~}{\omega }^n)\stackrel{~}{\omega }_x^n\stackrel{~}{\omega }_y^n`$ on $`^2`$. Then $`H_{\mathrm{}}0`$ and $`lim\; inf_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n)\frac{4\pi }{3H_{\mathrm{}}^2}`$. Proof. From the assumption $`(ii)`$, there exists $`\omega XL^{\mathrm{}}`$ such that, for a subsequence, $`\omega ^n\omega `$ weakly in $`(L^2(^2,^3))^2`$. Thanks to the invariance of $`H`$-systems with respect to dilations, translations and Kelvin transform, we may also assume that $`\stackrel{~}{\omega }^n_{\mathrm{}}=\|\stackrel{~}{\omega }^n(0)\|=1`$. Then, arguing as in the proof of Proposition A.1, using the hypotheses $`(i)`$$`(iii)`$, one can show that $`\omega `$ is an $`H_{\mathrm{}}`$-bubble, and $`\stackrel{~}{\omega }^n\omega `$ strongly in $`H_{loc}^1(^2,^3)`$ and in $`L_{loc}^{\mathrm{}}(^2,^3)`$. In particular it must be $`H_{\mathrm{}}0`$ (there exists no $`0`$-bubble in $`X`$). Furthermore, for every $`R>0`$, one has $`_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n,D_R)`$ $``$ $`_H_{\mathrm{}}(\omega ,D_R)`$ (4.28) $`{\displaystyle _{D_R}}\stackrel{~}{\omega }^n{\displaystyle \frac{\stackrel{~}{\omega }^n}{\nu }}`$ $``$ $`{\displaystyle _{D_R}}\omega {\displaystyle \frac{\omega }{\nu }},`$ (4.29) where, in (4.28), we used the notation: $$_H(u,\mathrm{\Omega })=\frac{1}{2}_\mathrm{\Omega }\|u\|^2+2_\mathrm{\Omega }m_H(u)uu_xu_y.$$ Now, fixing $`ϵ>0`$, let $`R>0`$ be such that $$\|_H_{\mathrm{}}(\omega ,^2D_R)\|<ϵ,\left(\|H_{\mathrm{}}\|\omega _{\mathrm{}}+1\right)_{^2D_R}\|\omega \|^2<ϵ.$$ Multiplying $`\mathrm{\Delta }\omega =2H_{\mathrm{}}\omega _x\omega _y`$ by $`\omega `$ and integrating over $`^2D_R`$ we find $$\left\|_{D_R}\omega \frac{\omega }{\nu }\right\|=\left\|_{^2D_R}\left(\omega \mathrm{\Delta }\omega +\|\omega \|^2\right)\right\|ϵ.$$ Then, by (4.29), one has that $$\left\|_{D_R}\stackrel{~}{\omega }^n\frac{\stackrel{~}{\omega }^n}{\nu }\right\|ϵ+o(1).$$ (4.30) Now we multiply $`\mathrm{\Delta }\stackrel{~}{\omega }^n=2\stackrel{~}{H}_n(\stackrel{~}{\omega }^n)\stackrel{~}{\omega }_x^n\stackrel{~}{\omega }_y^n`$ by $`\stackrel{~}{\omega }^n`$ and we integrate over $`^2D_R`$ to get $`{\displaystyle _{D_R}}\stackrel{~}{\omega }^n{\displaystyle \frac{\stackrel{~}{\omega }^n}{\nu }}`$ $`=`$ $`{\displaystyle _{^2D_R}}\|\stackrel{~}{\omega }^n\|^2+2{\displaystyle _{^2D_R}}\stackrel{~}{H}_n(\stackrel{~}{\omega }^n)\stackrel{~}{\omega }^n\stackrel{~}{\omega }_x^n\stackrel{~}{\omega }_y^n`$ (4.31) $`=`$ $`3_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n,^2D_R){\displaystyle \frac{1}{2}}{\displaystyle _{^2D_R}}\|\stackrel{~}{\omega }^n\|^2`$ $`+2{\displaystyle _{^2D_R}}(\stackrel{~}{H}_n(\stackrel{~}{\omega }^n)3m_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n))\stackrel{~}{\omega }^n\stackrel{~}{\omega }_x^n\stackrel{~}{\omega }_y^n`$ $``$ $`3_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n,^2D_R)\left({\displaystyle \frac{1}{2}}\mu _n\rho \right){\displaystyle _{^2D_R}}\|\stackrel{~}{\omega }^n\|^2`$ where $`\rho =sup_n\stackrel{~}{\omega }^n_{\mathrm{}}`$, and $`\mu _n=sup_{\|u\|\rho }\|\stackrel{~}{H}_n(u)3m_{\stackrel{~}{H}_n}(u)\|`$. Hence (4.30) and (4.31) imply $$_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n,^2D_R)\frac{ϵ}{3}+o(1),$$ because, by $`(i)`$, $`\mu _n0`$. Finally, we have $`_H_{\mathrm{}}(\omega )ϵ`$ $``$ $`_H_{\mathrm{}}(\omega ,D_R)=_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n,D_R)+o(1)`$ $`=`$ $`_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n)_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n,^2D_R)+o(1)_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n)+{\displaystyle \frac{ϵ}{3}}+o(1).`$ Hence, by the arbitrariness of $`ϵ>0`$, one obtains $`lim\; inf_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n)_H_{\mathrm{}}(\omega )`$ and the thesis follows by Remark 2.6. Proof of Theorem 1.1. Let $`(H_n)C^1(^3)`$ be the sequence given by Lemma 4.3. From Theorem 4.1, for every $`n`$ there exists an $`H_n`$-bubble $`\omega ^n`$ such that $`_{H_n}(\omega ^n)=c_{H_n}`$. By Lemma 2.11, one has that $$\underset{n+\mathrm{}}{lim\; sup}_{H_n}(\omega ^n)c_H.$$ (4.32) We point out that if we prove that $`sup_n\omega ^n_{\mathrm{}}=R<+\mathrm{}`$, then we have concluded, since for $`n`$ large, $`H(u)=H_n(u)`$ as $`\|u\|R`$. To this goal, as a first step, we show that $$\omega ^n\omega _{\mathrm{}}^n_{\mathrm{}}C_1\left(1+_^2\|\omega ^n\|^2\right)$$ (4.33) where $`\omega _{\mathrm{}}^n=lim_{\|z\|\mathrm{}}\omega ^n(z)`$ and $`C_1>0`$ depends only on $`H_{\mathrm{}}`$. This is a consequence of an a priori $`L^{\mathrm{}}`$ estimate proved by Grüter (see also Theorem 4.8 in ). More precisely, fixing an arbitrary $`\delta >0`$, for every $`n`$ there exists $`\rho _n>0`$, depending on $`\delta `$, such that if $`\|z\|\rho _n`$ then $`\|\omega ^n(z)\omega _{\mathrm{}}^n\|\delta `$. Let us set $`\gamma ^n(z)`$ $`=`$ $`\omega ^n(\rho _nz)\omega _{\mathrm{}}^n\mathrm{as}zD`$ $`u^n(z)`$ $`=`$ $`\omega ^n(\rho _nz)\omega _{\mathrm{}}^n\mathrm{as}zD.`$ Thus $`u^n`$ is a smooth and conformal solution to $$\{\begin{array}{cc}\mathrm{\Delta }u^n=2\stackrel{~}{H}_n(u^n)u_x^nu_y^n\hfill & \text{in }D\hfill \\ u^n=\gamma ^n\hfill & \text{on }D,\hfill \end{array}$$ where $`\stackrel{~}{H}_n(u)=H_n(u+\omega _{\mathrm{}}^n)`$. Hence, by , $$u^n_{L^{\mathrm{}}(D)}\gamma ^n_{L^{\mathrm{}}(D)}+C\left(1+_D\|\omega ^n\|^2\right)$$ with $`C>0`$ depending on $`\stackrel{~}{H}_n_{\mathrm{}}=H_n_{\mathrm{}}`$. Since $`H_nH`$ uniformly on $`^3`$, actually, $`C`$ is independent of $`n`$, but depends only on $`H_{\mathrm{}}`$. Then $$\omega ^n\omega _{\mathrm{}}^n_{L^{\mathrm{}}(^2)}\delta +u^n_{L^{\mathrm{}}(D)}2\delta +C\left(1+_^2\|\omega ^n\|^2\right).$$ Therefore (4.33) holds true. As a second step, we show that for every $`n`$ $$_^2\|\omega ^n\|^2C_2$$ (4.34) where $`C_2>0`$ depends only on $`H`$. Indeed, by (2.9), using (2.6), one has $`(1M_{H_n})𝒟(\omega ^n)3_{H_n}(\omega ^n)`$. Since $`M_{H_n}M_H`$, from (4.32) it follows that $$\underset{n+\mathrm{}}{lim\; sup}_^2\|\omega ^n\|^2\frac{6c_H}{1M_H},$$ and thus (4.34) is proved. Consequently, by (4.33), one obtains $$\omega ^n\omega _{\mathrm{}}^n_{\mathrm{}}C_3$$ (4.35) with $`C_3>0`$ independent of $`n`$. As a last step, let us show that $`sup_n\|\omega _{\mathrm{}}^n\|<+\mathrm{}`$. We argue by contradiction, assuming that (for a subsequence) $`\|\omega _{\mathrm{}}^n\|+\mathrm{}`$. Since $`H_nH`$ uniformly on $`^3`$, by $`(𝐡_{\mathrm{}})`$, we have that $`\stackrel{~}{H}_nH_{\mathrm{}}`$ uniformly on compact sets. Moreover, $`\stackrel{~}{\omega }^n(z)=\omega ^n(z)\omega _{\mathrm{}}^n`$ is an $`\stackrel{~}{H}_n`$-bubble and, thanks to (4.35) and (4.34), we can apply Lemma 4.4, to infer that $`H_{\mathrm{}}0`$ and $$\frac{4\pi }{3H_{\mathrm{}}^2}lim\; inf_{\stackrel{~}{H}_n}(\stackrel{~}{\omega }^n)=lim\; inf_{H_n}(\omega ^n).$$ Then (4.32) implies that $`\frac{4\pi }{3H_{\mathrm{}}^2}c_H`$, contrary to the condition $`()`$. Therefore, we have that $`sup\|\omega _{\mathrm{}}^n\|<+\mathrm{}`$, that, together with (4.35), gives the desired estimate. This concludes the proof. We end the work, by making some comments about the case of radially symmetric curvatures. ###### Example 4.5 Let $`HC^1(^3)`$ be a radial function satisfying $`(𝐡_\mathrm{𝟏})`$ and $`(𝐡_{\mathrm{}})`$ with $`H_{\mathrm{}}0`$. Given $`\varphi :^2𝕊^2`$ defined by (2.3), and $`\rho >0`$, the mapping $`\rho \varphi `$ is a solution to (1.1), i.e., it is a radial $`H`$-bubble, if and only if $`\rho \|H(\rho )\|=1`$. In this case the energy of this radial $`H`$-bubble is $`\frac{4\pi }{3H(\rho )^2}`$. Clearly, since $`H`$ is regular and $`H_{\mathrm{}}0`$, the equation $`\rho \|H(\rho )\|=1`$ always admits positive solutions. Now, suppose, in addition that the condition $`()`$ holds true. This happens, for instance, if $`H(\rho )>H_{\mathrm{}}>0`$ for $`\rho `$ large. Then, there exist $`H`$-bubbles with minimal energy $`c_H<\frac{4\pi }{3H_{\mathrm{}}^2}`$. Hence, these minimal $`H`$-bubbles cannot be radial if $`\|H(\rho )\|H_{\mathrm{}}`$ whenever $`\rho \|H(\rho )\|=1`$. Acknowledgments Work supported by M.U.R.S.T. progetto di ricerca “Metodi Variazionali ed Equazioni Differenziali Nonlineari” (cofin. 2001/02) ## Appendix A Appendix ### A.1 Convergence of approximating solutions in a <br>Sacks-Uhlenbeck type setting In this Appendix we study the behaviour of sequences of solutions of approximating problems of the type $$\{\begin{array}{cc}\mathrm{div}((1+\|u\|^2)^{\alpha 1}u)=2H(u)u_xu_y\hfill & \text{in }D\hfill \\ u=0\hfill & \text{on }D\hfill \end{array}$$ in the limit as $`\alpha 1_+`$. More precisely, we assume that for every $`\alpha (1,\overline{\alpha })`$ a function $`u^\alpha H_0^{1,2\alpha }`$ is given, in such a way that $`d_H^\alpha (u^\alpha )=0,`$ (A.36) $`\underset{\alpha (1,\overline{\alpha })}{sup}\left(u^\alpha _{\mathrm{}}+u^\alpha _2\right)<+\mathrm{},`$ (A.37) $`\underset{\alpha (1,\overline{\alpha })}{inf}u^\alpha _2>0.`$ (A.38) The first main result is non-variational and concerns a blow up analysis of sequences of approximating solutions. We point out that this result applies to any sequence of functions satisfying (A.36)–(A.38). ###### Proposition A.1 Let $`HC^1(^3)L^{\mathrm{}}`$ and for every $`\alpha (1,\overline{\alpha })`$ let $`u^\alpha H_0^{1,2\alpha }`$ satisfy (A.36)–(A.38). Then, there exist sequences $`(ϵ_\alpha )(0,+\mathrm{})`$, $`(z_\alpha )\overline{D}`$, a number $`\lambda (0,1]`$, and a function $`\omega XL^{\mathrm{}}`$ such that, setting $`v^\alpha (z)=u^\alpha (ϵ_\alpha z+z_\alpha )`$, for a subsequence, one has: $`ϵ_\alpha 0`$ and $`ϵ_\alpha ^{2(\alpha 1)}\lambda `$, $`v^\alpha \omega `$ strongly in $`H_{loc}^1(^2,^3)`$ and uniformly on compact sets of $`^2`$, $`\omega `$ is a nonconstant solution to $`\mathrm{\Delta }\omega =2\lambda H(\omega )\omega _x\omega _y`$ on $`^2`$. Notice that, according to Proposition A.1, in the limiting problem the curvature function is $`\lambda H`$, with $`\lambda (0,1]`$, and not necessarily $`\lambda =1`$. The second important result of this Appendix is variational and states a semicontinuity property, under an additional assumption on $`H`$, involving the value $`\overline{M}_H`$ defined by (2.10). ###### Proposition A.2 Let $`HC^1(^3)`$ be such that $`\overline{M}_H<1`$. For $`\alpha (1,\overline{\alpha })`$ let $`u^\alpha H_0^{1,2\alpha }`$ satisfy (A.36)–(A.38), and let $`\lambda (0,1]`$ and $`\omega XL^{\mathrm{}}`$ be given by Proposition A.1. Then $$_{\lambda H}(\omega )\lambda \underset{\alpha 1}{lim\; inf}_H^\alpha (u^\alpha ).$$ To prove Proposition A.1, first of all we need some local estimates on the family $`(u^\alpha )`$. This will be developed in Subsection A.1. Then the proof of Proposition A.1 will be performed in Subsection A.2. Finally, Proposition A.2 will be proved in Subsection A.3. ### A.2 Local estimates ($`\epsilon `$-regularity) Here we study the regularity properties of critical points for $`_H^\alpha `$, following the arguments by Sacks and Uhlenbeck . The first (minor) difference with respect to the framework of Sacks and Uhlenbeck paper lies in the nonlinear term. In the Euler-Lagrange equation for the harmonic map problem involves the second fundamental form of the embedding of the target space $`N`$ into an Euclidean space, instead of the curvature term. This is far to lead to any extra difficulty, since the invariance of the curvature term with respect to dilations makes computations even easier, in this case. The main difference with concerns the $`L^{\mathrm{}}`$ bound on the maps $`u`$ under consideration. In their paper, Sacks and Uhlenbeck deal with maps $`u`$ whose target space is a compact Riemannian manifold, and therefore, they have a natural $`L^{\mathrm{}}`$ bound on all maps $`u`$. On the contrary, the target space of our maps $`u`$ is the noncompact space $`^3`$, and hence we have no natural a priori bound. Therefore we have to ask it as an hypothesis. Another difference with respect to the proof of Sacks and Uhlenbeck is due to the presence of a boundary in the domain. However, this does not lead any extra difficulty. One can argue, for example, as in Struwe , Proposition 2.6. The first result concerns global regularity for fixed $`\alpha >1`$, and it can be obtained as in , using Theorem 1.11.$`1^{}`$ in and Struwe , proof of Proposition 2.6, for the regularity up to the boundary. ###### Lemma A.3 Let $`HC^1(^3)`$ and let $`uH_0^{1,2\alpha }`$ be a critical point of $`_H^\alpha `$ for some $`\alpha >1`$. Then $`u`$ belongs to $`W^{2,q}(D,^3)`$ for every $`q[1,+\mathrm{})`$ and solves $$\mathrm{\Delta }u=\frac{2(\alpha 1)}{1+\|u\|^2}(^2u,u)u+\frac{2H(u)}{(1+\|u\|^2)^{\alpha 1}}u_xu_y\mathrm{𝑖𝑛}D.$$ (A.39) The second result of this Section concerns some local estimates for the solutions of the approximating problems ($`\epsilon `$-regularity) which are actually the same as in the celebrated paper , and which are stated in the following Lemma (compare also with Lemma A.1 in ). We restrict ourselves to make estimates in the interior of the disk, thanks to the extension argument by Struwe . ###### Lemma A.4 (Main Estimate) Let $`HC^1(^3)L^{\mathrm{}}`$. Then there exist $`\overline{\epsilon }=\overline{\epsilon }(H_{\mathrm{}})>0`$, and for every $`p(1,+\mathrm{})`$ an exponent $`\alpha _p>1`$ and a constant $`C_p=C_p(H_{\mathrm{}})>0`$, such that if $`\alpha [1,\alpha _p)`$ and $`uW_{loc}^{2,p}(D,^3)`$ solves (A.39), then $$u_{L^2(D_R(z))}\overline{\epsilon }u_{H^{1,p}(D_{R/2}(z))}C_pR^{\frac{2}{p}2}u_{L^2(D_R(z))}$$ for every disc $`\overline{D_R(z)}D`$. Proof. Our arguments strictly follow the original proof in . Let $`u`$ be a solution to (A.39) for some $`\alpha 1`$. Fixing $`zD`$, for $`R(0,1\|z\|)`$ we expand $`D_R(z)`$ to the unit disc $`D`$, and we define a map $`\omega :D^3`$ by setting $$\omega (\zeta )=u(R\zeta +z)_{D_R(z)}u$$ A direct computation shows that $`\omega `$ is a regular solution in $`D`$ to the system $$\mathrm{\Delta }\omega =\frac{2(\alpha 1)}{R^2+\|\omega \|^2}(^2\omega ,\omega )\omega +2\frac{H_R(\omega )}{(R^2+\|\omega \|^2)^{\alpha 1}}\omega _x\omega _y$$ (A.40) where $`H_R(\omega )=R^{2(\alpha 1)}H(\omega +_{D_R(z)}u)`$. Note also that $`R^{2(\alpha 1)}H_R_{\mathrm{}}\overline{H}=:H_{\mathrm{}}`$. Now fix four radii $`\frac{1}{2}=r_0<r_1<r_2<r_3=1`$ and three cut-off functions $`\phi _iC^{\mathrm{}}(^2,[0,1])`$ such that $`\phi _i1`$ on $`D_{r_{i1}}`$, $`\phi _i0`$ on $`^2D_{r_i}`$ ($`i=1,2,3`$). Let $`K=\mathrm{max}_i\left(\phi _i_{\mathrm{}}+^2\phi _i_{\mathrm{}}\right)`$. Our aim is to use equation (A.40) in order to obtain some estimate on $`\phi _i\omega `$. First, we point out some simple inequalities: $$\|\mathrm{\Delta }(\phi _i\omega )\|\phi _i\|\mathrm{\Delta }\omega \|+2\|\phi _i\|\|\omega \|+\|\mathrm{\Delta }\phi _i\|\|\omega \|,$$ (A.41) and $$\left\|\frac{\phi _i(^2\omega ,\omega )\omega }{R^2+\|\omega \|^2}\right\|\|\phi _i^2\omega \|\|^2(\phi _i\omega )\|+2\|\phi _i\|\|\omega \|+\|^2\phi _i\|\|\omega \|.$$ (A.42) In order to handle the curvature term in (A.40) we observe that $`2\phi _i(\omega _x\omega _y)=[(\phi _i\omega )_x\omega _y+\omega _x(\phi _i\omega )_y][(\phi _i)_x(\omega \omega _y)+(\phi _i)_y(\omega _x\omega )]`$ and hence $`\|2\phi _i(\omega _x\omega _y)\|2\|(\phi _i\omega )\|\|\omega \|+\|\phi _i\|\|\omega \|\|\omega \|`$. Therefore, we can estimate $$\left\|\frac{2\phi _iH_R(\omega )\omega _x\omega _y}{(R^2+\|\omega \|^2)^{\alpha 1}}\right\|2\overline{H}\|(\phi _i\omega )\|\|\omega \|+\overline{H}\|\phi _i\|\|\omega \|\|\omega \|.$$ (A.43) Multiplying (A.40) by $`\phi _i`$ and using (A.41)–(A.43) we obtain $`\|\mathrm{\Delta }(\phi _i\omega )\|`$ $``$ $`2(\alpha 1)\|^2(\phi _i\omega )\|+6K\chi _i(\|\omega \|+\|\omega \|)`$ $`+2\overline{H}\|(\phi _i\omega )\|\|\omega \|+\overline{H}\|\phi _i\|\|\omega \|\|\omega \|`$ where $`\chi _i`$ is the characteristic function of the set $`D_{r_i}`$. Thus, for all $`p(1,+\mathrm{})`$ we have $`\mathrm{\Delta }(\phi _i\omega )_{L^p(D_{r_i})}`$ $``$ $`2(\alpha 1)\phi _i\omega _{H^{2,p}(D_{r_i})}+6K\left(\omega _{L^p(D_{r_i})}+\omega _{L^p(D_{r_i})}\right)`$ (A.44) $`+`$ $`2\overline{H}\|(\phi _i\omega )\|\|\omega \|_{L^p(D_{r_i})}+\overline{H}K\|\omega \|\|\omega \|_{L^p(D_{r_i})}.`$ Since $`\omega `$ has zero mean value on $`D`$, we have that for every $`p(1,+\mathrm{})`$ $$\omega _{L^p(D_{r_i})}C_p\omega _{L^2(D)},$$ (A.45) where $`C_p`$ depends only on the Sobolev embedding constant of $`H^{1,2}(D)`$ into $`L^p(D)`$ and on the Poincaré constant on $`D`$. Taking $`p(1,2]`$, we plainly have $$\omega _{L^p(D_{r_i})}2\omega _{L^2(D)}.$$ (A.46) Moreover, for $`p(1,4)`$, using Hölder inequality and (A.45), we can estimate $`\|(\phi _i\omega )\|\|\omega \|_{L^p(D_{r_i})}(\phi _i\omega )_{L^4(D_{r_i})}\omega _{L^{4p/(4p)}(D_{r_i})},`$ (A.47) $`\|\omega \|\|\omega \|_{L^p(D_{r_i})}C_4\omega _{L^2(D)}\omega _{L^{4p/(4p)}(D_{r_i})}.`$ (A.48) Now we apply the standard regularity theory for linear elliptic equations. Denoting by $`c(p)`$ the norm of the operator $`\mathrm{\Delta }^1`$ as a map from $`L^p(D_{r_i})`$ into $`W^{2,p}H_0^1(D_{r_i})`$, and using (A.44)–(A.48), we obtain the following crucial inequality for $`p(1,2]`$ $`\beta _{p,\alpha }\phi _i\omega _{H^{2,p}(D_{r_i})}`$ $``$ $`\overline{C}_p\omega _{L^2(D)}+C_4\overline{H}K\omega _{L^2(D)}\omega _{L^{4p/(4p)}(D_{r_i})}`$ (A.49) $`+2\overline{H}(\phi _i\omega )_{L^4(D_{r_i})}\omega _{L^{4p/(4p)}(D_{r_i})}.`$ where we have set $`\beta _{p,\alpha }=c(p)^12(\alpha 1)`$ and $`\overline{C}_p=6K(C_p+2)`$. First, we use (A.49) taking $`p=2`$ and $`i=2`$. From (A.45), we have that $`(\phi _2\omega )_{L^4(D_{r_2})}C_4K\omega _{L^2(D)}+\omega _{L^4(D_{r_2})}`$. Then, we fix $`\overline{\alpha }>1`$ such that $`\beta _{2,\overline{\alpha }}>0`$, and we observe that $`\overline{\alpha }`$ depends only on the constants in elliptic regularity theory. Hence, if $`\alpha [1,\overline{\alpha }]`$, (A.49) with $`p=2`$ and $`i=2`$ yields $`\omega _{H^{2,2}(D_{r_1})}`$ $``$ $`\phi _2\omega _{H^{2,2}(D_{r_2})}`$ $``$ $`C_1(\overline{H})\left(\omega _{L^2(D)}+\omega _{L^4(D_{r_2})}\omega _{L^2(D)}+\omega _{L^4(D_{r_2})}^2\right),`$ where $`C_1(\overline{H})`$ depends only on $`\overline{H}`$. Now we show that $`\omega _{L^4(D_{r_2})}`$ can be controlled in terms of $`\omega _{L^2(D)}`$, if $`\omega _{L^2(D)}`$ is small enough. To do this, we use again (A.49) taking $`p=\frac{4}{3}`$ and $`i=3`$. We point out that the critical Sobolev exponent corresponding to $`p=\frac{4}{3}`$ is $`p^{}=4`$. Hence, there exists $`S_{4/3}>0`$ (independent of the domain) such that $`S_{4/3}\phi _3\omega _{H^{1,4}(D)}\phi _3\omega _{H^{2,4/3}(D)}`$. Therefore, reminding that $`r_3=1`$, (A.49), with $`p=\frac{4}{3}`$ and $`i=3`$, yields $$\left(\beta _{\frac{4}{3},\alpha }2\overline{H}S_{4/3}^1\omega _{L^2(D)}\right)\phi _3\omega _{H^{2,4/3}(D)}\overline{C}_{4/3}\omega _{L^2(D)}+C_4\overline{H}K\omega _{L^2(D)}^2.$$ Now, take a smaller $`\overline{\alpha }>1`$ in order that $`\beta _{4/3,\overline{\alpha }}=\overline{\beta }>0`$. Thus, for every $`\alpha [1,\overline{\alpha }]`$ we have $`\beta _{4/3,\alpha }\overline{\beta }`$. Then, take $`\overline{\epsilon }>0`$ small enough, such that $`\overline{\beta }2S_{4/3}^1\overline{H}\overline{\epsilon }>0`$. Notice that $`\overline{\epsilon }`$ depends only on $`\overline{H}`$. Therefore, we infer that $`\omega _{L^4(D_{r_2})}`$ $``$ $`(\phi _3\omega )_{L^4(D)}\phi _3\omega _{H^{1,4}(D)}`$ $``$ $`S_{4/3}^1\phi _3\omega _{H^{2,4/3}(D)}C_2(\overline{H})\omega _{L^2(D)},`$ if $`\omega _{L^2(D)}\overline{\epsilon }`$, with $`C_2(\overline{H})`$ depending only on $`\overline{H}`$. Going back to (A.2), we have proved that $$\omega _{H^{2,2}(D_{r_1})}C_3(\overline{H})\omega _{L^2(D)}$$ when $`\alpha [1,\overline{\alpha }]`$, provided that $`\omega _{L^2(D)}\overline{\epsilon }`$, being $`C_3(\overline{H})`$ a positive constant depending only on $`\overline{H}`$. Hence, by the Sobolev embeddings, for every $`q[1,+\mathrm{})`$ there exists a positive constant $`C_4(q,\overline{H})`$, depending also on $`q`$ such that $$\omega _{H^{1,q}(D_{r_1})}C_4(q,\overline{H})\omega _{L^2(D)}$$ (A.51) when $`\alpha [1,\overline{\alpha }]`$ and $`\omega _{L^2(D)}\overline{\epsilon }`$. For the last step, we apply (A.44) with $`i=1`$ and we use the following estimates, obtained with the Hölder inequality and with (A.45): $`\|(\phi _1\omega )\|\|\omega \|_{L^p(D_{r_1})}`$ $``$ $`K\|\omega \|\|\omega \|_{L^p(D_{r_1})}+\omega _{L^{2p}(D_{r_1})}^2,`$ $`\|\omega \|\|\omega \|_{L^p(D_{r_1})}`$ $``$ $`C_{2p}\omega _{L^2(D)}\omega _{L^{2p}(D_{r_1})}.`$ Then, arguing as for (A.49) we get $`\beta _{p,\alpha }\phi _1\omega _{H^{2,p}(D_{r_1})}`$ $``$ $`6K\omega _{L^2(D)}+6K\omega _{L^p(D_{r_1})}`$ $`+2\overline{H}\omega _{L^{2p}(D_{r_1})}+3\overline{H}KC_{2p}\omega _{L^2(D)}\omega _{L^{2p}(D_{r_1})}.`$ Finally, in order to estimate $`\omega _{L^p(D_{r_1})}`$ and $`\omega _{L^{2p}(D_{r_1})}`$, we use (A.51) with $`q=p`$ and $`q=2p`$. Thus, for fixed $`p(1,+\mathrm{})`$ we can find $`\alpha _p(1,\overline{\alpha }]`$ such that for $`\alpha [1,\alpha _p]`$ one has $`\beta _{p,\alpha }\beta _{p,\alpha _p}>0`$. Moreover, we can also find a constant $`C_5(p,\overline{H})>0`$ such that for $`\alpha [1,\alpha _p]`$, one has $$\omega _{H^{2,p}(D_{1/2})}\phi _1\omega _{H^{2,p}(D_{r_1})}C_5(p,\overline{H})\omega _{L^2(D)}$$ provided that $`\omega _{L^2(D)}\overline{\epsilon }`$. To conclude the proof, we just have to remark that $`\omega _{L^2(D)}=u_{L^2(D_R(z))}`$, and $`u_{H^{1,p}(D_{R/2}(z))}^p=R^{2p}\omega _{L^p(D_{1/2})}^p+R^{22p}^2\omega _{L^p(D_{1/2})}^pR^{22p}\omega _{H^{2,p}(D_{1/2})}^p`$, since $`R1`$. ### A.3 Passing to the limit (blow up analysis for $`(u^\alpha )`$) The first preliminary result concerns the behaviour of the starting sequence $`(u^\alpha )`$ satisfying (A.36)–(A.38). ###### Lemma A.5 $`u^\alpha 0`$ weakly in $`H_0^1`$ and $`u^\alpha _{\mathrm{}}+\mathrm{}`$ as $`\alpha 1`$. Proof. Since $`(u^\alpha )`$ is bounded in $`H_0^1`$ and in $`L^{\mathrm{}}`$, passing to a subsequence, we can assume that $`u^\alpha u`$ weakly in $`H_0^1`$, for some $`uH_0^1L^{\mathrm{}}`$. Let us prove that $`u`$ is a weak solution to the Dirichlet problem $$\{\begin{array}{cc}\mathrm{\Delta }u=2H(u)u_xu_y\hfill & \text{in }D\hfill \\ u=0\hfill & \text{on }D\hfill \end{array}$$ (A.52) To this aim, fix an open set $`\mathrm{\Omega }`$ with $`\overline{\mathrm{\Omega }}D`$. Arguing as in , proof of Proposition 4.3, we can find a finite set of points $`F\mathrm{\Omega }`$ such that $`_{D_R(z)}\|u\|^2\overline{\epsilon }`$ for $`zF`$ and $`R`$ small enough (depending on $`z`$), where $`\overline{\epsilon }>0`$ is given by Lemma A.4. Then, an application of Lemma A.4 gives a uniform bound for $`u^\alpha _{H^{1,2}(D_{R/2}(z))}`$. Noting also that $`(u^\alpha )`$ is bounded in $`L^q(D)`$ for every $`q[1,+\mathrm{}]`$, we infer that $`(u^\alpha )`$ is bounded in $`W^{2,2}(D_{R/2}(z))`$, and hence, by Rellich Theorem, $`u^\alpha u`$ strongly in $`H^1(D_{R/2}(z))`$ and in $`L^{\mathrm{}}(D_{R/2}(z))`$. This is sufficient to conclude that $`u`$ is a weak solution to the equation $`\mathrm{\Delta }u=2H(u)u_xu_y`$ in $`D_{R/2}(z)`$ and hence, since $`z`$ was arbitrarily chosen, in $`\mathrm{\Omega }F`$. Now we can repeat the proof of Theorem 3.6 in . Assume for simplicity that $`F=\{0\}`$. Let $`\eta C^{\mathrm{}}(,[0,1])`$ be such that $`\eta (s)=0`$ for $`s1`$ and $`\eta (s)=1`$ for $`s2`$, and set $`\eta ^k(s)=\eta (ks)`$. Given $`hC_c^{\mathrm{}}(\mathrm{\Omega },^3)`$ we set $`h^k(\zeta )=\eta ^k(\|\zeta \|)h(\zeta )`$. Notice that $`h^k`$ can be used as test for $`u`$ to get $$_\mathrm{\Omega }uh^k+2_\mathrm{\Omega }H(u)h^ku_xu_y=0.$$ (A.53) Now, since $`h^kh`$ weakly in $`L^{\mathrm{}}`$, we get $`_\mathrm{\Omega }H(u)h^ku_xu_y_\mathrm{\Omega }H(u)hu_xu_y`$. Also, $`_\mathrm{\Omega }uh^k_\mathrm{\Omega }uh`$, since, by Hölder inequality, $`_\mathrm{\Omega }\|u\eta ^k\|\|h\|Cu_{L^2(D_{2/k})}=o(1)`$ as $`k+\mathrm{}`$. Therefore, (A.53) yields in the limit $$_\mathrm{\Omega }uh+2_\mathrm{\Omega }H(u)hu_xu_y=0$$ for every test function $`hC_c^{\mathrm{}}(\mathrm{\Omega },^3)`$, that is, $`u`$ solves $`\mathrm{\Delta }u=2H(u)u_xu_y`$ in $`\mathrm{\Omega }`$. Finally, for the arbitrariness of $`\mathrm{\Omega }`$, we conclude that $`u`$ is a weak solution to problem (A.52). Then, by a Heinz regularity result , $`u`$ is smooth, and a nonexistence result by Wente , which holds also in case $`H`$ nonconstant, can be applied, to conclude that $`u0`$. Thus, we obtain that $`u^\alpha 0`$ weakly in $`H_0^1`$ and strongly in $`H_{loc}^1(DN)`$ where $`N`$ is a countable set of $`D`$. In particular $`u^\alpha 0`$ pointwise a.e. in $`D`$. Therefore, as a last step, we observe that if it were $`lim\; inf_{\alpha 1}u^\alpha _{\mathrm{}}<+\mathrm{}`$, then $`lim\; inf_{\alpha 1}u^\alpha _2=0`$, contrary to (A.38). Hence, it must be $`u^\alpha _{\mathrm{}}+\mathrm{}`$ as $`\alpha 1`$. Proof of Proposition A.1. For every $`\alpha (1,\overline{\alpha })`$ set $`ϵ_\alpha =u^\alpha _{\mathrm{}}^1`$, let $`z_\alpha \overline{D}`$ be such that $`\|u^\alpha (z_\alpha )\|=ϵ_\alpha ^1`$ and define $$v^\alpha (z)=u^\alpha (ϵ_\alpha z+z_\alpha ).$$ (A.54) Note that $`v^\alpha H_0^1(D_\alpha ,^3)`$ where $`D_\alpha =D_{ϵ_\alpha ^1}\left(\frac{z_\alpha }{ϵ_\alpha }\right)`$. Moreover the following facts hold: $`v^\alpha _{\mathrm{}}=u^\alpha _{\mathrm{}}`$ (A.55) $`v^\alpha _2=u^\alpha _2`$ (A.56) $`\|v^\alpha (0)\|=v^\alpha _{\mathrm{}}=1.`$ (A.57) Furthermore, $`v^\alpha W_{loc}^{2,q}(D_\alpha ,^3)`$ for every $`q[1,+\mathrm{})`$ and solves the system $$\mathrm{\Delta }v^\alpha =\frac{2(\alpha 1)}{ϵ_\alpha ^2+\|v^\alpha \|^2}(^2v^\alpha ,v^\alpha )v^\alpha +\frac{2ϵ_\alpha ^{2(\alpha 1)}H(v^\alpha )}{(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}}v_x^\alpha v_y^\alpha \mathrm{in}D_\alpha .$$ (A.58) Since $`ϵ_\alpha 0`$ as $`\alpha 1`$, one has that $`0<ϵ_\alpha ^{2(\alpha 1)}<1`$, and then, for a subsequence, $`ϵ_\alpha ^{2(\alpha 1)}\lambda `$ for some $`\lambda [0,1]`$. Moreover, setting $`\rho _\alpha =ϵ_\alpha ^1\mathrm{dist}(z_\alpha ,D)`$, we may also assume that there exists $`lim_{\alpha 1}\rho _\alpha [0,+\mathrm{}]`$. Let $`\mathrm{\Omega }_{\mathrm{}}`$ be the union of all compact sets in $`^2`$ contained in $`D_\alpha `$ as $`\alpha 1`$. Note that $`\mathrm{\Omega }_{\mathrm{}}`$ is a half-plane if $`\rho _\alpha \mathrm{}[0,+\mathrm{})`$, while $`\mathrm{\Omega }_{\mathrm{}}=^2`$ if $`\rho _\alpha +\mathrm{}`$. From (A.37), (A.55) and (A.56) it follows that there exists $`\omega XL^{\mathrm{}}`$ such that, for a subsequence, $`v^\alpha \omega `$ weakly in $`(L^2(^2,^3))^2`$. Moreover, by (A.57) one has that $`v^\alpha \omega `$ strongly in $`L_{loc}^{\mathrm{}}(^2,^3)`$. Let $`\overline{\epsilon }>0`$ be given by Lemma A.4. Take an arbitrary compact set $`K`$ in $`\mathrm{\Omega }_{\mathrm{}}`$ and set $`R_K=\mathrm{dist}(K,\mathrm{\Omega }_{\mathrm{}})`$. Then, let $`R(0,\mathrm{min}\{1,R_K,\frac{\overline{\epsilon }}{\sqrt{\pi }}\})`$. Hence, there exists $`\alpha _K>1`$ such that $`KD_\alpha `$ for $`\alpha (1,\alpha _K)`$ and, consequently, for every $`zK`$, one has $`\overline{D_R(z)}D_\alpha `$ and $`v^\alpha _2\overline{\epsilon }`$. Because of the definition (A.54) of $`v^\alpha `$, one can apply Lemma A.4, in order to conclude that $`v^\alpha _{H^{1,p}(D_{R/2}(z))}`$ is uniformly bounded with respect to $`\alpha (1,\alpha _K)`$, for every $`p>1`$. Using (A.55) and (A.37), we infer that $`(v^\alpha )`$ is bounded in $`H^{2,p}(D_{R/2}(z))`$. Therefore we can conclude that $`\omega H^{2,p}(D_{R/2}(z))`$, $`v^\alpha \omega `$ strongly in $`H^1(D_{R/2}(z))`$, and $`v^\alpha \omega `$ pointwise everywhere in $`D_{R/2}(z)`$. Since $`z`$ is an arbitrary point in $`K`$ and $`K`$ is any compact set in $`\mathrm{\Omega }_{\mathrm{}}`$, a standard diagonal argument yields that $`\omega H_{loc}^{2,p}(\mathrm{\Omega }_{\mathrm{}})`$ for every $`p<+\mathrm{}`$, $`v^\alpha \omega `$ strongly in $`H_{loc}^1(\mathrm{\Omega }_{\mathrm{}})`$, and $`v^\alpha \omega `$ pointwise everywhere in $`^2`$. In particular, by (A.57), $`\|\omega (0)\|=\omega _{\mathrm{}}=1`$, and thus $`\omega `$ is nonconstant. Now we test (A.58) on an arbitrary function $`hC_c^{\mathrm{}}(D_{R/2}(z),^3)`$ and we pass to the limit as $`\alpha 1`$. First, we have $$_^2\mathrm{\Delta }v^\alpha h_^2\omega h,$$ (A.59) because of the weak convergence $`v^\alpha \omega `$. Secondly, using the estimate $$\left\|_^2\frac{(^2v^\alpha ,v^\alpha )v^\alpha h}{ϵ_\alpha ^2+\|v^\alpha \|^2}\right\|_^2\|^2v^\alpha \|\|h\|^2v^\alpha _{L^p(D_{R/2}(z))}h_{L^p^{}}$$ and the fact that $`v^\alpha `$ is uniformly bounded in $`H^{2,p}(D_{R/2}(z))`$ as $`\alpha (1,\alpha _K)`$, we obtain that $$2(\alpha 1)_^2\frac{(^2v^\alpha ,v^\alpha )v^\alpha h}{ϵ_\alpha ^2+\|v^\alpha \|^2}0,$$ (A.60) as $`\alpha 1`$. Lastly, setting $$w^\alpha =ϵ_\alpha ^{2(\alpha 1)}\left(\frac{1}{(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}}1\right)v_x^\alpha v_y^\alpha $$ one has $$_^2\frac{ϵ_\alpha ^{2(\alpha 1)}H(v^\alpha )}{(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}}hv_x^\alpha v_y^\alpha =_^2H(v^\alpha )hw^\alpha +ϵ_\alpha ^{2(\alpha 1)}_^2H(v^\alpha )hv_x^\alpha v_y^\alpha .$$ Since $`ϵ_\alpha ^{2(\alpha 1)}\lambda `$, $`H(v^\alpha )H(\omega )`$ uniformly on $`\overline{D_{R/2}(z)}`$ and $`v^\alpha \omega `$ pointwise in $`D_{R/2}(z)`$, by (A.57), on one hand we infer that $$ϵ_\alpha ^{2(\alpha 1)}_^2H(v^\alpha )hv_x^\alpha v_y^\alpha \lambda _^2H(\omega )h\omega _x\omega _y.$$ On the other hand, since $`ϵ_\alpha (0,1)`$, we observe that $$\|w^\alpha \|(1+ϵ_\alpha ^{2(\alpha 1)})\|v_x^\alpha \|\|v_y^\alpha \|\|v^\alpha \|^21$$ and $`w^\alpha (\zeta )0`$ for every $`\zeta D_{R/2}(z)`$. Indeed, if $`\omega (\zeta )=0`$ then $`\|w^\alpha \|\|v^\alpha \|^20`$, while if $`\omega (\zeta )0`$ then $`(ϵ_\alpha ^2+\|v^\alpha (\zeta )\|^2)^{\alpha 1}0`$. In conclusion, by the dominated convergence Theorem, we obtain that $`_^2H(v^\alpha )hw^\alpha 0`$ and then $$_^2\frac{2ϵ_\alpha ^{2(\alpha 1)}H(v^\alpha )}{(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}}hv_x^\alpha v_y^\alpha 2\lambda _^2H(\omega )h\omega _x\omega _y,$$ (A.61) as $`\alpha 1`$. Then (A.58)–(A.61) imply that $$_^2\omega h+2\lambda _^2H(\omega )h\omega _x\omega _y=0$$ for every $`hC_c^{\mathrm{}}(D_{R/2}(z),^3)`$, for every $`zK`$ and for every compact set $`K`$ in $`\mathrm{\Omega }_{\mathrm{}}`$, that is, $`\omega `$ solves $`\mathrm{\Delta }\omega =2\lambda H(\omega )\omega _x\omega _y`$ in $`\mathrm{\Omega }_{\mathrm{}}`$. Suppose that $`\mathrm{\Omega }_{\mathrm{}}`$ is a half-plane. Since $`v^\alpha =0`$ on $`D_\alpha `$, one has that $`\omega =0`$ on $`\mathrm{\Omega }_{\mathrm{}}`$. Moreover, since a half-plane is conformally equivalent to a disc, $`\omega `$ gives arise to a nonconstant solution to the Dirichlet problem $$\{\begin{array}{cc}\mathrm{\Delta }u=2\lambda H(u)u_xu_y\hfill & \text{in }D\hfill \\ u=0\hfill & \text{on }D.\hfill \end{array}$$ (A.62) As already noted in the proof of Lemma A.5, the only solution to (A.62) is $`u0`$, and this gives a contradiction, since $`\omega `$ is nonconstant. Hence, it must be $`\mathrm{\Omega }_{\mathrm{}}=^2`$, that is, $`\omega `$ is a $`\lambda H`$-bubble. Finally, we observe that $`\lambda >0`$, since the only bounded solutions to $`\mathrm{\Delta }u=0`$ on $`^2`$ with $`𝒟(u)<+\mathrm{}`$ are the constant functions, and we already know that $`\omega `$ is nonconstant. This concludes the proof. ### A.4 Proof of Proposition A.2 For every domain $`\mathrm{\Omega }`$ in $`^2`$, $`\alpha (1,\overline{\alpha })`$, and $`\lambda (0,1]`$, set $`\stackrel{~}{}_H^\alpha (v^\alpha ,\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{2\alpha }}{\displaystyle _\mathrm{\Omega }}\left((ϵ_\alpha ^2+\|v^\alpha \|^2)^\alpha ϵ_\alpha ^{2\alpha }\right)+2ϵ_\alpha ^{2(\alpha 1)}{\displaystyle _\mathrm{\Omega }}m_H(v^\alpha )v^\alpha v_x^\alpha v_y^\alpha `$ $`_{\lambda H}(\omega ,\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Omega }}\|\omega \|^2+2\lambda {\displaystyle _\mathrm{\Omega }}m_H(\omega )\omega \omega _x\omega _y.`$ Notice that $`\stackrel{~}{}_H^\alpha (v^\alpha ,D_\alpha )=ϵ_\alpha ^{2(\alpha 1)}_H^\alpha (u^\alpha )`$ and $`v^\alpha `$ solves the system $$\mathrm{div}\left((ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}v^\alpha \right)=2ϵ_\alpha ^{2(\alpha 1)}H(v^\alpha )v_x^\alpha v_y^\alpha .$$ (A.63) Now, multiplying (A.63) by $`v^\alpha `$, we obtain $`\mathrm{div}\left((ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}v^\alpha v^\alpha \right)`$ $`=`$ $`(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}\|v^\alpha \|^2`$ (A.64) $`+2ϵ_\alpha ^{2(\alpha 1)}H(v^\alpha )v^\alpha v_x^\alpha v_y^\alpha .`$ Integrating (A.64) on a domain $`\mathrm{\Omega }`$ and using the divergence theorem we infer that $`{\displaystyle _\mathrm{\Omega }}(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}v^\alpha {\displaystyle \frac{v^\alpha }{\nu }}`$ $`=`$ $`{\displaystyle _\mathrm{\Omega }}(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}\|v^\alpha \|^2`$ (A.65) $`+2ϵ_\alpha ^{2(\alpha 1)}{\displaystyle _\mathrm{\Omega }}H(v^\alpha )v^\alpha v_x^\alpha v_y^\alpha .`$ Using (2.10) and the definition of $`\stackrel{~}{}_H^\alpha (v^\alpha ,\mathrm{\Omega })`$ we can estimate $`2ϵ_\alpha ^{2(\alpha 1)}{\displaystyle _\mathrm{\Omega }}H(v^\alpha )v^\alpha v_x^\alpha v_y^\alpha `$ $``$ $`ϵ_\alpha ^{2(\alpha 1)}{\displaystyle \frac{\overline{M}_H}{2}}{\displaystyle _\mathrm{\Omega }}\|v^\alpha \|^2+3\stackrel{~}{}_H^\alpha (v^\alpha ,\mathrm{\Omega })`$ $`{\displaystyle \frac{3}{2\alpha }}{\displaystyle _\mathrm{\Omega }}\left((ϵ_\alpha ^2+\|v^\alpha \|^2)^\alpha ϵ_\alpha ^{2\alpha }\right).`$ (A.66) Hence, setting $`I_\alpha (\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle _\mathrm{\Omega }}(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}v^\alpha {\displaystyle \frac{v^\alpha }{\nu }}`$ $`I_\alpha (\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{2\alpha }}{\displaystyle _\mathrm{\Omega }}\left((ϵ_\alpha ^2+\|v^\alpha \|^2)^\alpha ϵ_\alpha ^{2\alpha }\right){\displaystyle \frac{1}{3}}{\displaystyle _\mathrm{\Omega }}(ϵ_\alpha ^2+\|v^\alpha \|^2)^{\alpha 1}\|v^\alpha \|^2`$ $`ϵ_\alpha ^{2(\alpha 1)}{\displaystyle \frac{\overline{M}_H}{6}}{\displaystyle _\mathrm{\Omega }}\|v^\alpha \|^2,`$ by (A.66) the equation (A.65) becomes $$\stackrel{~}{}_H^\alpha (v^\alpha ,\mathrm{\Omega })I_\alpha (\mathrm{\Omega })+I_\alpha (\mathrm{\Omega }).$$ (A.67) With algebraic computations, one has $`I_\alpha (\mathrm{\Omega })`$ $``$ $`\left({\displaystyle \frac{1}{2\alpha }}{\displaystyle \frac{1}{3}}\right){\displaystyle _\mathrm{\Omega }}\left((ϵ_\alpha ^2+\|v^\alpha \|^2)^\alpha ϵ_\alpha ^{2\alpha }\right)ϵ_\alpha ^{2(\alpha 1)}{\displaystyle \frac{\overline{M}_H}{6}}{\displaystyle _\mathrm{\Omega }}\|v^\alpha \|^2`$ $``$ $`ϵ_\alpha ^{2(\alpha 1)}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{\alpha }{3}}{\displaystyle \frac{\overline{M}_H}{6}}\right){\displaystyle _\mathrm{\Omega }}\|v^\alpha \|^2.`$ Then, since $`\overline{M}_H<1`$, one obtains that $`I_\alpha (\mathrm{\Omega })0`$ for $`\alpha >1`$ sufficiently close to 1, whatever $`\mathrm{\Omega }`$ is. Hence, (A.67) reduces to $$\stackrel{~}{}_H^\alpha (v^\alpha ,\mathrm{\Omega })I_\alpha (\mathrm{\Omega }).$$ (A.68) Now we take $`\mathrm{\Omega }=^2D_R`$. First, we observe that, since $`v^\alpha \omega `$ strongly in $`H_{loc}^1(^2,^3)`$ and uniformly on compact sets, and $`ϵ_\alpha ^{2(\alpha 1)}\lambda `$, it holds that $`\underset{\alpha 1}{lim}\stackrel{~}{}_H^\alpha (v^\alpha ,D_R)`$ $`=`$ $`_{\lambda H}(\omega ,D_R)`$ $`\underset{\alpha 1}{lim\; sup}\left\|I_\alpha (D_R)\right\|`$ $``$ $`{\displaystyle \frac{1}{3}}\left\|{\displaystyle _{D_R}}\omega {\displaystyle \frac{\omega }{\nu }}\right\|`$ for every $`R>0`$. Then, by (A.68), we obtain $`_{\lambda H}(\omega ,D_R)`$ $`=`$ $`\stackrel{~}{}_H^\alpha (v^\alpha )\stackrel{~}{}_H^\alpha (v^\alpha ,^2D_R)+o(1)`$ $``$ $`ϵ_\alpha ^{2(\alpha 1)}_H^\alpha (u^\alpha )+{\displaystyle \frac{1}{3}}\left\|{\displaystyle _{D_R}}\omega {\displaystyle \frac{\omega }{\nu }}\right\|+o(1)`$ where $`o(1)0`$ as $`\alpha 1`$, for every $`R>0`$. Hence $$\lambda \underset{\alpha 1}{lim\; inf}_H^\alpha (u^\alpha )_{\lambda H}(\omega ,D_R)\frac{1}{3}\left\|_{D_R}\omega \frac{\omega }{\nu }\right\|$$ (A.69) for every $`R>0`$. Finally, notice that $`\left\|{\displaystyle _{D_R}}\omega {\displaystyle \frac{\omega }{\nu }}\right\|`$ $`=`$ $`\left\|{\displaystyle _{^2D_R}}\left(\omega \mathrm{\Delta }\omega +\|\omega \|^2\right)\right\|`$ $`=`$ $`\left\|{\displaystyle _{^2D_R}}\left(2\lambda H(\omega )\omega \omega _x\omega _y+\|\omega \|^2\right)\right\|`$ $``$ $`\left(\lambda H_{\mathrm{}}\omega _{\mathrm{}}+1\right){\displaystyle _{^2D_R}}\|\omega \|^2.`$ Then, passing to the limit as $`R+\mathrm{}`$, from (A.69) the thesis follows. References
warning/0006/hep-ph0006065.html	ar5iv	text	# 1 Introduction ## 1 Introduction It is fair to say that, at present, the Standard Model belongs to the category of most thoroughly tested and best confirmed physical theories . But the secret is that no one truly understands it. At least not in the way we understand QED or QCD. It appears that the vast number of free parameters, speculations about the number of Higgs doublets, and the ad-hoc definition of the Higgs potential defy easy analysis. With each new generation of fundamental fermions the unknowns have multiplied. The neutrinos may or may not have masses. It remains undecided whether nature provides another CKM matrix for the leptons. In the long run, the struggle for a better understanding will perhaps be resolved by resorting to string theory (M-theory or other oracles). In the meantime we might be content with modest explanations using constructions in ordinary (commutative) differential geometry. One concept, which convincingly illuminates the role of the Higgs field, comes under the heading superconnection. Implied is the concept of a generalized Dirac operator which unites gauge and Yukawa couplings in one term. In we showed how the Higgs field fits into the framework of superconnections on some superbundle with structure group $`U(n)`$. We then proposed to take $`n=2`$ to construct a model of gauge fields and two Higgs doublets. In we added leptons to the model to see how gauge and Yukawa interactions can be combined by passing from a superconnection to the associated Dirac operator. Constructing such a model was a tentative step towards an understanding of the structure of more realistic theories. Of course, without quarks the leptonic model was not free of anomalies and thus called for an extension incorporating essential features of the Standard Model. A short account of such an attempt appeared in where the gauge group $`G`$ was assumed to be a subgroup of $`SU(5)`$. Parallel to this work we opened a discussion in and on the mathematical background which should serve as a reference when we now resume the analysis of the Standard Model begun in . Moreover, since we are dealing here with a chiral model, the question of consistency (absence of local nonabelian anomalies) arises. This problem has been dealt with and settled in . The present exposition of the subject aims to provide motivation and practical tools rather than mathematical abstraction. Its goal is to arrive at predictions with minimal technical machinery. ## 2 The Gauge Group of the Standard Model The Standard Model, extending of the earlier Weinberg-Salam model of electroweak interactions, is a gauge theory based on the Lie algebra $$\text{Lie }G\mathrm{𝐬𝐮}(3)\mathrm{𝐬𝐮}(2)𝐮(1).$$ (1) Though the gauge group $`G`$ is well defined locally, its global structure remains obscure unless we add further assumptions. Obviously, no restriction on the spectrum of the hypercharge $`Y`$ is to be expected on the basis of $`𝐮(1)=i`$ alone. Nor would one be able to argue that the conditions $`Q+Y`$ and $`3Y`$ are satisfied where $`Q`$ denotes the electric charge. Still, in the vicinity of the unit, a possible difference between various choices of $`G`$ would not be felt at all, but globally it would: $`G`$ will dictate the subset of allowed representations of $`\text{Lie }G`$ and, therefore, the structure of particle multiplets admitted by the setup. The present approach borrows from the idea that any grand unified theory, perhaps any theory beyond the Standard Model, ought to encorporate the gauge group $`SU(5)`$ in one way or another. As for the minimal version of the Standard Model, we require that $`G`$ be a subgroup of $`SU(5)`$ consistent with (1), i.e., we define $$G=\{(u,v)U(3)\times U(2)\|detudetv=1\}$$ (2) and let the embedding $`GSU(5)`$ be given by $$(u,v)\left(\begin{array}{cc}u& 0\\ 0& v\end{array}\right).$$ It is easy to see that the Lie algebra $`\text{Lie }G`$ has indeed the required structure given by (1). The relation of the group $`G`$ to the color group $`SU(3)`$ and the electroweak group $`U(2)`$ is expressed by the following exact sequence $$1SU(3)\stackrel{j}{}G\stackrel{s}{}U(2)1$$ (3) where $`j(u)=(u,1)`$ and $`s(u,v)=v`$. Inspite of the relationship (3), the group $`G`$ cannot be identified with the direct product $`SU(3)\times U(2)`$. It is still correct to say that the color group $`SU(3)`$ of quantum chromodynamics is embedded in $`G`$ as a subgroup. But the gauge group $`U(2)`$ of the Salam-Weinberg Theory is recovered here only as the quotient $`G/SU(3)`$. This fact deserves careful attention: it explains why the hypercharge $`Y`$ assumes fractional values. More specifically, the quotient structure accounts for the existence of values that are multiples of 1/3. The theory of leptons, gauge and Higgs particles is built around the assumption that the hypercharge $`Y`$ is the generator of $`U(1)`$, subgroup and center of $`U(2)`$. As these groups constitute proper symmetries (before spontaneous symmetry breaking takes place), the hypercharge is integer-valued for all leptonic states. When quarks are added, the picture changes. With quarks, the group $`U(1)`$ fails to be a subgroup of $`G`$ and hence cannot be regarded a symmetry though there is a related group $`\stackrel{~}{U}(1)`$ which can. To see more clearly the emergence of a fractional spectrum we consider the exact sequence $$1_3\stackrel{j}{}\stackrel{~}{U}(1)\stackrel{s}{}U(1)1$$ (4) obtained from the groups in (3) by restricting to the centers. In more detail: * The center $`\stackrel{~}{U}(1)`$ of $`G`$ consists of elements $$\text{diag}(e^{i\beta },e^{i\beta },e^{i\beta },e^{i\alpha },e^{i\alpha })SU(5)$$ satisfying $`(e^{i\beta })^3(e^{i\alpha })^2=1`$. It may conveniently be looked upon as a one-dimensional closed subgroup of the two-torus: $$\stackrel{~}{U}(1)=\{(e^{i\beta },e^{i\alpha })\|3\beta +2\alpha =0mod2\pi \}.$$ * The cyclic group $`_3`$ of order three is formed by the complex solutions of $`z^3=1`$. The injection $`j`$ takes $`z`$ into $`(z,1)\stackrel{~}{U}(1)`$. * The surjection $`s`$ maps $`(e^{i\beta },e^{i\alpha })`$ to $`e^{i\alpha }U(1)`$ where $`U(1)`$ relates to the hypercharge. Viewed geometrically, the group $`\stackrel{~}{U}(1)`$ describes a closed curve on the 2-torus. Suppose we unwind the torus to obtain its covering plane with real coordinates $`\alpha `$ and $`\beta `$. Then the curve appears as a straight line with slope parameter $`2/3`$: The end points of the line have to be identified to form a closed curve on the torus. As we run once through this curve, the angle $`\alpha `$ assumes all values from 0 to $`6\pi `$. Phrased more formally, $`\stackrel{~}{U}(1)`$ is a threefold cover of the group $`U(1)=\{e^{i\alpha }\}`$. When it comes to particle multiplets, we must focus on the symmetry group $`\stackrel{~}{U}(1)`$ since its generator has an integer-valued spectrum. In unitary irreducible representations of the gauge group $`G`$, the hypercharge $`Y`$ assumes a constant value subject to the constraint $$e^{i6\pi Y}=1\text{oder}3Y,$$ as can be inferred from the behavior of the variable $`\alpha `$. In other words, the group $`\stackrel{~}{U}(1)`$ is not connected in a direct manner with $`Y`$ but rather with $`3Y`$. Phrased more formally, the covering map $`s:\stackrel{~}{U}(1)U(1)`$ admits a local inverse $$s^1(e^{i\alpha })=(e^{i2\alpha /3},e^{i\alpha }).$$ (5) and hence, at least locally (for small $`\alpha `$), the group $`U(1)`$ is represented by a phase factor $`e^{i\alpha Y}`$ in any unitary irreducible representation of $`G`$, in such a way that $`3Y`$ becomes an integer. In the leptonic sector, spontaneous symmetry breaking selects another one-parameter subgroup of $`U(2)`$ which remains unbroken and gives rise to the concept of electric charge. By convention, this subgroup is $$U(1)_Q=\left\{\left(\begin{array}{cc}1& 0\\ 0& e^{i\alpha }\end{array}\right)\right\}$$ (6) Therefore, the charge $`Q`$ assumes integer values in this sector. Again, with quarks the situation changes. The group $`U(1)_Q`$ being no longer a symmetry is replaced by its threefold cover, $$\stackrel{~}{U}(1)_Q=\left\{(u,v)\right\|u=e^{i\beta }1\mathrm{l}_3,v=\left(\begin{array}{cc}1& 0\\ 0& e^{i\alpha }\end{array}\right),3\beta +\alpha =0mod2\pi \}$$ (7) Locally, we may put $`\beta =\alpha /3`$ and $`\stackrel{~}{U}(1)_Q=e^{i\alpha Q}`$ such that $`3Q`$. Summarizing, the hypercharge and the electric charge are represented on $`^5`$ by the following traceless matrices: $`Y`$ $`=`$ $`\text{diag}(\frac{2}{3},\frac{2}{3},\frac{2}{3},1,1)`$ $`Q`$ $`=`$ $`\text{diag}(\frac{1}{3},\frac{1}{3},\frac{1}{3},0,1)`$ Though many states, especially those invariant under the color group $`SU(3)`$, carry integer charges, $`Q`$ and $`Y`$ assume fractional values in general, still obeying $`Q+Y`$. ## 3 The $`G`$-Supermodule of Fermions By construction, there is a natural irreducible unitary action of the gauge group $`G`$ on the space $$^5=^3^2$$ (8) with subspaces $`^3`$ and $`^2`$ carrying fundamental representations of the color group SU(3) and the weak-isospin group $`SU(2)`$ respectively. But, passage to the $`_2`$-graded exterior algebra $$^5=\underset{k=0}{\overset{5}{}}^k^5=^+^5^{}^5,^\pm ^5=\underset{(1)^k=\pm 1}{}^k^5$$ is very essential if we want to let $`G`$ act on a superspace. For a general discussion and details concerning the exterior algebra as a superspace we refer to . It is apparent that the induced unitary representation $``$ of $`G`$ on $`^5`$ is reducible. We take the view that $`^5`$ is the basic $`G`$-module for fermions. Quarks and leptons of one generation will be grouped according to the irreducible constituents of the representation $``$. In addition, there is another space $`^3`$, not a $`G`$-module, which describes the flavor degrees of freedom. While we refer to $`^5`$ as the $`G`$-supermodule of fermions, we call the tensor product $$^5^3$$ the inner space, because it incorporates all inner degrees of freedom. We now turn to the structure of the $`G`$-supermodule. From (8) and the natural isomorphism (of vector spaces) $$(^3^2)^3^2,$$ where $$^3=_{p=0}^3^p^3,^2=_{q=0}^2^q^3,$$ we obtain $`(u,v)=uv`$ for $`(u,v)G`$ and hence $$^k(u,v)=_{p+q=k}^pu^qv,k=0,\mathrm{},5.$$ We call $$\kappa =(1)^k=(1)^{p+q}$$ (9) the parity operator in $`^5`$. Within a generation of fermions, each multiplet (left- or right-handed) is associated with one of the following irreducible representations of $`G`$, $$^{p,q}=^p^qp=0,1,2,3,q=0,1,2$$ whose dimension is $`\left(\genfrac{}{}{0pt}{}{p}{3}\right)\left(\genfrac{}{}{0pt}{}{q}{2}\right)`$. To find its hypercharge we use Eq. (5), $$e^{i\alpha Y}=^{p,q}\left(s^1(e^{i\alpha })\right)=\mathrm{exp}(i2p\alpha /3+iq\alpha ),$$ and thus obtain the fundamental relation $$Y=\frac{2}{3}pq.$$ (10) We distinguish \| lepton fields \| : $`p=0`$ or 3 \| \| --- \| --- \| \| quark fields \| : $`p=1`$ or 2. \| The electric charge then satisfies the formula of Gell-Mann-Nishijima $$Q=I_3+\frac{1}{2}Y$$ where $`I_3`$ denotes the third component of the weak isospin. Clearly, $`I_3=0`$ if $`q=0,2`$ and $`I_3=\pm \frac{1}{2}`$ if $`q=1`$. With each generation we associate a generalized Dirac field $`\psi `$ having $`2^5=32`$ elementary Weyl spinors as its components. Spinors that enter $`\psi `$ are characterized by three different “parities” owing to the $`_2`$-gradings of $`^5`$, $`^3`$, and $`^2`$. Their interpretation is as follows (recall that $`k=p+q)`$): \| $`k=`$even \| : right-handed \| $`k=`$odd \| : left-handed \| \| --- \| --- \| --- \| --- \| \| $`p=`$even \| : matter \| $`p=`$odd \| : antimatter \| \| $`q=`$even \| : singlets \| $`q=`$odd \| : doublets. \| Since charge conjugation acts on the inner space by complex conjugation, it passes from $`[p,q]`$ to $`[3p,2q]`$. It thus interchanges left and right, matter and antimatter, and reverses the signs of $`Y`$, $`I_3`$ and $`Q`$, but takes singlets into singlets and doublets into doublets. ## 4 Choosing a Basis in $`^5`$ To describe the field $`\psi `$ in more conventional terms we need to construct a basis of eigenvectors in $`^5`$. Such a construction starts from a basis in $`^5`$. Let us emphazise: besides being a 5-dimensional complex linear space, $`^5`$ is endowed with a Hermitian structure and also comes with a distinguished orthonormal basis $`(e_i)_{i=1}^5`$ such that $`(e_i)_j=\delta _{ij}`$. We shall refer to it as the standard basis in $`^5`$. An induced basis $`e_I`$ in $`^5`$ is then given by $$e_I=e_{i_1}\mathrm{}e_{i_k}^k^5,I=\{i_1,\mathrm{}i_k\},i_1<\mathrm{}<i_k,0k5$$ where $`I`$ runs over all subsets of $`\{1,2,3,4,5\}`$ including the empty set $`\mathrm{}`$. We simply have to identify the numbers $`k,p,q`$ previously introduced. Given any subset $`I`$, * $`k`$ is the number of elements in $`I`$ taken from $`\{1,2,3,4,5\}`$, * $`p`$ is the number of elements in $`I`$ taken from $`\{1,2,3\}`$, * $`q`$ is the number of elements in $`I`$ taken from $`\{4,5\}`$. We shall also write $`\|I\|`$ in place of $`k`$, and the complement of $`I`$ in $`\{1,2,3,4,5\}`$ is denoted $`I^c`$. The following table lists all 32 basis vectors and groups them according to their $`p`$ and $`q`$ values. For convenience, we write $`I`$ where we really mean $`e_I`$. \| $`I`$ \| $`q=0`$ \| $`q=2`$ \| $`q=1`$ \| \| \| --- \| --- \| --- \| --- \| --- \| \| $`p=0`$ \| $`\mathrm{}`$ \| 45 \| 4 \| 5 \| \| \| 23 \| 2345 \| 234 \| 235 \| \| $`p=2`$ \| 13 \| 1345 \| 134 \| 135 \| \| \| 12 \| 1245 \| 124 \| 125 \| \| \| 1 \| 145 \| 14 \| 15 \| \| $`p=1`$ \| 2 \| 245 \| 24 \| 25 \| \| \| 3 \| 345 \| 34 \| 35 \| \| $`p=3`$ \| 123 \| 12345 \| 1234 \| 1235 \| (11) The basis vectors of the first two columns have $`I_3=0`$ while those of the third and fourth column have $`I_3=\frac{1}{2}`$ and $`I_3=\frac{1}{2}`$ respectively. By construction, each basis vector $`e_I`$ is an eigenvector of $`Y`$ and $`Q`$. The following two tables provide the hypercharges (left table) and the electric charges (right table) associated with the basis vectors: \| $`Y`$ \| $`q=0`$ \| $`q=2`$ \| $`q=1`$ \| \| \| --- \| --- \| --- \| --- \| --- \| \| $`p=0`$ \| $`0`$ \| -2 \| -1 \| -1 \| \| \| 4/3 \| -2/3 \| 1/3 \| 1/3 \| \| $`p=2`$ \| 4/3 \| -2/3 \| 1/3 \| 1/3 \| \| \| 4/3 \| -2/3 \| 1/3 \| 1/3 \| \| \| 2/3 \| -4/3 \| -1/3 \| -1/3 \| \| $`p=1`$ \| 2/3 \| -4/3 \| -1/3 \| -1/3 \| \| \| 2/3 \| -4/3 \| -1/3 \| -1/3 \| \| $`p=3`$ \| 2 \| 0 \| 1 \| 1 \| \| $`Q`$ \| $`q=0`$ \| $`q=2`$ \| $`q=1`$ \| \| \| --- \| --- \| --- \| --- \| --- \| \| $`p=0`$ \| $`0`$ \| -1 \| 0 \| -1 \| \| \| 2/3 \| -1/3 \| 2/3 \| -1/3 \| \| $`p=2`$ \| 2/3 \| -1/3 \| 2/3 \| -1/3 \| \| \| 2/3 \| -1/3 \| 2/3 \| -1/3 \| \| \| 1/3 \| -2/3 \| 1/3 \| -2/3 \| \| $`p=1`$ \| 1/3 \| -2/3 \| 1/3 \| -2/3 \| \| \| 1/3 \| -2/3 \| 1/3 \| -2/3 \| \| $`p=3`$ \| 1 \| 0 \| 1 \| 0 \| After this assignment of charges, the basis vectors $`e_I`$ can be put in a 1:1 correspondence with the 32 Weyl spinors of the first generation: \| 1. gener. \| $`q=0`$ \| $`q=2`$ \| $`q=1`$ \| \| \| --- \| --- \| --- \| --- \| --- \| \| $`p=0`$ \| $`\nu _{eR}`$ \| $`e_R`$ \| $`\nu _{eL}`$ \| $`e_L`$ \| \| \| $`u_{1R}`$ \| $`d_{1R}`$ \| $`u_{1L}`$ \| $`d_{1L}`$ \| \| $`p=2`$ \| $`u_{2R}`$ \| $`d_{2R}`$ \| $`u_{2L}`$ \| $`d_{2L}`$ \| \| \| $`u_{3R}`$ \| $`d_{3R}`$ \| $`u_{3L}`$ \| $`d_{3L}`$ \| \| \| $`d_{1L}^c`$ \| $`u_{1L}^c`$ \| $`d_{1R}^c`$ \| $`u_{1R}^c`$ \| \| $`p=1`$ \| $`d_{2L}^c`$ \| $`u_{2L}^c`$ \| $`d_{2R}^c`$ \| $`u_{2R}^c`$ \| \| \| $`d_{3L}^c`$ \| $`u_{3L}^c`$ \| $`d_{3R}^c`$ \| $`u_{3R}^c`$ \| \| $`p=3`$ \| $`e_L^c`$ \| $`\nu _{eL}^c`$ \| $`e_R^c`$ \| $`\nu _{eR}^c`$ \| (12) Two more tables of this kind exist for the second and third generation. Some remarks are in order: * The symbols (including their sign) stand for the components $`\psi _I`$ of the field $`\psi `$ where the appropriate subset $`I\{1,2,3,4,5\}`$ is displayed in the table (11). For instance, $$\psi _{13}=u_{2R}\text{etc.}$$ By assumption, all fermions are massless to begin with. Therefore, each entry also represents a particle or antiparticle. The upper half of the table contains the particles (“matter”) while the lower half contains the antiparticles (“antimatter”). * Quarks such as $`u`$(up) and $`d`$(down) come in three colors: $`i=1,2,3`$. While quarks transform under the color group according to the representation $`\overline{3}`$, antiquarks transform according to the representation 3. Since both 3 and $`\overline{3}`$ are fundamental irreps of $`SU(3)`$, interchanging their role, as done here, has no physical effect. * Together with each spinor the charged conjugate spinor (with upper index $`^c`$) also enters the table (12) and so enters the field $`\psi `$. Since charge conjugation changes the chirality, we have to make precise what the symbols in (12) really mean. With the $`d`$-quark as an example our convention is: $$d_L^c:=(d^c)_L=(d_R)^c,d_R^c:=(d^c)_R=(d_L)^c.$$ * Contrary to the traditional formulation of the Minimal Standard Model, there is room for a right-handed neutrino. Note the presence of the unconventional field $`\nu _{eR}`$ (together with its charge conjugate $`\nu _{eL}^c`$) in the table (12) and the fact that it transform trivially under the gauge group $`G`$. Thus, the right-handed neutrino does not couple to any gauge field whatsoever which makes it hard to detect it in experiments. It only couples to the Higgs field and so acquires a mass after symmetry breaking. * Algebraic reasoning has led us to include certain minus signs in the table (12). One of the reasons is that we want the following condition to be satified: $$\sigma _I\psi _I^c=\psi _{I^c}$$ (13) where $`\sigma _I`$ is the sign of the permutation taking $`\{I,I^c\}`$ to its normal order $`\{1,2,3,4,5\}`$. This in particular guarantees that, if $`(\nu _e,e)_L`$ transforms as a $`SU(2)`$ doublet, so does the pair $$(e^c,\nu _e^c)_R=(e_L,\nu _{eL})^c$$ after charge conjugation. ## 5 The Concept of a Generalized Majorana Field The field $`\psi `$ is thought of as some generalized Dirac field having sufficiently many components $`\psi _I`$ to as to be able to describe all fundamental fermions of one generation. For its mathematical construction we need to introduce the spinor space $`S`$, basic to any Dirac field. Since its structure depends merely on the choice of spacetime (of even dimension in any case), we call $`S`$ the outer space. In the language of , $`S`$ is a Clifford supermodule, i.e., a linear space on which the $`\gamma `$ matrices act, carrying a $`_2`$-grading $$S=S^+S^{}$$ (14) given by the chirality, the eigenvalues $`\pm 1`$ of $`\gamma _5`$. Field components taking values in $`S^+`$ ($`S^{}`$) are said to be right-handed (left-handed). In addition, we have assumed that there is another superspace, $`^5`$, graded by the parity $`\kappa `$ of exterior powers. This space, specific to the Standard Model, is the same for all fermion generations. Since we wish to relate the chirality in $`S`$ to the parity in $`^5`$, the field $`\psi `$ is required to take values in the even part of the tensor product $`E=^5S`$ which is $$E^+=(^5S)^+=(^+^5S^+)(^{}^5S^{}).$$ (15) In order to write $`\psi `$ in terms of its components $`\psi _I`$ we use the basis $`e_I`$ for the inner space as constructed in the previous section. With respect to the tensor product $`^5S`$, we decompose the field $`\psi `$ as $$\psi (x)=\underset{I}{}e_I\psi _I(x)E^+,\psi _I(x)S.$$ (16) The condition $`\psi E^+`$ translates into $$\gamma _5\psi _I=(1)^{\|I\|}\psi _I.$$ (17) Hence, $`\psi _I`$ is right(left)-handed depending on whether $`\|I\|`$ is even(odd). Two operations of similar nature, one in $`^5`$ and one in $`S`$, will play an important role: * The Hodge operator $`:^5^5`$ is antilinear and acts on the basis as $$e_I=\sigma _Ie_{I^c}.$$ (18) Recall that $`\sigma _I`$ is the sign of the permutation taking $`\{I,I^c\}`$ to its normal order $`\{1,2,3,4,5\}`$. Since $`\|I\|+\|I^c\|=`$odd, the Hodge operator is parity changing: $$\kappa =\kappa .$$ (19) Note that $`\sigma _{I^c}=\sigma _I`$ (valid in odd dimensions) making the Hodge $``$ an involutive operator: $`^2=1\mathrm{l}`$ * The charge conjugation $`SS`$, $`ss^c`$, is antilinear, involutive and reverses the chirality: $`S^\pm S^{}`$. The identification $$\text{parity in }^5=\text{chirality in }S$$ suggests to couple both operations, resulting in a single antilinear operator $`:E^+E^+`$ such that $$\psi (x)=\underset{I}{}e_I\psi _I(x)\psi (x)=\underset{I}{}e_I\psi _I^c(x).$$ (20) Though the $``$ operator now changes matter into antimatter and vice versa, it should not be confused with “charge conjugation” in the traditional sense, $$\psi ^c(x)=\underset{I}{}e_I\psi _I^c(x),$$ which is not a symmetry. Using (13) and (18) we find $`\psi (x)`$ $`=`$ $`{\displaystyle \underset{I}{}}\sigma _Ie_{I^c}\psi _I^c(x)={\displaystyle \underset{I}{}}e_{I^c}\sigma _I\psi _I^c(x)`$ $`=`$ $`{\displaystyle \underset{I}{}}e_{I^c}\psi _{I^c}(x)={\displaystyle \underset{I}{}}e_I\psi _I(x)=\psi (x)`$ A generalized Dirac field $`\psi `$ based on the $`G`$-module $`^5`$ is said to be selfdual<sup>1</sup><sup>1</sup>1The definition of selfduality works for spaces $`^n`$ where $`n`$ is odd. If $`n=1`$, the concept reduces to that of an ordinary Majorana field. or a generalized Majorana field if it satisfies the relation $`\psi =\psi `$. Experimentally, three generations of fundamental fermions have been found: from the decays of the $`Z^0`$ boson one infers that there are exactly three generations (i.e., three is the number of neutrinos with masses below 45 GeV). It is a trivial matter to combine the Dirac fields $`\psi _f`$ ($`f=1,2,3`$) of three generations to a single master field: $$\mathrm{\Psi }=\psi _1ϵ_1+\psi _2ϵ_2+\psi _3ϵ_3$$ (21) where $`(ϵ_i)_{i=1}^3`$ denotes the standard basis in $`^3`$, the flavor space. For consistency, the $``$ operator must also act on $`^3`$ in an antilinear manner. We let it coincide with complex conjugation. Hence, if $`A`$ is some matrix in $`\text{End}^3`$, then $`A`$ stands for the complex conjugate matrix. The fundamental relation $`\psi _f=\psi _f`$, valid in each generation, is now equivalent to stating that $`\mathrm{\Psi }=\mathrm{\Psi }`$. ## 6 Analysis of Operators The complex space $`^5`$ we started from not only provides a natural basis $`(e_i)_{i=1}^5`$ but is also equipped with a Hermitian structure given by the standard scalar product. Supposing $`V`$ is any Hermitian vector space, one associates a multiplication operator $`ϵ(v)`$ to any $`vV`$, acting on the exterior algebra, $$ϵ(v)a=va(aV),$$ and lets $`\iota (v)`$ be its adjoint with respect to the induced Hermitian structure in $`V`$. In the language of Fock spaces, these operators are said to be creation and annihilation operators respectively. It is evident that they are of odd type, i.e., they change the parity of elements in $`V`$. We state this property as $$ϵ(v),\iota (v)\text{End}^{}V.$$ With $`V=^5`$ we put $`b_i=ϵ(e_i)`$ and $`b_i^{}=\iota (e_i)`$ so as to obtain the usual anticommutation relations: $$\{b_i,b_k^{}\}=\delta _{ik},\{b_i,b_k\}=0,\{b_i^{},b_k^{}\}=0.$$ The key observation is that lifting matrices $`a=(a_{ik})\text{End}^5`$ to observables now admits an explicit description: $$\theta (a)=\underset{ik}{}a_{ik}b_ib_k^{}\text{End}^+^5.$$ (22) To put it more formally, $`\theta (a)`$ is a derivation of the algebra $`^5`$. As for us, $`\theta (a)`$ is simply a linear operator of even type. Being a $`G`$-module, $`^5`$ carries a representation $`\theta `$ of Lie$`G`$. Note that the map $`\theta `$ constructed above extends the representation of the Lie algebra. Some of the previously introduced observables receive a new description: $`Q`$ $`=`$ $`\frac{1}{3}(b_1b_1^{}+b_2b_2^{}+b_3b_3^{})b_5b_5^{}`$ $`Y`$ $`=`$ $`\frac{2}{3}(b_1b_1^{}+b_2b_2^{}+b_3b_3^{})b_4b_4^{}b_5b_5^{}`$ $`I_3`$ $`=`$ $`\frac{1}{2}(b_4b_4^{}b_5b_5^{})`$ (23) $`p`$ $`=`$ $`b_1b_1^{}+b_2b_2^{}+b_3b_3^{}`$ $`q`$ $`=`$ $`b_4b_4^{}+b_5b_5^{}.`$ The significance of the operators $`p`$ and $`q`$ is that they generate the maximal algebra of operators, invariant under gauge transformations. This algebra, also referred to as the commutant $`(G)^{}`$, is abelian and 12-dimensional: there are 4 and 3 possible values for $`p`$ resp. $`q`$. Global gauge transformations act on the operators $`b_i`$ as they would act on the basis $`e_i`$: $$gb_ig^1=\underset{j}{}g_{ji}b_j(gG).$$ Consequently, $$gb_i^{}g^1=\underset{j}{}(g^{})_{ij}b_j^{}.$$ The Hodge operator $``$ has already been seen to play a prominent role. We shall now elaborate on its properties a little further. Recall first that $`^2=1\mathrm{l}`$. Second, we have $$p=3p,q=2q.$$ Third, the Hodge operator interchanges $`b_i`$ and $`b_i^{}`$ apart from possible sign change: $$b_i=\kappa b_i^{},b_i^{}=b_i\kappa .$$ See (9) for the definition of the parity operator $`\kappa `$. As a consequence, we get $$b_ib_j^{}=\kappa b_i^{}b_j\kappa =b_i^{}b_j=\delta _{ij}b_jb_i^{}.$$ Let $`a\text{End}^5`$ be arbitrary. From (22) and the antilinearity of $``$, $`\theta (a)`$ $`=`$ $`{\displaystyle \underset{ij}{}}\overline{a}_{ij}b_ib_j^{}`$ $`=`$ $`{\displaystyle \underset{ij}{}}(a^{})_{ji}(\delta _{ij}b_jb_i^{})=\text{tr}a^{}\theta (a^{}).`$ In particular, $$\theta (a)=\theta (a),a\mathrm{𝐬𝐮}(\mathrm{𝟓})$$ (24) owing to the relations $`a^{}=a`$ and $`\text{tr}a=0`$, and ultimately: $$g=g,g=e^aG.$$ (25) This is to demonstrate that we cannot dispense with the trace condition $`\text{tr}a=0`$, i.e., demanding that $`G`$ be a subgroup of $`SU(5)`$ (rather than of $`U(5)`$) if we want the consistency condition (25) to be satisfied. While elements of the Lie algebra are unchanged under the $``$ operation, their Hermitian counterparts $`i\theta (a)`$ pick up a minus sign owing to antilinearity, and so do the charges: $$Q=Q,Y=Y,I_3=I_3.$$ This is in accord with the conception, formulated before, that the Hodge $``$ converts particles into antiparticles and vice versa. ## 7 The Higgs Field There is no other way than to assume that particles receive their masses through the Higgs mechanism. As explained in and , the mechanism works well only if the Higgs field is an odd operator on the internal $`_2`$-graded space. In the present situation, this space is $$V=^5^3$$ endowed with the obvious grading $`V^\pm =^\pm ^5^3`$. With one Higgs doublet<sup>2</sup><sup>2</sup>2Normally, the Higgs field is thought of as a $`(Y=1)`$ doublet $`(\varphi _+,\varphi _0)`$. We prefer to work instead with $`\varphi _1=\varphi _0^{}`$ and $`\varphi _2=\varphi _+^{}`$. Note that $`\varphi _1`$ has zero electric charge while $`\varphi _2`$ has the charge $`Q=1`$. Note also that the scalar $`\varphi ^{}\varphi =\|\varphi _1\|^2+\|\varphi _2\|^2`$ is gauge invariant., $$\varphi =\left(\begin{array}{c}\varphi _1\\ \varphi _2\end{array}\right),$$ of hypercharge $`Y=1`$, there is more than one choice for the Higgs field $`\mathrm{\Phi }`$ (when written in terms of $`\varphi _1`$ and $`\varphi _2`$) if we want $`\mathrm{\Phi }`$ to act as an operator on the (rather large) internal space $`V`$, merely requiring that $`\varphi `$ transforms properly under the gauge group. As perceived by the founders of the Standard Model, this freedom of choice shows up in the appearance of a variety of undetermined Yukawa coupling constants, one for each elementary fermion field. Traditional thinking forbids to combine the elementary fermion fields into a single mathematical entity, as the constituents couple differently to the Higgs doublet $`\varphi `$. In essence, what we suggest here is to look at the freedom of fixing parameters (such as Yukawa couplings) from a different perspective. Recalling that the commutant $`(G)^{}\text{End}^5`$ is nontrivial, we are free to choose some invariant operator $$h(G)^{}\text{End}^3$$ (26) and to define the Higgs field by $$\mathrm{\Phi }(x)=h^{}\left(\varphi _1(x)b_4+\varphi _2(x)b_5\right)$$ (27) so as to have some field acting on the inner space $`^5^3`$. A gauge transformation, taking $`\mathrm{\Phi }`$ into $`g\mathrm{\Phi }g^1`$, reveals that the complex scalar fields $`\varphi _i`$ transform as desired. In general, the invariant operator (26) is determined by providing 12 matrices $$h(p,q)\text{End}^3(p=0,\mathrm{},3,q=0,1,2)$$ (28) The fact the matrices $`h(p,0)`$ do not contribute to the Higgs field (27) reduces the number 12 to 8. A further reduction to 4 is achieved with help of the constraint $$\mathrm{\Phi }(x)=\mathrm{\Phi }(x)^{}$$ (29) implying that particles couple to $`\mathrm{\Phi }`$ in the same way as antiparticles couple to $`\mathrm{\Phi }^{}`$. The condition (29) is equivalent to $$h=h^{}$$ (30) where another invariant operator $`h^{}`$, associated with $`h`$, has been introduced so as to satisfy the equation $`b_ih^{}=h^{}b_i\kappa `$ ($`i=4,5`$). Written out, the definition is $$h^{}(p,q)=(1)^{p+q}h(p,q+1)^{}(p=0,\mathrm{},3,q=0,1).$$ (31) Recall that $`h(p,q)`$ means complex conjugation of matrix elements in conjunction with the replacements $`p3p`$ and $`q2q`$. Therefore, another way to write the condition (30) is $$h(3p,2q)=(1)^{p+q}h(p,q+1)^T(p=0,\mathrm{},3,q=0,1)$$ (32) where $`A^T`$ denotes the transpose of a matrix $`A\text{End}^3`$. Therefore, among the matrices $`h(p,q)`$ $`(q>0)`$, only four have to be fixed in order to determine them all. Summarizing: > We are working with a common Yukawa coupling set equal to unity and assemble all free parameters, such as the entries of the fermion mass matrix, in the operator $`h`$. Presently, there does not seem to exist a convincing theoretical argument that would settle the question as to the origin of $`h`$ and completely determine the matrices $`h(p,q)`$ entering the mass operator. We need first of all understand the mathematical origin of the flavor symmetry (when there is no Higgs condensate). Given the form of $`h`$, the next step is to allege that the Higgs field $`\mathrm{\Phi }`$ enters the Dirac operator in a symmetrized form $$L=i(\mathrm{\Phi }+\mathrm{\Phi }^{})\text{End}^{}V$$ (33) so as to satisfy $$L^{}=L,L=L.$$ (34) Spontaneous symmetry breaking gives rise to a condensate, $$\mathrm{\Phi }_c=m^{}b_4,m=r^{1/2}h,r=\|\varphi _1\|^2$$ (35) and hence to a fermion mass operator $$M=iL_c=m^{}b_4+b_4^{}m,M=M=M^{},m=m^{}.$$ (36) Note the difference: while the symmetry breaking parameter $`r`$ is but an ordinary constant (to be obtained from the Higgs potential), the condensate $`\mathrm{\Phi }_c`$ and the mass operator $`M`$ are still operators of odd type. In terms of the mass matrices $`m(p,q)\text{End}^3`$ ($`q>0`$) the relation $`m=m^{}`$ may be written: $$m(3p,2q)=(1)^{p+q}m(p,q+1)^T(p=0,\mathrm{},3,q=0,1)$$ (37) The spectrum of $`M^2`$ describes the fermion masses (squared). The components $`m(p,q)`$ of $`M`$ refer to subgroups of particles, each group containing three particles of the same electric charge $`Q`$ but different flavors: \| $`m(0,1):`$ \| $`Q=0`$ \| $`m(0,2):`$ \| $`Q=1`$ \| \| --- \| --- \| --- \| --- \| \| $`m(2,1):`$ \| $`Q=3/2`$ \| $`m(2,2):`$ \| $`Q=1/3`$ \| Here, we we have listed only those matrices $`m(p,q)`$ for which $`p`$ is even as the relation (37) says that matrices, for which $`p`$ is odd, are related to the former. At the end, the relation simply guarantees that matter fields ($`p=`$even) and antimatter fields ($`p=`$odd) receive equal masses. ## 8 Superconnections, Generalized Dirac Operators, <br>and the Fermionic Action Recall the interpretation<sup>3</sup><sup>3</sup>3For the present purpose we avoid a global bundel-theoretic formulation and shall be content with formulas that hold locally. This is possible since, locally, all bundels are trivial. of the gauge field $`A`$ as a connection 1-form, i.e., Euclidean spacetime is modelled by some four-dimensional Riemannian manifold $`M`$, and $`A`$ takes values in $`T^{}M\text{Lie}G`$ where $`T^{}M`$ stands for the cotangent bundle. Our main interest lies in the lifted field $`\widehat{A}=\theta (A)`$ taking values in $`T^{}M\text{End}^+^5`$. In local coordinates, $`A`$ $`=`$ $`dx^\mu A_\mu ,A_\mu (x)\text{Lie}G`$ $`\widehat{A}`$ $`=`$ $`dx^\mu \widehat{A}_\mu ,\widehat{A}_\mu (x)\text{End}^+^5.`$ A superconnection (see for details) extends the notion of a gauge connection and is given by some first-order differential operator of odd type, $$ID=D+L,D=d+\widehat{A},$$ (38) acting on sections<sup>4</sup><sup>4</sup>4Sections are referred to as $`^5`$-valued differential forms. of the bundle $`T^{}M^5`$, where $`d`$ denotes the exterior derivative, $`D`$ the covariant derivative, and $`L`$ the Higgs field. Generally speaking, $`L`$ could also include $`n`$-forms ($`n1`$) complying with the oddness<sup>5</sup><sup>5</sup>5The property of being odd or even is defined with reference to the total $`_2`$-grading of the space $`T^{}M^5`$ on which these operators act. of $`ID`$. With an $`2n`$-dimensional Riemannian spin<sup>c</sup> manifold one associates a Clifford bundle $`C(M)`$ and a spin bundle which locally coincides with $`M\times S`$, $`S`$ being the spinor space, a complex vector space of dimension $`2^n`$. At $`xM`$, there is an isomorphism $`c:C(T_x^{}M)\text{End}S`$ and hence a way to construct $`\gamma `$ matrices, $$\gamma ^\mu =c(dx^\mu ),\{\gamma ^\mu ,\gamma ^\nu \}=2g^{\mu \nu },$$ involving $`g^{\mu \nu }`$, the metric tensor. Spinor fields $`\mathrm{\Psi }(x)`$ take values in the space $$F=VS=^5^3S.$$ To use the language of , $`F`$ is a twisted Clifford supermodule with $`V`$ the twisting space. With a superconnection $`ID`$ one associates the generalized Dirac operator $`ID/`$ which, roughly speaking, is obtained from $`ID`$ by replacing the basis elements $`dx^\mu `$ by $`\gamma ^\mu `$ wherever they occur. Therefore, if $`L`$ is scalar, $$ID/=/+\widehat{A}/+L=\gamma ^\mu (_\mu +\widehat{A}_\mu )+L.$$ (39) The Dirac operator acts on spinor fields such that $$\mathrm{\Psi }(x)F^\pm ID/\mathrm{\Psi }(x)F^{}.$$ Suppose the manifold $`M`$ is orientable and $`\omega _0`$ is a volume form. Then the fermionic action, considered as a functional of the master field $`\mathrm{\Psi }`$, is<sup>6</sup><sup>6</sup>6To put a factor $`\frac{1}{2}`$ in front is necessary because each elementary fermion field $`\psi _I`$ enters the action together with its charge conjugate $`\psi _I^c`$, both giving equal contributions. $$S_F=\frac{1}{2}_M\overline{\mathrm{\Psi }}iID/\mathrm{\Psi }\omega _0.$$ (40) By construction, the integrand incorporates both gauge and Yukawa interactions. Following we regard $`\mathrm{\Psi }\overline{\mathrm{\Psi }}`$ as an antilinear map into the dual space which reverses the chirality: $`\overline{(\mathrm{\Psi }_L)}=(\overline{\mathrm{\Psi }})_R`$. We also emphasize that, with regard to functional integration, the (Grassmann) variables $`\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}`$ are dependent, and a reasonable way to write the functional measure is $$d\mathrm{\Psi }=\underset{f,I}{}d\psi _{fI}$$ (41) with $`f=1,2,3`$ and $`I`$ running over the subsets of $`\{1,2,3,4,5\}`$. Two important features characterize the ansatz (38) for the Dirac operator. 1. The operator $`p`$ (but not $`q`$) commutes with $`ID/`$: $$pID/=ID/p.$$ (42) One consequence is that the baryon number and the lepton number are preserved in interactions: leptons and quarks do not couple at vertices of a Feynman diagram. Another is that matter is not converted into antimatter. Though matter and antimatter may annihilate to yield gauge and Higgs particles. 2. The $``$ operator anticommutes with $`ID/`$ and hence commutes with $`iID/`$: $$(iID/)=iID/.$$ (43) This fact follows from<sup>7</sup><sup>7</sup>7The relation $`\gamma ^\mu =\gamma ^\mu `$ is our way of stating that $`(\gamma ^\mu \psi )^c=\gamma ^\mu \psi ^c`$ in ordinary Dirac theory. $`\gamma ^\mu =\gamma ^\mu `$, (24) and (34). The property (43) of the Dirac operator implies that, if $`\mathrm{\Psi }`$ is a generalized Majorana field, so is $`\mathrm{\Psi }^{}=iID/\mathrm{\Psi }`$: $$\mathrm{\Psi }=\mathrm{\Psi }\mathrm{\Psi }^{}=\mathrm{\Psi }^{},$$ and from $$\overline{\mathrm{\Psi }}\mathrm{\Psi }^{}=\underset{I}{}\overline{\psi _I^c}\psi _I^{}{}_{}{}^{c}=\underset{I}{}\overline{\psi _I^{}}\psi _I=\overline{\mathrm{\Psi }^{}}\mathrm{\Psi },$$ we obtain $$\overline{\mathrm{\Psi }}iID/\mathrm{\Psi }=\overline{iID/\mathrm{\Psi }}\mathrm{\Psi }.$$ (44) The two properties of the Dirac operator stated above suggest decomposing $`\mathrm{\Psi }`$ into (anti)matter fields by projecting onto the $`p=`$even(odd) parts: $`\psi _M+\psi _A`$ $`=`$ $`\mathrm{\Psi }`$ $`\psi _M\psi _A`$ $`=`$ $`(1)^p\mathrm{\Psi }.`$ The fact that the $``$ operator switches between matter and antimatter is now reflected by the relation $`\psi _M=\psi _A`$. The antimatter field is thus seen to be a redundant variable and, if desired, may be eliminated from the action functional using $`\overline{\mathrm{\Psi }}iID/\mathrm{\Psi }`$ $`=`$ $`\overline{\psi _M}iID/\psi _M+\overline{\psi _A}iID/\psi _A`$ $`=`$ $`\overline{\psi _M}iID/\psi _M+\overline{\psi _M}iID/\psi _M`$ $`=`$ $`\overline{\psi _M}iID/\psi _M+\overline{\psi _M}(iID/\psi _M)=\overline{\psi _M}iID/\psi _M+\overline{iID/\psi _M}\psi _M.`$ The functional measure (41) now assumes the standard form $`d\psi _Md\overline{\psi }_M`$. The matter field may be decomposed even further so as to extract singlet and doublet components: $`\psi _{MS}+\psi _{MD}`$ $`=`$ $`\psi _M`$ $`\psi _{MS}\psi _{AD}`$ $`=`$ $`(1)^q\psi _M.`$ Owing to the term $`L`$, part of the Dirac operator and anticommuting with $`q`$, the Dirac operator induces transitions $`SD`$ and $`DS`$. Also, since $$(1)^{p+q}\psi _M=(1)^q\psi _M,$$ the field $`\psi _{MS}`$ is right-handed while $`\psi _{MD}`$ is left-handed. ## 9 Currents It is instructive to see how currents emerge from the ansatz (40). For this we need only evaluate $`\widehat{A}/=\theta (A/)`$ in some basis. Let $`(it_a)_{a=1}^{12}`$ be any basis in $`\text{Lie}G`$ (so that $`t_a^{}=t_a`$). If the gauge field $`A`$ has real components $`A_\mu ^a`$ given by $`iA_\mu =A_\mu ^at_a`$, then $$i\widehat{A}/=A_\mu ^aT_a\gamma ^\mu ,T_a=\theta (t_a).$$ The representation of gauge couplings in terms of currents is an immediate result: $$\overline{\psi }_Mi\widehat{A}/\psi _M=j_a^\mu A_\mu ^a,j_a^\mu =\overline{\psi }_MT_a\gamma ^\mu \psi _M.$$ (45) Inspection shows that all currents preserve both $`p`$ and $`q`$ in the following sense: $$pT_a=T_ap,qT_a=T_aq$$ Each current can therefore be decomposed into constituents with a definite $`(p,q)`$ assignment, $$j_a^\mu =\underset{p,q}{}j_a^\mu (p,q),$$ meaning that fermions with different $`(p,q)`$ assignments have no common vertex. There is no current for the right-handed neutrino: $`j_a^\mu (0,0)=0`$. Currents may be characterized according to their behavior with respect to chirality. Definition. A current $`j_a^\mu =\overline{\psi }_MT_a\gamma ^\mu \psi _M`$ (for fixed index $`a`$) is said to be vectorlike if $`T_ab_4=b_4T_a`$ and chiral otherwise. Likewise, the interaction $`j_a^\mu A_\mu ^a`$ is said to be vectorlike resp. chiral if the current is. The reason for the above definition is the observation that the operators $`b_4`$ and $`b_4^{}`$ switch between left- and right-handed fields of the same kind. No matter what gauge group $`G`$, vectorlike currents are abundant. They form a linear subspace of the space of all currents. While the total space is 12-dimensional, the subspace is 9-dimensional for our choice of $`G`$. It is spanned by the electric current and eight currents pertaining to the color group. In other words, the interactions with the photon and the gluons are vectorlike. There is no unique choice of generators $`T_a`$ for the 3-dimensional quotient space. One choice could be: $`I_1`$ $`=`$ $`\frac{1}{2}(b_5b_4^{}+b_4b_5^{})`$ $`I_2`$ $`=`$ $`\frac{i}{2}(b_5b_4^{}b_4b_5^{})`$ (46) $`I_3`$ $`=`$ $`\frac{1}{2}(b_4b_4^{}b_5b_5^{})`$ It is indeed easily verified that $`I_ab_4b_4I_a`$ ($`a=1,2,3`$). The observation made here corresponds to the fact that nature provides three basic parity violating interactions of matter corresponding to the exchange of $`W^\pm `$ and $`Z^0`$ bosons. The coupling to the massive vector bosons specify the selection of $`T`$’s not commuting with $`b_4`$. In passing we mention that, with the group $`SU(5)`$ replacing $`G`$, the quotient space of chiral currents would be 9-dimensional, admitting further parity violating interactions. ## 10 The Bosonic Action In the same way as the Yang-Mills connection<sup>8</sup><sup>8</sup>8We use the terms connection and covariant derivative interchangeably. $`D`$ gives rise to the concept of curvature, $$F=D^2=\frac{1}{2}dx^\mu dx^\nu F_{\mu \nu }^a(iT_a),$$ the superconnection $`ID`$ gives rise to a generalized curvature , $$IF=ID^2=F+DL+L^2,$$ (47) where we have identified $`DL`$, the covariant derivative of the Higgs field, with the supercommutator<sup>9</sup><sup>9</sup>9Note that both $`D`$ and $`L`$ are odd operators. $`[[D,L]]=\{D,L\}`$. Without the Higgs field the bosonic action would consist of nothing but the Yang-Mills term: $$S_B=\frac{1}{2}F^2=\frac{1}{2}_M\|F\|^2\omega _0(L=0).$$ The precise nature of the invariant $`\|F\|^2`$ better be such that it reduces to $$\|F\|^2=\frac{1}{2}\underset{a,\mu \nu }{}(F_{\mu \nu }^a)^2$$ in a flat Euclidean universe. Constructing an invariant $`\|F\|^2`$ for 2-forms $`F`$ with the required property is a well-known procedure. An extension to $`IF`$, properly treating $`p`$-forms of arbitrary order, has been given and discussed in . It requires various steps. 1. Step. A basis $`it_a`$ in $`\text{Lie}G`$ is chosen such that $$\text{Tr}T_aT_b=\delta _{ab}(a,b=1,\mathrm{},12)$$ (48) where $`T_a=\theta (t_a)`$ as before. 2. Step. Since any operator $`A`$ in $`^5`$ (like $`T_a`$) extends to an operator in $$V=^5^3,$$ we shall not distinguish between $`A`$ and its extended version, $`A1\mathrm{l}`$. However when it comes to studying their traces, there would be a distinction unless ‘Tr’ in $`^3`$ is given another meaning: $$\text{Tr}C=\frac{1}{3}\underset{i=1}{\overset{3}{}}C_{ii},C=(C_{ik})\text{End}^3.$$ The modified trace is but an average. We stipulate that traces, although written $`\text{Tr}A`$, always include an extra factor $`\frac{1}{3}`$ for operators $`A`$ in $`V`$. This precaution is necessary in order to preserve the validity of relations like (48) or $`\text{Tr}q=32`$ and many others. 3. Step. There is no need for a differential form to be homogeneous (i.e., a $`p`$-form). As the construction of $`IF`$ shows, we are dealing here with very general forms $`B`$ that are sections of the algebra $$=T^{}M\widehat{}\text{End}V.$$ Since both $`T^{}M`$ and $`\text{End}V`$ are superalgebras, so is their product, and the symbol $`\widehat{}`$ accounts for that fact: the tensor product is special for $`_2`$-graded algebras . Any element $`B`$ taking values in $``$ has the following structure<sup>10</sup><sup>10</sup>10Again, we will be content here with a local description.: $$B=\underset{I}{}dx^IB_I,B_I\text{End}V.$$ As $`\text{dim}M=4`$, the sum runs over the subsets $`I\{1,2,3,4\}`$ and $$dx^I=dx^{\mu _1}\mathrm{}dx^{\mu _r}$$ where $$I=\{\mu _1,\mathrm{},\mu _r\},\mu _1<\mathrm{}<\mu _r,r=\|I\|.$$ We let $$\|B\|^2=\underset{I,J}{}g^{IJ}\text{Tr}B_I^{}B_J$$ (49) where $`g^{IJ}=(dx^I,dx^J)`$, to be obtained from the Riemannian metric, and $`g^{IJ}=\delta ^{IJ}`$ in a flat Euclidean universe. Similar to the procedure in we let the gauge-invariant bosonic action be given by $$S_B=_M\frac{1}{2}\|IF+\mu ^2C\|^2\omega _0$$ (50) with $`\mu `$ some mass parameter and $$C(G)^{}\text{End}^3,C=C=C^{}.$$ (51) Below we shall learn that it suffices to take $`C=1\mathrm{l}`$. Without the shift operator $`\mu ^2C`$ the minimum of $`S_B`$ would be attained for a flat superconnection: $`IF=0`$. By contrast, the ansatz (50) when taken as classical field theory predicts a constant curvature $`IF`$ in the ground state. Thanks to the fact that $`IF`$ splits into $`p`$-forms with $`p=0,1,2`$, the integrand above consists of three terms only, $$\frac{1}{2}\|IF+\mu ^2C\|^2=\frac{1}{2}\|F\|^2+\frac{1}{2}\|DL\|^2+\frac{1}{2}\|L^2+\mu ^2C\|^2$$ which are easily identified as the Yang-Mills term, the covariant kinetic term of the Higgs field, and the Higgs potential. In order to guarantee the correct behavior of the kinetic term, i.e., $$\frac{1}{2}\|dL\|^2=g^{\mu \nu }(_\mu \varphi )^{}(_\nu \varphi ),$$ a normalization condition must be satisfied: $$\frac{1}{2}\text{Tr}qhh^{}=1.$$ (52) As we shall see (in Section 13), this condition has important consequences. ## 11 The Higgs Potential and Symmetry Breaking In this section, we shall work out the details of the Higgs potential $$V(\varphi )=\frac{1}{2}\|L^2+\mu ^2C\|^2=\frac{1}{2}\text{Tr}(L^2+\mu ^2C)^2.$$ Starting from (27) and (33) we first evaluate some traces using $`C=C`$, $`\mathrm{\Phi }=\mathrm{\Phi }^{}`$, and the $``$-invariance of the trace: $`\text{Tr}L^4`$ $`=`$ $`2\text{Tr}(\mathrm{\Phi }\mathrm{\Phi }^{})^2=(\varphi ^{}\varphi )^2\text{Tr}q(hh^{})^2`$ $`\text{Tr}L^2C`$ $`=`$ $`\text{Tr}(\mathrm{\Phi }\mathrm{\Phi }^{}+\mathrm{\Phi }^{}\mathrm{\Phi })C=\text{Tr}\mathrm{\Phi }\mathrm{\Phi }^{}(C+C)`$ $`=`$ $`2\text{Tr}\mathrm{\Phi }\mathrm{\Phi }^{}C=\varphi ^{}\varphi \text{Tr}qhh^{}C.`$ As a result of this calculation the Higgs potential may be written in terms of three constants: $$V(\varphi )=V_0+\frac{\lambda }{4}(\varphi ^{}\varphi r)^2.$$ (53) Apart from the irrelevant $`V_0`$, the other two constants are \| the Higgs coupling \| $`\lambda =2\text{Tr}q(hh^{})^2`$ \| \| --- \| --- \| \| the condensate \| $`r=\mu ^2\text{Tr}qhh^{}C/\text{Tr}q(hh^{})^2`$ . \| (54) The invariant operator $`C`$ is seen to enter the Higgs potential via two constants only, $`V_0`$ and $`r`$. There will be no loss of generality when we decide to work $`C=1\mathrm{l}`$ from now on. This fixes the parameter $`\mu `$ setting the mass scale, and so $$\lambda r=4\mu ^2$$ (55) owing to the normalization condition (52). The techniques presented here are, for the most part, quite standard. However, we have decided to work with $`r=v^2/2`$ where $`v`$ is, by convention, the symmetry breaking parameter used in most texts on the subject. Note that in the present framework the Higgs coupling constant $`\lambda `$ can never be negative or zero. It cannot be arbitrarily small either. In fact, $$\lambda \frac{1}{4}.$$ (56) The lower bound can be derived as follows. Define averages $$W=\text{Tr}qW/\text{Tr}q$$ for operators $`W`$ acting on $`^5^3`$. The obvious inequality $`W^2W^2`$ for $`W=hh^{}`$ together with $`\text{Tr}q=32`$ and (52) leads to the lower bound for $`\lambda `$. As $`r>0`$, we observe a breakdown of symmetry, and choosing the unitary gauge, we have $$\varphi _1=2^{1/2}\phi +r^{1/2},\varphi _2=0$$ where $`\phi `$ is the neutral scalar Higgs field. To extract its mass, we need only expand the Higgs potential up to second-order terms: $$V(\varphi )=V_0+\frac{1}{2}m_H^2\phi ^2+O(\phi ^3),m_H=2\mu .$$ (57) ## 12 Masses and Coupling Constants <br>of Vector Bosons To give the particulars of vector bosons we must specify the basis $`it_a`$ in $`\text{Lie}G`$ obeying (48): \| \| $`a=1,2,3`$ \| $`a=4`$ \| $`a=5,\mathrm{},12`$ \| \| --- \| --- \| --- \| --- \| \| $`t_a`$ \| $`\frac{1}{4}\left(\begin{array}{cc}0& 0\\ 0& \sigma _a\end{array}\right)`$ \| $`\frac{1}{4}\sqrt{\frac{3}{5}}\left(\begin{array}{cc}\frac{2}{3}1\mathrm{l}_3& 0\\ 0& 1\mathrm{l}_2\end{array}\right)`$ \| $`\frac{1}{4}\left(\begin{array}{cc}\lambda _{a4}& 0\\ 0& 0\end{array}\right)`$ \| We have written down the generators $`t_a`$ as matrices acting on $`^3^2`$ and made use of the Gell-Mann matrices $`\lambda _a`$, which operate on $`^3`$, and the Pauli matrices $`\sigma _a`$, which operate on $`^2`$. Checking the conditions (48) is facilitated by the relation $$\text{Tr}T_aT_b=8\text{tr}t_at_b.$$ With respect to the given basis, the mass matrix<sup>11</sup><sup>11</sup>11A matrix whose eigenvalues are the masses squared. $`m^2`$ of the 12 vector bosons has matrix elements $$m_{ab}^2=\text{Tr}[T_a,L_c][T_b,L_c]$$ (58) (see for a derivation) with $`L_c`$, the ‘condensate of $`L`$’, as in (36). By a straightforward computation using (52), we obtain $$m_{ab}^2=r\{t_a,t_b\}_{44},(a,b=1,\mathrm{},12).$$ (59) where $`\{t_a,t_b\}_{ik}`$ are the matrix elements of $`\{t_a,t_b\}`$ as an operator on $`^5`$. By inspection, $`m_{ab}^2=0`$ if either $`a5`$ or $`b5`$ (gluons do not get masses) and we are left with some nontrivial $`4\times 4`$ matrix $$(m^2)_{a,b=1}^4=\frac{r}{8}\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& \sqrt{3/5}\\ 0& 0& \sqrt{3/5}& 3/5\end{array}\right)$$ the eigenvalues of which provide the masses of the four vector bosons of the electroweak sector: \| mass squared of the $`W^\pm `$ \| : \| $`m_W^2=r/8`$ \| \| --- \| --- \| --- \| \| mass squared of the $`Z^0`$ \| : \| $`m_Z^2=r/5`$ \| \| mass squared of the $`\gamma `$ \| : \| 0 \| (60) The eigenvalue $`r/8`$, when identified with $`m_W^2`$, together with (55) and (57) leads to an expression for the Higgs coupling constant, $$\lambda =\frac{m_H^2}{8m_W^2},$$ and together with (56) to a remarkable inequality concerning the Higgs mass: $$m_H\sqrt{2}m_W.$$ (61) The eigenvectors corresponding to the eigenvalues $`r/5`$ and 0 are $`t_z`$ $`=`$ $`\mathrm{cos}\theta _Wt_3\mathrm{sin}\theta _Wt_4,\mathrm{cos}\theta _W=\sqrt{5/8}`$ $`t_0`$ $`=`$ $`\mathrm{sin}\theta _Wt_3+\mathrm{cos}\theta _Wt_4,\mathrm{sin}\theta _W=\sqrt{3/8}`$ with $`\theta _W`$ the Weinberg angle. The relation $`\mathrm{sin}^2\theta _W=3/8`$ is typical for a $`SU(5)`$-oriented gauge theory. The change of basis leads to the photon field $`A_\mu ^0(x)`$ and the field $`Z_\mu (x)`$ of the $`Z^0`$ particle given by $`Zt_z+A^0t_0=A^3t_3+A^4t_4`$. Hence, $`Z`$ $`=`$ $`\mathrm{cos}\theta _WA^3\mathrm{sin}\theta _WA^4`$ $`A^0`$ $`=`$ $`\mathrm{sin}\theta _WA^3+\mathrm{cos}\theta _WA^4`$ Let us now investigate the structure of the basis vector we associate with the photon: $`t_0`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{8}}}{\displaystyle \frac{1}{4}}\left(\begin{array}{cc}0& 0\\ 0& \sigma _3\end{array}\right)+\sqrt{{\displaystyle \frac{5}{8}}}{\displaystyle \frac{1}{4}}\sqrt{{\displaystyle \frac{3}{5}}}\left(\begin{array}{cc}\frac{2}{3}1\mathrm{l}_3& 0\\ 0& 1\mathrm{l}_2\end{array}\right)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\sqrt{{\displaystyle \frac{3}{2}}}\text{diag}(\frac{1}{3},\frac{1}{3},\frac{1}{3},0,1).`$ Physics takes place not in $`^5`$ but in $`^5`$. Letting $`T_0=\theta (t_0)`$, we find a relationship with the operator of electric charge, $$T_0=\frac{1}{4}\sqrt{\frac{3}{2}}Q.$$ (62) In the same manner we deduce, for $`T_z=\theta (t_z)`$, $$T_z=\sqrt{\frac{2}{5}}\left(I_3\frac{3}{8}Q\right).$$ Consult (23) and (46) for the structure of $`Q`$, $`Y`$, and $`I_i`$. Together with $`2\sqrt{2}T_+=b_4b_5^{}`$ and $`2\sqrt{2}T_{}=b_5b_4^{}`$ derived from $$T_\pm =\frac{1}{\sqrt{2}}(T_1\pm iT_2)=\frac{1}{2\sqrt{2}}(I_1\pm iI_2)$$ we have thus obtained expressions for all currents of the electroweak sector: \| $`\overline{\psi }_M\gamma ^\mu T_0\psi _M`$ \| electromagnetic current (vectorlike) \| \| --- \| --- \| \| $`\overline{\psi }_M\gamma ^\mu T_z\psi _M`$ \| neutral weak current (chiral) \| \| $`\overline{\psi }_M\gamma ^\mu T_\pm \psi _M`$ \| charged weak currents (chiral) \| Consider the fine structure constant $`\alpha =e^2/(4\pi )`$ and recall the formula (62). To give the photon the coupling constant $`e`$ is to say that $$e=\frac{1}{4}\sqrt{\frac{3}{2}}\text{or}\alpha ^1=128\pi /3=134.04\mathrm{}$$ (63) Renormalization group equations suggest that $`\alpha ^1`$ decreases slightly with energy. The value above is in accord with this conception: it lies half way between $`\alpha ^1(m_Z)=128`$ and $`\alpha ^1(0)=137`$. Sparked by the success, we proceed extracting the coupling constants $`g`$ and $`g^{}`$ of the Salam-Weinberg (SW) model from the relation<sup>12</sup><sup>12</sup>12Within the SW formalism, one writes $`B`$ instead of $`A^4`$ and calls $`\frac{1}{2}Y`$ the hypercharge. $$\underset{a=1}{\overset{4}{}}A^aT_a=\underset{i=1}{\overset{3}{}}gA^iI_i+\frac{1}{2}g^{}A^4Y$$ with the following result: $$g=\frac{1}{2},g^{}=\frac{1}{2}\sqrt{\frac{3}{5}}.$$ (64) The value $`1/2`$ obtained for $`g`$ is consistent with the SW formula $`m_W=gv/2`$ taking the relations $`r=v^2/2`$ and $`m_W^2=r/8`$ into account. The value for $`g^{}`$ is consistent with the SW formula $`g^{}/g=\mathrm{tan}\theta _W`$ provided $`\mathrm{tan}\theta _W=\sqrt{3/5}`$. ## 13 Fermion Masses and CKM Matrices Recall the structure of the fermion mass operator $`M=M^{}`$ from Section 7; $`M`$ stays invariant under gauge transformations $`g`$ provided $`g`$ is an element of the residual gauge group, leaving $`b_4`$ invariant. It consists of complex $`3\times 3`$ mass matrices $`m(p,q)`$ where $`q>0`$, each one pertaining to a group of three fermions having different flavors but same quantum numbers otherwise. The relation $`M=M`$, equivalently the relation (37), reflects the matter-antimatter symmetry of $`M`$. A diagonalization can be achieved in the following sense. There is a unitary operator $`U`$ on the Hermitian vector space $`V`$, i.e. an element $`USU(V)`$, such that * the transformed mass operator $$M^{}=UMU^{}=m_{}^{}{}_{}{}^{}b_4+b_4^{}m^{}$$ has diagonal mass matrices $`m^{}(p,q)`$ for all $`p`$ and $`q`$, * the operator $`U`$ commutes with gauge transformations from the residual gauge group and may be represented as $$U=U_0b_4^{}b_4+U_1b_4b_4^{}$$ with operators $`U_i(G)^{}\text{End}^3`$ where $`^3`$ is the flavor space, * the relation $`U=U`$ is satisfied, and hence $`U_1=U_0`$, * when replacing $`T`$ by $`T^{}=UTU^{}`$, we observe that currents such as $`\overline{\psi }_M\gamma ^\mu T\psi _M`$ remain unchanged for most generators $`T`$ except when $`T=T_\pm `$. The unitary transformation thus creates a flavor changing charged current associated with the transformed generator $$T_+^{}=UT_+U^{}=T_+U_M,U_M=U_1U_0^{}$$ where $`U_M`$ is known as the CKM matrix. We shall now comment on these items. Passage to the diagonal form of all mass matrices means that we construct the mass eigenstates: $`m^{}(0,1)`$ $`=`$ $`\text{diag}(m_{\nu _e},m_{\nu _\mu },m_{\nu _\tau })`$ $`m^{}(0,2)`$ $`=`$ $`\text{diag}(m_e,m_\mu ,m_\tau )`$ $`m^{}(2,1)`$ $`=`$ $`\text{diag}(m_u,m_c,m_t)`$ $`m^{}(2,2)`$ $`=`$ $`\text{diag}(m_d,m_s,m_b)`$ To keep this list short we confined attention to $`p=`$ even. The structure of $`U`$ follows at once from the splitting $$VV_0V_1,V_0V_1$$ corresponding to the eigenvalues 0 and 1 of the operator $`b_4b_4^{}`$. This allows us to write $$b_4\left(\begin{array}{cc}0& 0\\ 1\mathrm{l}& 0\end{array}\right),b_4^{}\left(\begin{array}{cc}0& 1\mathrm{l}\\ 0& 0\end{array}\right),U\left(\begin{array}{cc}U_0& 0\\ 0& U_1\end{array}\right),M\left(\begin{array}{cc}0& m\\ m^{}& 0\end{array}\right)$$ from which all results can be inferred at ease. The operators $`U_0`$ and $`U_1`$ have components $$U_0(p,q)=\{\begin{array}{cc}\text{unitary}\hfill & \text{if }q=0,1\hfill \\ 0\hfill & \text{if }q=2\hfill \end{array}U_1(p,q)=\{\begin{array}{cc}\text{unitary}\hfill & \text{if }q=1,2\hfill \\ 0\hfill & \text{if }q=0\hfill \end{array}$$ and the unitary transformation $`M^{}=UMU^{}`$ diagonalizes each individual mass matrix $`m(p,q)`$ by means of a biunitary transformation: $$m^{}(p,q)=U_0(p,q1)m(p,q)U_1^{}(p,q)(q>0).$$ Performing these transformations eliminates many irrelevant parameters from $`M`$ and hence from the input operator $`h`$. The analysis above also reveals the structure of the CKM matrix: 1. $`U_M`$ is thought of as an element of $`(G)^{}\text{End}^3`$, having components $$U_M(p,q)=U_0(p,q)U_1^{}(p,q)=\{\begin{array}{cc}\text{unitary}\hfill & q=1\hfill \\ 0\hfill & q=0,2\hfill \end{array}$$ 2. As the CKM matrix satisfies the relation $`U_M^{}=U_M`$, we have $$U_M(3p,1)=U_M(p,1)^T,p=0,\mathrm{},3.$$ 3. There are precisely two independent unitary $`3\times 3`$ matrices which completely determine $`U_M`$: $`U_M(0,1)`$ $`=`$ CKM matrix of the leptons $`U_M(2,1)`$ $`=`$ CKM matrix of the quarks They enter the charged currents and thus induce flavor changing interactions: $$j_+^\mu =\overline{\psi }_M\gamma ^\mu T_+\psi _M^{},j_{}^\mu =\overline{\psi }_M^{}\gamma ^\mu T_{}\psi _M,\psi _M^{}=U_M\psi _M.$$ There are no flavor changing neutral currents in such a theory. From the fact that traces are invariant under unitary transformations we obtain the following trace formulas $`\text{Tr}m(0,1)m(0,1)^{}`$ $`=`$ $`\frac{1}{3}(m_{\nu _e}^2+m_{\nu _\mu }^2+m_{\nu _\tau }^2)`$ $`\text{Tr}m(0,2)m(0,2)^{}`$ $`=`$ $`\frac{1}{3}(m_e^2+m_\mu ^2+m_\tau ^2)`$ $`\text{Tr}m(2,1)m(2,1)^{}`$ $`=`$ $`\frac{1}{3}(m_u^2+m_c^2+m_t^2)`$ $`\text{Tr}m(2,2)m(2,2)^{}`$ $`=`$ $`\frac{1}{3}(m_d^2+m_s^2+m_b^2).`$ Consequently, taking (37) into account, $`\text{Tr}M^2`$ $`=`$ $`\text{Tr}qmm^{}=_{p,q}q\left(\genfrac{}{}{0pt}{}{p}{3}\right)\left(\genfrac{}{}{0pt}{}{q}{2}\right)\text{Tr}m(p,q)m(p,q)^{}`$ (65) $`=`$ $`\frac{4}{3}(m_{\nu _e}^2+m_{\nu _\mu }^2+m_{\nu _\tau }^2+m_e^2+m_\mu ^2+m_\tau ^2)`$ $`+4(m_u^2+m_c^2+m_t^2+m_d^2+m_s^2+m_b^2).`$ In the same manner, one finds $`\text{Tr}M^4`$ $`=`$ $`\text{Tr}q(mm^{})^2=_{p,q}q\left(\genfrac{}{}{0pt}{}{p}{3}\right)\left(\genfrac{}{}{0pt}{}{q}{2}\right)\text{Tr}\left(m(p,q)m(p,q)^{}\right)^2`$ (66) $`=`$ $`\frac{4}{3}(m_{\nu _e}^4+m_{\nu _\mu }^4+m_{\nu _\tau }^4+m_e^4+m_\mu ^4+m_\tau ^4)`$ $`+4(m_u^4+m_c^4+m_t^4+m_d^4+m_s^4+m_b^4).`$ Two equations relate the $`W`$ mass and the Higgs mass to the fermion masses: $`16m_W^2`$ $`=`$ $`\text{Tr}M^2`$ (67) $`4m_H^2m_W^2`$ $`=`$ $`\text{Tr}M^4.`$ (68) The first equality uses (65), (35), (52), and (60): $$\text{Tr}M^2=\text{Tr}qmm^{}=r\text{Tr}qhh^{}=2r=16m_W^2,$$ while the second uses (66), (35), (54), (60), (55), and (57): $$\text{Tr}M^4=\text{Tr}q(mm^{})^2=r^2\text{Tr}q(hh^{})^2=\frac{1}{2}r^2\lambda =4\lambda rm_W^2=4m_H^2m_W^2.$$ Empirically , the mass of the top quark is dominant among the fermion masses, and so $$\text{Tr}M^24m_t^2,\text{Tr}M^44m_t^4.$$ The relations (67) and (68) therefore tell us that $$m_Hm_t2m_W.$$ Prior to its observation, the value of the top quark mass has been predicted on theoretical grounds within a 10% error bracket, certainly one of the greatest triumphs of the Standard Model. Again, the value $`2m_W`$ for the top mass obtained above is off the empirical value $`(174\pm 7)`$GeV by 10%. A value for the Higgs mass around $`2m_W`$ has already been predicted in . The fact that $`m_H`$ depends strongly on the top quark mass, suggesting $`m_H158`$GeV, has also been noted by Okumura who argues on the basis of noncommutative geometry. Similarly, Pirogov and Zenin , using the renormalization group approach, find that a cutoff equal to the Planck scale would give the Higgs boson a mass around $`160`$GeV. Surprisingly, the Higgs mass values offered in the literature, though dependent on very different schemes, all seem to to converge. References 1. S. Catani et.al.: The QCD and the Standard Model Working Group: Summary Report. hep-ph/0005114 M.W. Grunewald, Phys.Rept. 322 (1999), 125 W. Hollik, Acta Phys.Polon. B30 (1999), 1787 C. Dionisi, Nucl.Phys.Proc.Suppl. 38 (1995), 125 S.J. Brodsky: Precision Tests of QCD and the Standard Model, hep-ph/9506322 P.H. Chankowski: Precision Tests of the MSSM, hep-th/9505304 2. G. Roepstorff: Superconnections and the Higgs Field, hep-th/9801040 and J.Math.Phys. 40 (1999) 2698 3. G. Roepstorff: Superconnections and Matter, hep-th/9801045 4. G. Roepstorff: Superconnections: an Interpretation of the Standard Model, hep-th/9907221, Electronic Journal of Diff. Equ., Conf. 04, 2000, pp.165-174 5. G. Roepstorff and Ch. Vehns: An Introduction to Clifford Supermodules, math-ph/9908029 6. G. Roepstorff and Ch. Vehns: Generalized Dirac Operators and Superconnections, math-ph/9911006 7. G. Roepstorff: A Class of Anomaly-Free Gauge Theories, hep-th/0005079 8. CDF and DO Collaboration (G. Brooijmans for the collaboration): Top Quark Mass Measurements at the Tevatron, hep-ex/0005030 9. E. Laenen, J. Smith, and W. van Neerven, Phys.Lett. B321 (1994), 254 E.L. Berger, H. Contopanagos, Phys.Rev. D54 (1996), 3085 S. Catani, M.L. Mangano, P. Nason, L. Trentadue, Phys.Lett. B378 (1996), 329 10. Y. Okumura: An estimation of the Higgs boson mass in the two loop approximation in a noncommutative differential geometry, hep-ph/9707350 11. Yu.F. Pirogov and O.V. Zenin: Two-loop renormalization group profile of the standard model and a new generation, hep-ph/9808414
warning/0006/astro-ph0006383.html	ar5iv	text	# THEORY OF LOW-MASS STARS AND SUBSTELLAR OBJECTS ## 1 INTRODUCTION Interest in the physics of objects at the bottom of and below the main sequence (MS) originated in the early demonstration by Kumar (1963) that hydrogen-burning in a stellar core no longer occurs below a certain mass and that below this limit hydrostatic equilibrium against gravitational collapse is provided by electron degeneracy. Simple analytical arguments, based on the balance between the classical ionic thermal pressure and the quantum electronic pressure yield for this H-burning minimum mass (HBMM) $`m_{\mathrm{HBMM}}0.1M_{}`$, while the first detailed evolutionary calculations gave $`m_{\mathrm{HBMM}}0.085M_{}`$ (Grossman et al 1974). Tarter (1975) proposed the term ”brown dwarfs” (BD) for objects below this H-burning limit. D’Antona and Mazzitelli (1985) first pointed out that the luminosity below $`m_{\mathrm{HBMM}}`$ would stretch by a few orders of magnitude over a few hundredths of a solar mass, making the observation of BDs a tremendously difficult task. Subsequent benchmarks in low-mass star (LMS) and BD theory are due primarily to VandenBerg et al (1983), D’Antona and Mazzitelli (1985), Dorman et al (1989), Nelson et al (1986, 1993a) and the Tucson group (Lunine et al 1986, 1989, Burrows et al 1989, 1993). In spite of this substantial theoretical progress, all of these models failed to reproduce the observations at the bottom of the MS and thus could not provide a reliable determination of the characteristic properties (mass, age, effective temperature, luminosity) of low-mass stellar and substellar objects. A noticeable breakthrough came from the first calculations by a few groups of synthetic spectra and atmosphere models characteristic of cool ($`T_{\mathrm{eff}}<\mathrm{\hspace{0.17em}4000}`$ K) objects (Allard 1990, Brett & Plez 1993, Saumon et al 1994, Allard & Hauschildt 1995, Brett 1995, Tsuji et al 1996, Hauschildt et al 1999). This allowed the computation of consistent non-grey evolutionary models (Saumon et al 1994, Baraffe et al 1995) and the direct confrontation of theory and observation in photometric passbands and color-magnitude diagrams, thus avoiding dubious transformations of observations into theoretical L-$`T_{\mathrm{eff}}`$ Hertzsprung-Russell diagrams. In the meantime, the search for faint (sub)stellar objects has bloomed over the past few years. Several BDs have now been identified since the first discoveries of bona-fide BDs (Rebolo et al 1995, Oppenheimer et al 1995) either in young clusters (see Basri, this volume and Martín 1999 for reviews, and references therein) or in the Galactic field (Ruiz et al 1997, Delfosse et al 1997, Kirkpatrick et al 1999a). The steadily increasing number of identified extra-solar giant planets (EGPs) since the discovery of 51PegB (Mayor & Queloz 1995, see Marcy & Butler 1998 for a review) has opened up a new era in astronomy. Ongoing and future ground-based and space-based optical and infrared surveys of unprecedented faintness and precision are likely to reveal hundreds more of red dwarfs, brown dwarfs and giant planets which will necessitate the best possible theoretical foundation. The correct understanding of the physical properties of these objects bears major consequences for a wide range of domains of physics and astrophysics: dense matter physics, planet and star formation and evolution, galactic evolution, missing mass. A general outline of the basic physics entering the structure and the evolution of BDs can be found in the excellent reviews of Stevenson (1991) and Burrows & Liebert (1993). It is the aim of this review to summarize the most recent progress realized in the theory of LMS, BDs and EGPs. We will focus mainly on the internal structure and the evolution of these objects since a comprehensive review on the atmosphere of LMS and BDs has appeared recently (Allard et al 1997). We will also consider the implications of LMS and BDs in a more general galactic context and evaluate their contribution to the Galactic mass budget. ## 2 THE PHYSICS OF DENSE OBJECTS ### 2.1 Interior Physics #### 2.1.1 Equation of State Central conditions for LMS, hereafter identified generically as objects below a solar mass, and substellar objects (SSO) for solar composition range from a maximum density $`\rho _c10^3\mathrm{gcm}^3`$ at the hydrogen-burning limit ($`m0.07M_{}`$, see below) to $`\rho _c10\mathrm{gcm}^3`$ for Saturn ($`=5\times 10^4M_{}`$) at 5 Gyr, and from $`T_c10^7`$ K for the Sun to $`T_c10^4`$ K for Saturn at the same age, spanning several orders of magnitudes in mass, density and temperature. Effective temperatures range from $`6000`$ K to $`2000`$ K in the stellar domain and extend down to $`100`$ K for Saturn. Molecular hydrogen and other molecules become stable for $`kT<\mathrm{\hspace{0.17em}3}\times 10^2`$ Ryd ($`T5\times 10^3`$ K), a condition encountered in the atmosphere, or even in deeper layers, of most of these objects. Under these conditions, the interior of LMS and SSO is essentially a fully ionized H<sup>+</sup>/He<sup>++</sup> plasma characterized by coupling parameters $`\mathrm{\Gamma }_i=Z^{5/3}e^2/a_ekT0.1`$-$`50`$ for classical ions and $`r_S=Z^{1/3}a/a_00.11`$ for degenerate electrons. With $`\mathrm{\Gamma }_i=Z^{5/3}(2.693\times 10^5\mathrm{K}/T)n_{24}^{1/3}`$, $`r_S=1.39/(\rho /\mu _e)^{1/3}=1.172n_{24}^{1/3}`$, $`n_{24}n_e/10^{24}`$ cm<sup>-3</sup>$`(\rho /1.6605\mathrm{gcm}^3)\mu _\mathrm{e}^1`$ the electron number-density, $`a=(\frac{3}{4\pi }\frac{V}{N_i})^{1/3}=a_eZ^{1/3}`$ the mean interionic distance, $`\mu _e^1=Z/A`$ the electron mean molecular weight, $`\rho `$ the mass-density, $`A`$ the atomic mass, $`a_0`$ the electronic Bohr radius and $`<>`$ denotes the average number-fraction $`X=_ix_iX_i`$ ($`x_i=N_i/_iN_i`$). The temperature is of the order of the electron Fermi temperature $`T_F`$ so that the degeneracy parameter $`\psi =kT/kT_F3.314\times 10^6T(\mu _e/\rho )^{2/3}`$ is of the order of unity. The classical (Maxwell-Boltzman) limit corresponds to $`\psi +\mathrm{}`$, whereas $`\psi 0`$ corresponds to complete degeneracy. The afore-mentioned thermodynamic conditions yield $`\psi 20.05`$ in the interior of LMS and BDs implying that finite-temperature effects for the electrons must be included to describe accurately the thermodynamic properties of the correlated, partially degenerate electron gas. Moreover, the Thomas-Fermi wavelength $`\lambda _{TF}=\left(kT_F/(6\pi n_ee^2)\right)^{1/2}`$ is of the order of $`a`$, so that the electron gas is polarized by the external ionic field and electron-ion coupling must be taken into account in the plasma hamiltonian. Last but not least, the electron average binding energy can be of the order of the Fermi energy $`Ze^2/a_0ϵ_F`$ so that pressure-ionization takes place along the internal profile. Figure 1 illustrates central characteristic quantities for LMS and SSOs from the Sun to Jupiter. Above $`0.4M_{}`$, the structure evolves slowly with increasing mass from a $`n`$=$`3/2`$ towards a $`n=3`$ polytrope, which yields the correct $`P_c`$ for the Sun, due to the growing central radiative core. This leads to increasing central pressures and densities for increasing mass (increasing polytropic index) in this mass range. Below $`0.3`$-$`0.4M_{}`$, the core becomes entirely convective and follows the behaviour of a $`n=3/2`$ polytrope. Since the gas is still in the classical regime ($`\psi >\mathrm{\hspace{0.17em}1}`$), $`mR`$ and the central density increases with decreasing mass, $`\rho _cm^2`$. Below about the H-burning limit, electron degeneracy becomes dominant ($`\psi <\mathrm{\hspace{0.17em}0.1}`$), so that one approaches the relation $`mR^3`$ (for $`\psi `$=0) and density decreases again with decreasing mass, $`\rho _cm^2`$. These various effects yield a non-monotomic behaviour of the central density and pressure with mass with a minimum in the stellar regime around 0.4 $`M_{}`$ and a maximum near the H-burning limit. The equation of state (EOS) of low-mass objects thus requires a detailed description of strongly correlated, polarisable, partially degenerate classical and quantum plasmas, plus an accurate treatment of pressure partial ionization, a severe challenge for dense matter physicists. Several steps towards the derivation of such an accurate EOS for dense astrophysical objects have been realized since the pioneering work of Salpeter (1961), with major contributions by Fontaine et al (1977), Magni and Mazzitelli (1979), Marley and Hubbard (for objects with $`m<0.2M_{}`$), Hümmer & Mihalas (1988) (primarily devoted to conditions characteristic of solar-type stellar envelopes) and Saumon and Chabrier (Chabrier 1990, Saumon & Chabrier 1991, 1992, Saumon et al 1995, hereafter SCVH). We refer the reader to Saumon (1994) and SCVH for a detailed review and an extensive comparison of these different EOS in the domain of interest. The SCVH EOS presents a consistent treatment of pressure ionization and includes significant improvements with respect to previous calculations in the treatment of the correlations in dense plasmas. It compares well with available high-pressure shock-wave experiments in the molecular domain, and with Monte-Carlo simulations in the fully-ionized, metallic domain (see references above for details). Recently laser-driven shock-wave experiments on D<sub>2</sub> have been conducted at Livermore (Da Silva et al 1997, Collins et al 1998) which reach for the first time the pressure-dissociation and ionization domain and thus probe directly the thermodynamic properties of dense hydrogen under conditions characteristic of BDs and giant planets (GP). The relevance of these experiments for the interior of SSOs can be grasped from Figure 1 of Collins et al (1998). These experiments have shown the good agreement between the predictions of the Saumon-Chabrier model and the data (although there is certainly room for improvement). In particular the strong compression factor arising from hydrogen pressure-dissociation and ionization observed in the experiment ($`\rho /\rho _i5.8`$) agrees well with the predicted theoretical value (see Figure 3 of Collins et al 1998). Any EOS devoted to the description of the interior of dense astrophysical objects must now be confronted with these available data. The SCVH EOS is a pure hydrogen and helium EOS, based on the so-called additive-volume-law between the pure components (H and He). The accuracy of the additive-volume-law has been examined in detail by Fontaine et al (1977). The invalidity of this approximation to describe accurately the thermodynamic properties of the mixture is significant only in the partial ionization region (see SCVH), which concerns only a few percent of the stellar mass under LMS and BD conditions. The effect of metals on the structure and the evolution of these objects has been examined in detail by Chabrier & Baraffe (1997, §2.1). As shown by these authors, because of their negligible number-abundance ($`0.2\%`$), metals do not contribute significantly to the EOS and barely modify the structure and evolution ($`1\%`$ in $`T_{\mathrm{eff}}`$ and $`4\%`$ in $`L`$) of these objects. In the low-density limit characteristic of the atmosphere, the SCVH EOS recovers the perfect gas limit and thermal contributions from various atomic or molecular species can be added within the afore-mentioned additive-volume law formalism, which is exact in this regime. Only in the (denser) metal-depleted atmospheres ($`Z<\mathrm{\hspace{0.17em}10}^2Z_{}`$) of the densest objects ($`<\mathrm{\hspace{0.17em}0.1}M_{}`$), one finds a slight departure from ideality, with a 1% to 4% effect on the adiabatic gradient $`_{ad}`$ (Chabrier & Baraffe 1997). #### 2.1.2 Nuclear rates. Screening factors Although the complete PP chain is important for nucleosynthesis, the thermonuclear processes relevant from the energetic viewpoint under LMS and BD conditions are given by the PPI chain (see e.g. Burrows & Liebert 1993, Chabrier & Baraffe 1997) : $`p+pd+e^++\nu _e;p+d^3He+\gamma ;^3\mathrm{He}+^3\mathrm{He}^4\mathrm{He}+\mathrm{\hspace{0.17em}2}p`$ (1) Below $`0.7M_{}`$, the PPI chain contributes to more than 99% of the energy generation on the zero-age Main Sequence (ZAMS), and the PPII chain to less than 1%. The destruction of $`{}_{}{}^{3}\mathrm{He}`$ by the above reaction is important only for T $`>6\times 10^6`$ K, i.e masses $`m>\mathrm{\hspace{0.17em}0.15}M_{}`$ for ages $`t<`$ 10 Gyr, for which the lifetime of this isotope against destruction becomes eventually smaller than a few Gyr. As examined in §4.6, the abundance of light elements ($`D,Li,Be,B`$) provides a powerful diagnostic to identify the age and/or the mass of SSOs. The rates for these reactions (see e.g. Nelson et al 1993b) in the vacuum, or in an almost perfect gas where kinetic energy largely dominates the interaction energy, are given by Caughlan and Fowler (1988) and Ushomirsky et al (1998) for updated values. The reaction rate $`R_0`$ (in cm<sup>-3</sup>$`s^1`$) in the vacuum is given by $`R_0e^{3ϵ_0/kT}`$ where $`ϵ_0`$ corresponds to the Gamow-peak energy for non-resonant reactions (Clayton 1968). However, as mentioned in the previous section, non-ideal effects dominate in the interior of LMS and BDs and lead to polarization of the ionic and electronic fluids. These polarization effects due to the surrounding particles yield an enhancement of the reaction rate, as first recognized by Schatzman (1948) and Salpeter (1954). The distribution of particles in the plasma reads : $`n(r)=\overline{n}e^{Ze\varphi (r)/kT},\mathrm{with}\varphi (r)={\displaystyle \frac{Ze}{r}}+\psi (r)`$ (2) where $`\psi (r)`$ is the induced mean field potential due to the polarization of the surrounding particles. This induced potential lowers the Coulomb barrier between the fusing particles and thus yields an enhanced rate in the plasma $`R=E\times R_0`$ where $`E=lim_{r0}\left\{g_{12}(r)exp({\displaystyle \frac{Z_1Z_2e^2}{rkT}})\right\}`$ (3) is the enhancement (screening) factor and $`g_{12}(r)`$ the pair-distribution function of particles in the plasma. Under BD conditions, not only ion screening must be considered but also electron screening, i.e. $`E=E_i\times E_e`$. Both effects are of the same order ($`E_iE_e`$ a few) and must be included in the calculations for a correct estimate of the light element-depletion factor (see §4.6). Figure 2 portrays the evolution of the central temperature for objects respectively above, at the limit of and below the hydrogen-burning minimum mass, with lines indicating the hydrogen, lithium and deuterium burning temperatures in the plasma. #### 2.1.3 Transport Properties Energy in the interior of LMS below $`0.4M_{}`$, BDs and GPs is transported essentially by convection (see e.g. Stevenson 1991). According to the mixing length theory (MLT), the convective flux reads (Cox and Giuli 1968) : $$F_{\mathrm{conv}}\rho v_{\mathrm{conv}}C_\mathrm{P}\delta T(Q^{1/2}/\mathrm{\Gamma }_1^{1/2})(\frac{l}{H_\mathrm{P}})^2\rho C_\mathrm{P}c_\mathrm{S}T(_\mathrm{e})^{3/2}$$ (4) where $`Q=(\frac{\mathrm{ln}\rho }{\mathrm{ln}T})_P`$ is the volume expansion coefficient, $`C_\mathrm{P}`$ is the specific heat at constant pressure, $`c_\mathrm{S}=(\mathrm{\Gamma }_1P/\rho )^{1/2}`$ is the adiabatic speed of sound, $`\delta T`$ is the temperature excess between the convective eddy and the surrounding ambient medium, $``$ the temperature gradient, $`v_{\mathrm{conv}}`$ is the convective velocity, $`l\mathrm{and}H_\mathrm{P}`$ denote the mixing length and the pressure scale height, respectively, and $`_\mathrm{e}_{\mathrm{ad}}`$ is the eddy temperature gradient. The last term in parenthesis on the r.h.s. of (4) defines the fluid superadiabaticity, i.e. the fractional amount by which the real temperature gradient exceeds the adiabatic temperature gradient. Below a certain mass, the inner radiative core vanishes and the star becomes entirely convective (VandenBerg et al 1983, D’Antona & Mazzitelli 1985, Dorman et al 1989). This minimum mass, calculated with consistent non-grey atmosphere models and with the most recent OPAL radiative opacities (Iglesias and Rogers 1996) for the interior, is found to be $`m_{conv}=0.35M_{}`$ within the metallicity range $`10^2Z/Z_{}1`$ (Chabrier & Baraffe 1997 §3.2). The Rayleigh number is defined as $`Ra=\frac{gQH_P^3}{\xi \nu }(_{\mathrm{ad}})`$, where $`g10^3`$-$`10^5`$ cm s<sup>-2</sup> is the surface gravity, $`\xi `$ is the thermal diffusivity (conductive or radiative diffusivity) and $`\nu `$ is the kinematic viscosity. In the interior of LMS, $`Ra10^{25}`$ so that convection is almost perfectly adiabatic and the MLT provides a fairly reasonable description of this transport mechanism. Variation of the MLT parameter $`\alpha =l/H_P`$ between 1 and 2 in the interior is found to be inconsequential below $`0.6M_{}`$ (see Baraffe et al 1997). However convection can be affected or even inhibited by various mechanisms. Maximum rotation velocities for LMS and BDs are usually in the range $`2030\mathrm{km}\mathrm{s}^1`$ (see e.g Delfosse et al 1998a) with values as large as 50 $`\mathrm{km}\mathrm{s}^1`$ and 80 $`\mathrm{km}\mathrm{s}^1`$ in the extreme case of Kelu 1 (Ruiz et al 1997, Basri et al 2000). This corresponds to an angular velocity $`\mathrm{\Omega }5\times \mathrm{\hspace{0.17em}10}^4`$ rad $`\mathrm{s}^1`$ for a characteristic radius $`R0.1R_{}`$ ($`\mathrm{\Omega }_{Jup}=1.76\times 10^4`$ rad $`\mathrm{s}^1`$). Within most ($`>95\%`$ in mass) of the interior of LMS and BDs, $`v_{\mathrm{conv}}`$ is of the order of $`10^2`$ cm $`\mathrm{s}^1`$ ($`<<c_S`$) and $`H_\mathrm{P}10^9`$ cm. Thus, the Rossby number $`Ro=v_{\mathrm{conv}}/\mathrm{\Omega }l10^3`$, where $`lH_P`$, and convection can in principle be inhibited by rotation. However, the fact that lithium is not observed in objects with $`m>\mathrm{\hspace{0.17em}0.06}M_{}`$ (see §4.6) suggests that some macroscopic transport mechanism (e.g. meridional circulation, turbulence or convection) remains efficient throughout the star. The Reynolds number $`Re=v_{conv}l/\nu `$ remains largely above unity in stellar and SSO interiors so that convection is not inhibited by viscosity. Magnetic inhibition requires a magnetic velocity $`v_A=(B^2/4\pi \rho )^{1/2}>v_{conv}`$ (Stevenson 1979, 1991), i.e. $`B>\mathrm{\hspace{0.17em}10}^4`$ Gauss under the conditions of interest, only slightly above the predicted fields in these objects (see §2.3). However, it has been proposed by Stevenson (1979) that in a rapidly rotating fluid with a significant magnetic field, the Proudman-Taylor theorem no longer applies so that vertical convective motions could be possible. Lastly, convection efficiency can be diminished by the presence of a density gradient. As shown by Guillot (1995), this may occur in the interior of gaseous planets, due to a gradient of heavy elements, leading to inefficient convective transport in the envelope of these objects. The treatment of convection in the outer (molecular) layers, above and near the photosphere, is a more delicate question. The Rayleigh number in this region is $`Ra10^{15}`$, a rather modest value by the usual standards in turbulence. Since the MLT, by definition, applies to the asymptotic regime $`Ra\mathrm{}`$, it is no longer valid near the photosphere. Much work has been devoted to the improvement of this formalism and to the derivation of a non-local treatment of convection. One of the most detailed attempts is due to Canuto and Mazitelli (Canuto & Mazzitelli 1991, CM). The original CM formalism was based on a linear stability analysis, whereas energy transport by turbulence is a strongly non-linear process. It has been improved recently by including to some extent non-linear modes in the energy rate (Canuto et al 1996) but it yields similar results as the initial CM model. The formalism, however, still requires the calibration of a free parameter, which represents a characteristic ”mixing” scale and in this sense resembles the MLT. A more severe shortcoming of the CM formalism is that the predicted outermost limit of the convection zone for the Sun and the amount of superadiabaticity in this zone are in significant disagreement, both quantitatively and qualitatively, with the results obtained from coupled hydrodynamic-radiation 3D simulations (Demarque et al 1999, Nordlund & Stein 1999). The 3D simulations yield excellent agreement with the observational data at several diagnostic levels: the thermal structure (and the inferred depth of the convection zone), the dynamical structure (vertical velocity amplitude and spectral line synthesis) and the p-mode frequencies and amplitudes for the Sun, without the use of free parameters (see e.g. Spruit & Nordlund 1990 for a review, and Nordlund & Stein 1999). This provides now a high degree of confidence in 3-D hydrodynamic models of stellar surface layers and the inferred transition from convective to radiative energy transport. The thermal structure, for example, is a very robust property of the numerical models, since it depends relatively little on the level of turbulence (and thus on numerical resolution) (Stein & Nordlund 2000; Nordlund & Stein 1999). These results illustrate the fact that any formalism based on a homogeneous description of convection (e.g. MLT, CM) can not describe accurately this strongly inhomogeneous process. Interestingly enough, the standard MLT is found to compare reasonably well with these simulations, at least for the thermal profile (see afore-mentioned references), and thus seems to offer a reasonable (or least worst!) overall description of convective transport, even in the small-efficiency convective regions. Indeed, in contrast to 1-D models computed with the CM formalism, the 3-D models do not have a steeper and more narrow superadiabatic structure for the Sun than obtained with the classical Böhm-Vitense mixing length recipe (Nordlund & Stein 1999). It it thus fair to say that, in the absence of a correct non-local treatment of convection in LMS and BD interiors, the standard MLT is probably the most reasonable choice, at least for the present objects, LMS and SSOs. Clearly the developpement of 3-D hydrodynamics models of the atmosphere of these objects, and the calibration of the mixing length from these simulations, as done for solar-type stars (Ludwig et al 1999), represents one of the next major challenges in the theory. Another possible mechanism of energy transport in stellar interiors is conduction. Its efficiency can be estimated as follows: the distance $`l`$ over which the temperature changes significantly is $`l(\chi t)^{1/2}`$, where $`\chi `$ is the thermal diffusivity and $`t`$ is the time during which the temperature change occurs. For $`\chi 10^1`$ cm<sup>2</sup>s<sup>-1</sup>, characteristic of metals, and $`t10^9`$ yr, $`l10^2`$ km. Heat can thus possibly be transported by conduction only over a limited range of the interior, providing the density is high enough and the temperature low enough for electron conductivity to become important. Indeed, below the HBMM $`0.07M_{}`$, the interior becomes degenerate enough during cooling that the conductive flux $`F_{\mathrm{cond}}=K_{\mathrm{cond}}T`$, where $`K_{\mathrm{cond}}=\frac{4ac}{3}\frac{T^3}{\rho \kappa _{\mathrm{cond}}}`$ and $`\kappa _{\mathrm{cond}}`$ is the conductive opacity, becomes larger than the convective flux. Old enough BDs in the mass-range 0.02-0.07 $`M_{}`$ become degenerate enough to develop a conductive core, which slows down the cooling $`L(t)`$ (Chabrier et al 2000). ### 2.2 Atmosphere Knowledge of the atmosphere is needed for two reasons : (i) as a boundary condition for the interior profile in the optically-thick region and (ii) as a description of the emergent radiative flux. A comprehensive review of the physics of the atmosphere of LMS and BDs can be found in Allard et al (1997). Only a general outline of the main characteristics of these atmospheres will be mentioned in the present review. #### 2.2.1 Spectral distribution Surface gravity $`g=Gm/R^2`$ for LMS and SSOs range from $`\mathrm{log}g4.4`$ for a main sequence 1.0 $`M_{}`$ star to $`\mathrm{log}g3.4`$ for Jupiter, with a maximum $`\mathrm{log}g5.5`$ at the H-burning limit, for solar metallicity. This yields $`P_{ph}g/\overline{\kappa }0.110`$ bar and $`\rho _{ph}10^6`$-$`10^4\mathrm{gcm}^3`$ at the photosphere. Collision effects become significant under these conditions and induce molecular dipoles on e.g. H<sub>2</sub> or He-H<sub>2</sub>, yielding so-called collision-induced absorption (CIA) between roto-vibrational states of molecules which otherwise would have only quadrupolar transitions (Linsky 1969; Borysow et al 1997). At first order this CIA coefficient scales as $`\kappa _{CIA}n_{H_2}^2\times \kappa _{H_2H_2}`$ and thus becomes increasingly important as soon as H<sub>2</sub> molecules become stable. The CIA of H<sub>2</sub> suppresses the flux longward of 2 $`\mu `$m in the atmosphere of LMS, BDs and GPs. This is one of the reasons (the main one for metal-depleted objects) for the redistribution of the emergent radiative flux toward shorter wavelengths in these objects (Lenzuni, Chernoff & Salpeter 1991; Saumon et al 1994, Allard & Hauschildt 1995, Baraffe et al 1997). The CIA of H<sub>2</sub> and H<sub>2</sub>-He, and the bound-free and free-free opacities of H<sup>-</sup> and H$`{}_{}{}^{}{}_{2}{}^{}`$ provide the main continuum opacity sources below $`5000`$ K for objects with metal-depleted abundances. Below $`T_{\mathrm{eff}}<\mathrm{\hspace{0.17em}4000}`$ K, most of the hydrogen is locked into H<sub>2</sub> and most of the carbon in CO. Excess oxygen is bound in molecules such as TiO, VO and H<sub>2</sub>O, with some amounts also in OH and O (see e.g. Figure 1 of Fegley & Lodders 1996). Metal oxides and metal hydrides (FeH, CaH, MgH) are also present. The energy distribution of solar-abundance M-dwarfs is thus entirely governed by the line absorption of TiO and VO in the optical, and H<sub>2</sub>O and CO in the infrared, with no space for a true continuum. The VO and TiO band strength index are used to classify M-dwarf spectral types (see e.g. Kirkpatrick et al 1991). The strongly frequency-dependent absorption coefficient due to the line transitions of these molecules, as well as Rayleigh scattering ($`\nu ^4`$) shortward of $`0.4\mu `$m, yield a strong departure from a black-body energy distribution (see e.g. Figure 5 of Allard et al 1997). At $`T_{\mathrm{eff}}2000`$ K, near the H-burning limit, signatures of metal oxides and hydrides (TiO, VO, FeH, CaH bands) disappear from the spectral distribution (although some amount of TiO remains present in the atmosphere), as observed e.g. in GD 165B (Kirkpatrick et al 1999b). Alkalies are present under their atomic form. The disappearance of TiO-bands prompted astronomers to suggest a new spectral type classification, the ”L”-type, for objects with $`T_{\mathrm{eff}}`$ below the afore-mentioned limit (Martín et al 1997, 1998, Kirkpatrick et al 1999a). Below a local temperature $`T13001500`$ K for $`P310`$ bars, for a solar abundance-distribution, carbon monoxide CO was predicted to dissociate and the dominant equilibrium form of carbon was predicted to become CH<sub>4</sub> (Allard & Hauschildt 1995, Tsuji et al 1995, Fegley & Lodders 1996). Note that this transition occurs gradually with some of the two elements being present in the stability field of the other (see e.g. Fegley & Lodders 1996, Lodders 1999) so that cool stars could contain some limited amount of methane and CO may be visible in objects with $`T_{\mathrm{eff}}<1800`$ K, as indeed detected in Gl229B (Noll et al 1997; Oppenheimer et al 1998). This prediction has been confirmed by the spectroscopic observation of Gliese229B and the identification of methane absorption features at 1.7, 2.4 and 3.3 $`\mu `$m (Oppenheimer et al 1995). The presence of methane in its spectrum confirmed unambiguously its sub-stellar nature (Oppenheimer et al 1995, Allard et al 1996, Marley et al 1996). Objects like Gl229B, characterized by the strong signature of methane absorption in their spectra, define a new spectral class of objects, the ”methane” (sometimes called ”T”) BDs. Methane absorption in the H and K bands yields a steep spectrum at shorter wavelengths and thus blue near-infrared colors, with $`JK<\mathrm{\hspace{0.17em}0}`$ but $`IJ>\mathrm{\hspace{0.17em}5}`$ (Kirkpatrick et al 1999a, Allard 1999). Although the transition temperature between ”L” and ”methane” dwarfs is not determined precisely yet, it should lie in the range 1000$`<T_{\mathrm{eff}}<`$ 1700 K (see §4.5.2). Below $`T_{\mathrm{eff}}2800`$ K, complex O-rich compounds condense in the atmosphere, slightly increasing the carbon:oxygen abundance ratio (see e.g. Tsuji et al 1996, Fegley & Lodders 1996, Allard et al 1997). Different constituents will condense at a certain location in the atmosphere with an abundance determined by chemical equilibrium conditions (although non-equilibrium material may form, depending on the time scale of the reactions, as mentioned below) between the gas phase and the condensed species. The formation of condensed species depletes the gas phase of a number of molecular species (e.g. VO, TiO which will be sequestered into perovskite CaTiO<sub>3</sub>), modifying significantly the emergent spectrum (see e.g. Fegley & Lodders 1994 for the condensation chemistry of refractory elements in Jupiter and Saturn). The equilibrium abundances can be determined from the Gibbs energies of formation either by minimization of the total Gibbs energy of the system (Sharp & Huebner 1990, Burrows & Sharp 1999), or by computing the equilibrium pressures of each grain species (Grossman 1972, Alexander & Ferguson 1994, Allard et al 1998). At each temperature, the fictitious pressure $`P_C`$ of each condensed phase under consideration is calculated from the partial pressures of the species $`i`$ which form the condensate (e.g. $`Al`$ and $`O`$ for corundum $`Al_2O_3`$) $`P_i=N_iT`$, where $`N_i`$ is the number of moles of species $`i`$ and $``$ is the gas constant, determined by the vapor phase equilibria. This fictitious pressure $`P_C`$ is compared to the equilibrium pressure $`P_{eq}`$, calculated from the Gibbs energy of formation of the condensate. The abundance of a condensed species is determined by the condition that this species be in equilibrium with the surrounding gas phase, $`P_CP_{eq}`$ (Grossman 1972). The opacities of the grains are calculated from Mie theory. As suggested by Fegley & Lodders (1996) and Allard et al (1998), refractory elements, Al, Ca, Ti, Fe and V are removed from the gaseous atmosphere by grain condensation at about the corundum (Al<sub>2</sub>O<sub>3</sub>) or perovskite (CaTiO<sub>3</sub>) condensation temperature $`T<\mathrm{\hspace{0.17em}1800}`$ K. Rock-forming elements (Mg, Si, Fe) condense as iron and forsterite (Mg<sub>2</sub>SiO<sub>4</sub>) or enstatite (MgSiO<sub>3</sub>) within about the same temperature range (depending on P). Therefore the spectral features of all these elements will disappear gradually for objects with $`T_{\mathrm{eff}}`$ below these temperatures (see e.g. Fegley & Lodders 1994, 1996, Lodders 1999, Burrows & Sharp 1999 for detailed calculations). For jovian-type effective temperatures ($`T_{\mathrm{eff}}125`$ K), H<sub>2</sub>O and NH<sub>3</sub> condense near and below the photosphere, and water and ammonia bands disappear completely for $`T_{\mathrm{eff}}<\mathrm{\hspace{0.17em}150}`$ K and 80 K, respectively (Guillot 1999). The gas abundance strongly depends on pressure and temperature so that the abundances of various species varies significantly with the mass (and $`T_{\mathrm{eff}}`$) of the astrophysical body. As shown e.g. in figure 2 of Lodders (1999), the condensation point of the dominant clouds lies much closer to the photosphere for M-dwarfs than for Gl229B or - worse - for Jupiter. In other words the location of a given grain condensation lies deeper in Jupiter than in Gl229B. For late M-dwarfs and for massive and/or young BDs, the main cloud formation is predicted to occur very near the photosphere. This is consistent with the fact that all the DENIS and 2MASS objects discovered near the bottom of and below the MS exhibit strong thermal heating and very red colors (Delfosse et al 1998b, Kirkpatrick et al 1999a)(see §4.5.2). Indeed the atmospheric heating due to the large grain opacity (the so-called greenhouse or backwarming effect), and the resulting enhanced H<sub>2</sub>O dissociation, yield a redistribution of the IR flux, as proposed initially by Tsuji et al (1996). A key question for the grain formation process is the size of the grains (see e.g. Alexander & Fergusson 1994). The suppression of the flux in the optical in Gl229B suggests grain sizes of $`0.1\mu `$m as a source of continuum opacity (Jones & Tsuji 1997, Griffith et al 1998), although the wings of alkali resonance lines (K I, Na I) is also a source of absorption in this region (Tsuji et al 1999, Burrows et al 2000). Recent calculations assume a grain-size distribution in the submicron range (Allard et al 1999). Inclusion of the grain absorption with this size distribution in the atmosphere of objects near the limit of the MS successfully reproduces the observed colors of GD 165B, Kelu-1 and of DENIS objects (Leggett et al 1998, Allard 1999, Goldman et al 1999, Chabrier et al 2000). However, the spectrum of GL229B shows no indication for dust in its IR spectrum from 1 to 5 $`\mu m`$ (Allard et al 1996, Marley et al 1996, Tsuji et al 1996, Oppenheimer et al 1998, Schultz et al 1998). This suggests ongoing dynamical processes such as grain settling in SSO atmospheres. Indeed, as noted by Chabrier et al (2000), although convection occurs only in optically-thick regions ($`\tau >1`$) for L-dwarfs, the top of the convection zone lies only about one pressure scale height or less from the photosphere ($`\tau 1`$) near which most grains condense for objects near the bottom of the MS. This can have important consequences on the formation and settling of atmospheric grains. Indeed, although the temperature at the top of the convective zone is found to be generally above the condensation temperature of all grains, convection-induced advection or turbulent diffusion could efficiently bring material upward to the region of condensation and maintain small-grain layers, which otherwise will settle gravitationally, In general, the grain formation process involves a balance between various dynamical timescales, such as the condensation, evaporation, coagulation, coalescence and convection timescales (see e.g. Rossow 1978), not mentioning large-scale hydrodynamical instabilities, typical of weather conditions on Earth. See for example Marley et al (1999) for an early attempt at modelling BD and EGP atmospheres. Moreover, although to an excellent first approximation the atmosphere of the solar jovian planets are in thermodynamical equilibrium, non-equilibrium species with chemical equilibrium timescales larger than the convection timescale can be dredged up to the photosphere by convection, as e.g. CO, PH<sub>3</sub>, GeH<sub>4</sub> in Jupiter (Fegley & Lodders 1994). A correct understanding and a reliable description of this complex grain formation process represent a major challenge for theorists and the observation of spectra of SSOs over a significant temperature-range is necessary to provide guidance through this maze. #### 2.2.2 Transport Properties As mentioned previously, below $`5000`$ K, i.e. $`m0.6`$-$`0.8M_{}`$ depending on the metallicity, hydrogen atoms recombine, $`n_{H_2}`$ increases, as does the opacity through H<sub>2</sub> CIA. The radiative opacity $`\kappa `$ increases by several orders of magnitude over a factor 2 in temperature, which in turn decreases the radiative transport efficiency ($`_{rad}\frac{}{\kappa }`$). On the other hand, the presence of molecules increases the number of internal degrees of freedom (vibration, rotation) and thus the molar specific heat $`C_p`$, which in turn decreases the adiabatic gradient $`(dT/dP)_{ad}=\frac{}{C_p}\frac{T}{P}`$. Both effects strongly favor the onset of convection in the optically-thin ($`\tau <1`$) atmospheric layers (Copeland et al 1970, Allard 1990, Saumon et al 1994, Baraffe et al 1995) with a maximum radial extension for the convection zone around $`T_{\mathrm{eff}}3000`$ K. The combination of relatively large densities, opacities and specific heat in layers where H<sub>2</sub> molecules are present contribute to a very efficient convective transport. Convection is found to be adiabatic almost up to the top of the convective region and the extension of the superadiabatic layers is very small compared to solar type stars (Brett 1995, Allard et al 1997). Thus a variation of the mixing length between $`H_\mathrm{P}`$ and 2$`H_\mathrm{P}`$ barely affects the thermal atmospheric profile (Baraffe et al 1997 §3). The presence of convection in the optically thin layers precludes radiative equilibrium in the atmosphere ($`_{tot}_{rad}`$) and requires the resolution of the transfer equation for radiative-convective atmosphere models for objects below $`0.7M_{}`$. This, and the afore-mentioned strong frequency-dependence of the molecular absorption coefficients, exclude grey model atmosphere and/or the use of radiative $`T(\tau )`$ relationships to determine the outer boundary condition to the interior structure. As demonstrated by Saumon et al (1994) and Baraffe et al (1995, 1997, 1998), such outer boundary conditions overestimate the effective temperature and the HBMM and yield erroneous $`m`$-$`T_{\mathrm{eff}}`$ and $`m`$-$`L`$ relationships, two key relations for the calibration of the temperature-scale and for the derivation of mass functions (see Chabrier & Baraffe 1997, §2.5 for a complete discussion). Correct evolutionary models for low mass objects require (i) the connection between the non-grey atmospheric ($`P`$,$`T`$) profile, characterized by a given ($`\mathrm{log}g,T_{\mathrm{eff}}`$) and the interior ($`P`$,$`T`$) profile at a given optical depth $`\tau `$, preferably large enough for the atmospheric profile to be adiabatic and (ii) consistency between the atmospheric profiles and the synthetic spectra used to determined magnitudes and colors. For a fixed mass and composition, only one atmosphere profile matches the interior profile for the afore-mentioned boundary condition. This determines the complete stellar model (mass, radius, luminosity, effective temperature, colors) for this mass and composition (Chabrier & Baraffe 1997, Burrows et al 1997 for objects below 1300 K (their figure 1)). ### 2.3 Activity The Einstein and Rosat surveys of the solar neighborhood and of several open clusters have allowed the identification of numerous M-dwarfs as faint X-ray sources, with typical luminosities of the order of the average solar luminosity $`L_X10^{27}`$ erg s<sup>-1</sup>, reaching up to $`10^{29}`$ erg s<sup>-1</sup> in young clusters (see e.g. Randich, 1999 and references therein for a recent review). This suggests that late M-dwarfs are as efficient coronal emitters as other cool stars in terms of $`L_X/L_{bol}`$, which expresses the level of X-ray activity. Moreover about 60% of M-dwarfs with spectral-type $`>M5`$ ($`M_{\mathrm{bol}}12`$) show significant chromospheric activity with $`\mathrm{log}(L_{H\alpha }/L_{bol})3.9`$ (Hawley et al 1996). These observations provide important empirical relationships between age, rotation and activity. The $`L_X/L_{bol}`$ ratio seems to saturate above a rotational velocity threshold at a limit $`\mathrm{log}(L_X/L_{bol})3`$, implying that the intrinsic coronal emission of LMS and BDs is quite low, decreases with $`L_{bol}`$ and thus mass and does not increase with increasing rotation above the threshold limit. This suggests a saturation relation in terms of a Rossby number, i.e. the ratio of the rotational period over the convective turnover time $`Ro=\frac{P}{t_{conv}}`$ (see §2.1.3). Because $`t_{conv}`$ increases for lower-mass stars, they will saturate at progressively larger periods, i.e. lower rotational velocities. The velocity threshold for a 0.4 $`M_{}`$ star is estimated around 5-6 km s<sup>-1</sup> (Stauffer et al 1997). Comparison of the Pleiades and the Hyades samples shows a steep decay of the X-ray activity in the Hyades. This result is at odds with a simple, monotonic rotation-activity relationship. The spin-down timescales for M-dwarfs are much longer than for solar-type stars. Hyades M-dwarfs do indeed show moderate or even rapid rotation ($`v\mathrm{sin}i10`$ km s$`1`$) and thus should show strong X-ray emission. The fainter $`L_X`$ in the Hyades than in the Pleiades thus reinforces the saturation scenario. Delfosse et al (1998a) obtained projected rotational velocities and fluxes in the H<sub>α</sub> and H<sub>β</sub> lines for a volume-limited sample of 118 field M-dwarfs with spectral-type M0-M6. They found a strong correlation between rotation and both spectral type (measured by $`RI`$ colors) and kinematic population: all stars with measurable rotation are later than M3.5 and have kinematic properties typical of the young disk population, or intermediate between young disk and old disk. They interpret this correlation as evidence for a spin-down timescale that increases with decreasing mass, with this timescale being a significant fraction of the age of the young disk ($`3`$ Gyr) at spectral type M4 ($`0.15M_{}`$). These data confirm the saturation relation inferred previously for younger or more massive stars: $`L_X/L_{bol}`$ and $`L_{H_\alpha }/L_{bol}`$ both correlate with $`v\mathrm{sin}i`$ for $`v\mathrm{sin}i<\mathrm{\hspace{0.17em}4}`$-5 km s<sup>-1</sup> and then saturate at 10<sup>-2.5</sup> and 10<sup>-3.5</sup>, respectively. Recently, Neuhäuser and Comeron (1998) have detected X-ray activity in young BDs with $`m0.04M_{}`$ in the Chamaleon star-forming region, with $`\mathrm{log}L_X=28.41`$ and $`\mathrm{log}L_X/L_{bol}=3.44`$. Coronal activity in BDs has been confirmed by Neuhäuser et al (1999), with all the objects belonging to young clusters or star-forming regions. All these observations show that there is no drop in $`L_X/L_{bol}`$ in field stars up to a spectral type of $`M7`$, which is well below the limit $`m0.3M_{}`$ where stellar interiors are predicted to become fully convective. However, recent observations of the 2MASS objects (Kirkpatrick et al 1999, Fig. 15a) seem to show a significant decrease of activity for L-type objects, which confirms the decline of activity suggested previously for very late M-spectral types ($`>M8`$)(Hawley et al 1996, Tinney & Reid 1998). In any event, X-ray emission does not disappear in fully convective stars and even at least young SSOs can support magnetic activity. A fossil field can survive only over a timescale $`\tau _dR^2/\eta `$ a few years for fully convective stars ($`\eta \eta _t`$ is the turbulent magnetic diffusivity and $`R`$ the stellar radius) so that a dynamo process is necessary to generate the magnetic field. The data suggest that (i) the dynamo-generation believed to be at work in the Sun does not apply for very-low mass objects, (ii) dynamo generation in fully convective stars is as efficient as in stars with a radiative core and (iii) whatever the dynamo mechanism is, it is bounded by the saturation condition at least for late-type stars. It is generally admitted that magnetic activity in the Sun results from the generation of a large-scale toroidal field by the action of differential rotation on a poloidal field at the interface between the convective envelope and the radiative core, where differential rotation is strongest, known as the tachocline (Spiegel & Weiss 1980, Spiegel & Zahn 1992). In this region a sufficiently strong magnetic field is stable against buoyancy. In this case the shear dominates the helicity, this is the so-called $`\alpha `$$`\mathrm{\Omega }`$ dynamo-generation, which predicts a correlation between activity and rotation, as observed in solar-type stars. As shown initially by Parker (1975), buoyancy prevents the magnetic field generated by the turbulent motions in a convective zone to become global in character. For this reason, the absence of a radiative core, i.e. of a region of weak buoyancy and strong differential rotation, precludes in principle the generation of a large-scale magnetic field. A dynamo generated by a turbulent velocity field that would generate chaotic magnetic fields in the absence of rotation has been proposed as an alternative process for fully convective stars (Durney et al 1993). They found that a certain level of magnetic activity can be maintained without the generation of a large-scale field. The turbulent velocity field can generate a self-maintained small-scale magnetic field providing the magnetic Reynolds number $`Re=(vl/\eta )^2`$ is large enough. The scale of this field is comparable to that of turbulence and in rough energy equipartition with the fluid motions i.e. $`\frac{1}{8\pi }(B_P^2+B_\varphi ^2)\frac{1}{2}\rho v_{conv}^2`$. Rotation is not essential in this case but it increases the generation rate of the field. Recently Küker and Chabrier (in preparation) explored an other possibility, namely the generation of a large-scale field by a pure $`\alpha ^2`$-effect. In the $`\alpha ^2`$-dynamo, helicity is generated by the action of the Coriolis force on the convective motions in a rotating, stratified fluid. The $`\alpha ^2`$-effect strongly depends on the Rossby number, or the equivalent Coriolis number $`\mathrm{\Omega }^{}=2t_{conv}\mathrm{\Omega }=4\pi /Ro`$. In low-mass objects, the convective turnover time is longer than the rotation period (see §2.1.3), so that $`\mathrm{\Omega }^{}>>1`$. They find that the $`\alpha ^2`$-dynamo is clearly supercritical. It generates a large-scale non-axisymmetric steady (co-rotating) field that is symmetric with respect to the equatorial plane. Equipartition of energy yields field strengths of several kiloGauss. The possible decrease of activity in very late-type stars, as observed for example with the small chromospheric activity of the M9.5 BD candidate BRI 0021-0214 in spite of its fast rotation ($`v\mathrm{sin}i=`$40 km s<sup>-1</sup>) (Tinney et al 1998), is a more delicate issue that requires the inclusion of dissipative processes. Indeed, these calculations show that $`\alpha ^2`$-dynamo can efficiently generate a large-scale magnetic field in the interior of fully convective stars. Since conductivity decreases in the outermost layers of the star, however, there no current would be created by the field in these regions, and thus no dissipative process and no activity (see e.g. Meyer & Meyer-Hofmeister 1999). The observed continuous transition in rotation and activity at the fully convective boundary suggests in fact that the $`\alpha ^2`$-dynamo is already at work in the convection zones of the more massive stars. An interesting possibility for verifying the present theory would be Doppler imaging of fast-rotating LMS and BDs. Turbulent dynamo is likely to yield a spatially uniform chromospheric activity whereas the large-scale $`\alpha ^2`$ process suggested by Küker & Chabrier generates asymmetry. Moreover, non-axisymmetric fields can propagate in longitudinal directions without any cyclic variation of the total field energy whereas dynamo waves generated by $`\alpha `$$`\mathrm{\Omega }`$ processes propagate only along the lines of constant rotation rate (Parker 1955). Therefore we do not expect cycles for uniformly rotating (fully convective) stars. ## 3 MECHANICAL AND THERMAL PROPERTIES ### 3.1 Mechanical Properties Figure 3 portrays the mass-radius behaviour of LMS and isolated SSOs from the Sun to Jupiter for $`t=6\times 10^7`$ (dash-dot line) and 5$`\times 10^9`$ yr (solid line) for $`Z=Z_{}`$ and $`Z=10^2\times Z_{}`$ (dash-line). The time required to reach the ZAMS, arbitrarily defined as the time when $`L_{\mathrm{nuc}}`$ = 95% $`L_{\mathrm{tot}}`$, for solar metallicity LMS is given in Table 1. The general $`m`$-$`R`$ behaviour reflects the physical properties characteristic of the interior of these objects, as inferred from Figure 1. For $`m>\mathrm{\hspace{0.17em}0.2}M_{}`$, $`\psi >1`$ for all ages so that the internal pressure is dominated by the classical perfect gas ion+electron contribution ($`P=\rho kT/\mu m_H`$, neglecting the molecular rotational/vibrational excited level contributions) plus the correcting Debye contribution arising from ion and electron interactions ($`P_{DH}\rho ^{3/2}/T^{1/2}`$) and $`Rm`$ in first order from hydrostatic equilibrium. When the density becomes large enough during pre-MS contraction so that $`\psi <1`$, which occurs at t $`>`$$``$ 50 Myr for $`m=0.15M_{}`$ and at t $`>`$$``$ 10 Myr for the HBMM $`m=0.075M_{}`$, the EOS starts to be dominated by the contribution of the degenerate electron gas ($`P\rho ^{5/3}`$). This yields a minimum in the $`m`$-$`R`$ relationship $`R_{\mathrm{min}}0.08R_{}`$ for $`m`$ 0.06-0.07 $`M_{}`$. Full degeneracy ($`\psi 0`$) would yield the well-known zero-temperature relationship $`Rm^{1/3}`$, as in white dwarf interiors, but partial degeneracy and the non-negligible contribution arising from the (classical) ionic Coulomb pressure (which implies $`Rm^{+1/3}`$) combine to yield a smoother relation $`R=R_0m^{1/8}`$ at t=5 Gyr, where $`R_00.06R_{}`$ for $`0.01M_{}<m<\mathrm{\hspace{0.17em}0.07}M_{}`$, i.e. an almost constant radius. For the age and metallicity of the solar system this radius is of the order of the Jupiter radius $`R_J0.10\times R_{}`$. At 5 Gyr, the radius reaches a maximum $`R0.11R_{}`$ for $`m4\times M_\mathrm{J}`$ (where $`M_\mathrm{J}=9.5\times 10^4M_{}`$ is the Jupiter mass) (Zapolsky & Salpeter 1969, Hubbard 1994, Saumon et al. 1996). Below this limit, degeneracy saturates (see Figure 1) and the classical ionic pressure contribution becomes important enough so that we recover a nearly classical behaviour. Figure 3 also displays the astrophysically determined radii of the two eclipsing binary systems YY-Gem (Leung & Schneider 1978) and CM-Dra (Metcalfe et al 1996), and of the white dwarf companion GD448 (Maxted et al 1998). The bump on the 60 Myr isochrone near $`m0.01M_{}`$ results from initial deuterium burning (see §3.2). The $`m`$-$`R`$ relationship is essentially determined by the EOS and is weakly sensitive to the outer boundary condition and the structure of the atmosphere, since the latter represents at most (for 1 $`M_{}`$) a few percent of the total radius of LMSs and SSOs (Dorman et al 1989, Chabrier & Baraffe 1997). The effect of metallicity $`Z`$ on the $`m`$R relation remains modest. A decrease in metallicity yields a slight decrease in the radius at a given age or at the same stage of nuclear burning (i.e the same H content). Indeed lower metallicity yields a larger $`T_{\mathrm{eff}}`$ at a given mass (see §3.2). This in turns implies an increase in the nuclear energy production to reach thermal equilibrium, and thus a larger central temperature $`T_\mathrm{c}`$. Since $`R\mu m/T`$ from hydrostatic equilibrium, where $`\mu `$ is the mean molecular weight, determined essentially by hydrogen and helium, a lower-metallicity star contracts more to reach thermal equilibrium. As shown by Beuermann et al (1998), the variation of $`R`$ with $`Z`$ predicted by the models of Baraffe et al (1998) is in agreement with the radii deduced from observations. ### 3.2 Thermal Properties Figure 4 exhibits the mass-effective temperature relationships for representative ages and metallicities. The arrows indicate the onset of formation of H<sub>2</sub> near the photosphere (Auman 1969, Copeland et al 1970, Kroupa et al 1990). As discussed in §2.2.2, molecular recombination favors convective instability in the atmosphere. Convection yields a smaller T-gradient $`(_{ad})`$ and thus a cooler structure in the deep atmosphere (Brett 1995, Allard & Hauschildt 1995, Chabrier & Baraffe 1997 Figure 5). Therefore, a model with atmospheric convection corresponds to a larger $`T_{\mathrm{eff}}`$ since the ($`P,T)`$ interior-atmosphere boundary is fixed for a given mass (§2.2.2). Furthermore, the adiabatic gradient in the regions of H<sub>2</sub> recombination decreases to a minimum value $`_{ad}0.1`$, compared to $`_{\mathrm{ad}}0.4`$ for an ideal monoatomic gas (Copeland et al 1970, Saumon et al 1995 Figure 17). Therefore, even if the atmosphere were already convective without H<sub>2</sub>, molecular recombination yields a flatter inner temperature gradient in the atmosphere, and thus enhances the former effect, i.e. a larger $`T_{\mathrm{eff}}`$ for a given mass. Formation of H<sub>2</sub> in the atmosphere occurs at higher $`T_{\mathrm{eff}}`$ for decreasing metallicity, because of the denser (more transparent) atmosphere (see Figure 4). The relation between the central and effective temperatures for LMS and SSOs can be inferred from Figures 1 and 4 and is commented on by Chabrier and Baraffe (1997 §4.2) and Baraffe et al. (1998 §2). As these authors demonstrate, a grey approximation or the use of a radiative $`T(\tau )`$ relationship significantly overestimates the effective temperature (from $`50`$ to 300 K depending on the atmospheric treatment). A representative $`T_\mathrm{c}`$-$`\rho _\mathrm{c}`$ diagram is shown in Figure 5 (see also Figure 8 of Burrows et al 1997) and illustrates the different evolutionary paths for LMS and SSOs in this diagram. The bumps appearing on the isochrones between $`10^6`$ and 10<sup>8</sup> yr and $`\mathrm{log}T_\mathrm{c}5.4`$-5.8 result from initial deuterium burning. For stars, $`T_\mathrm{c}`$ and $`\rho _\mathrm{c}`$ always increase with time until they reach the ZAMS. For BDs, $`T_\mathrm{c}`$ first increases for $`10^7`$-$`10^9`$ yr, for masses between $`0.01`$ and $`0.07M_{}`$, respectively. Then, when degeneracy becomes dominant, $`T_\mathrm{c}`$ reaches a maximum and decreases. For objects with $`m<\mathrm{\hspace{0.17em}5}\times 10^3M_{}`$, T<sub>c</sub> always decreases for $`t>\mathrm{\hspace{0.17em}1}`$ Myr. In the substellar regime there is no steady nuclear energy generation, by definition, and the evolution is governed by the change of internal energy $`_M\frac{dE}{dt}𝑑m`$ and the release of contraction work $`_M\frac{P}{\rho ^2}\frac{d\rho }{dt}𝑑m`$. Even when degeneracy becomes important, SSOs keep contracting, though very slowly, within a Hubble time. ## 4 EVOLUTION ### 4.1 Evolutionary tracks Figure 6 exhibits $`T_{\mathrm{eff}}(t)`$ obtained from consistent non-grey calculations for several masses, for $`[\mathrm{M}/\mathrm{H}]`$<sup>1</sup><sup>1</sup>1$`[\mathrm{M}/\mathrm{H}]=\mathrm{log}(\mathrm{Z}/\mathrm{Z}_{})`$=0 (helium mass fraction $`Y=0.275`$) and $`[\mathrm{M}/\mathrm{H}]=2.0`$ ($`Y=0.25`$), respectively. Initial deuterium burning (with an initial mass fraction $`[D]_0=2\times 10^5`$) for masses $`m>\mathrm{\hspace{0.17em}0.013}M_{}`$ proceeds quickly, at the early stages of evolution, and lasts about $`10^6`$-$`10^8`$ years. Objects below this limit are not hot enough to fuse deuterium in their core (see Figure 2). For masses above $`0.07M_{}`$ for $`[M/H]`$=0 and $`0.08M_{}`$ for $`[M/H]`$=-2.0, the internal energy provided by nuclear burning quickly balances the gravitational contraction energy, and after a few Gyr the lowest-mass star reaches complete thermal equilibrium ($`L=ϵ_{\mathrm{nuc}}𝑑m`$, where $`ϵ_{\mathrm{nuc}}`$ is the nuclear energy rate), for both metallicities. The lowest mass for which thermal equilibrium is reached defines the HBMM and the related hydrogen-burning minimum luminosity HBML. Stars with $`m0.4M_{}`$ develop a convective core near the ZAMS for a relatively short time, depending on the mass and metallicity, which results in the bumps for 0.6 and 1 $`M_{}`$ at respectively $`10^8`$ yr and $`3\times 10^7`$ yr (Chabrier and Baraffe 1997, §3.2). The dotted lines portray $`T_{\mathrm{eff}}(t)`$ for objects with solar abundance when grain opacity is taken into account in the atmosphere (see §4.5.2). The Mie opacity due to the formation of refractory silicate grains produces a blanketing effect that lowers the effective temperature and luminosity at the edge of the main sequence, an effect first noticed by Lunine et al (1989). However as a whole, grain formation only moderately affects the evolution near and below the bottom of the main sequence and thus the HBMM and HBML. Models with grainless atmospheres yield $`m=0.072M_{}`$, $`L=5\times 10^5L_{}`$ and $`T_{\mathrm{eff}}=1700`$ K at the H-burning limit, whereas models with grain opacity give $`m=0.07M_{}`$, $`L=4\times 10^5L_{}`$ and $`T_{\mathrm{eff}}=1600`$ K, for solar composition (Chabrier et al 2000). As shown in Figure 6, young and massive BDs can have the same effective temperature (or luminosity) as older very-low-mass stars, a possible source of contamination for the determination of the local stellar luminosity function at the bottom of the MS (see §5). Note the quick decrease of $`T_{\mathrm{eff}}`$ (and $`L`$) with time for objects below the HBMM, $`Lt^\alpha `$ with $`\alpha 5/4`$ (Stevenson 1991, Burrows et al 1994, 1997), with a small dependence of $`\alpha `$ on the presence of grains. Slightly below $`0.072M_{}`$ (resp. 0.083 $`M_{}`$) for $`[M/H]`$=0 (resp. $`[M/H]`$ $``$ -1), nuclear ignition still takes place in the central part of the star, but cannot balance the ongoing gravitational contraction (see Figure 2). Although these objects are sometimes called ”transition objects” we prefer to consider them as massive BDs, because strictly speaking they will never reach thermal equilibrium. Indeed, just below the HBMM, the contributions from the nuclear energy source $`ϵ_{\mathrm{nuc}}𝑑m`$ and the entropy source $`T\frac{dS}{dt}𝑑m`$ are comparable and cooling proceeds at a much slower rate than mentioned above. Below about $`0.07M_{}`$ (resp. $`0.08M_{}`$) for \[M/H\]=0 (resp. \[M/H\] $``$ -0.5), the energetic contribution arising from hydrogen-burning, though still present for the most massive objects, is orders of magnitude smaller than the internal heat content, which provides essentially all the energy of the star ($`ϵ_{\mathrm{nuc}}<<T\|\frac{dS}{dt}\|`$). Once on the ZAMS the radius for stars is essentially constant, whereas for BDs the contraction slows down when $`\psi <\mathrm{\hspace{0.17em}0.1}`$ (Figures 1 and 3). The effects of metallicity on the atmosphere structure (Brett 1995, Allard and Hauschildt 1995) and on the evolution (Saumon et al 1994, Chabrier and Baraffe 1997) have been discussed extensively in the previously-cited references and can be apprehended with simple arguments: the lower the metallicity $`Z`$, the lower the mean opacity $`\overline{\kappa }`$ and the more transparent the atmosphere so that the same optical depth lies at deeper levels and thus higher pressure ($`\frac{dP}{d\tau }=\frac{g}{\overline{\kappa }}`$). Therefore for a given mass ($`\mathrm{log}g`$) the $`(T,P)`$ interior profile matches, for a given optical depth $`\tau `$, an atmosphere profile with larger $`T_{\mathrm{eff}}`$ (dash-line on Figure 4), and thus higher luminosity $`L`$ since the radius depends only weakly on the metallicity. The consequence is a larger HBMM for lower $`Z`$ since a larger $`L`$ requires more efficient nuclear burning to reach thermal equilibrium, and thus a larger mass. A standard way to display evolutionary properties is a theoretical Hertszprung-Russell diagram (HRD). Such a HRD for LMS and SSOs is shown in Figure 7 for several masses and isochrones from 1 Myr to 5 Gyr, whereas Figure 8 shows evolutionary tracks in a $`\mathrm{log}g`$-$`T_{\mathrm{eff}}`$ diagram (see also Burrows et al 1997). These figures allow the determination of the mass and age of an object from the gravity and effective temperature inferred from its spectrum. ### 4.2 Mass-magnitude One of the ultimate goals of stellar theory is an accurate determination of the mass of an object for a given magnitude and/or color. Figure 9 (see also Baraffe et al 1998) shows the comparison between theory and observationally-determined masses in the $`K`$-band. The solid line corresponds to a 5$`\times 10^9`$ yr isochrone, for which the lowest-mass stars have settled on the MS, for a solar metal-abundance ($`[\mathrm{M}/\mathrm{H}]=0`$), whereas the dotted line corresponds to a 10<sup>8</sup> yr isochrone for the same metallicity and the dashed line corresponds to a 10<sup>10</sup> yr isochrone for $`[\mathrm{M}/\mathrm{H}]=0.5`$, representative of the thick-disk population. A striking feature is the weak metallicity-dependence in the K-band, compared to the strong dependence in the V-band (see Figure 3 of Baraffe et al 1998). As discussed in this paper, this stems from two different effects. On one hand, the increasing opacity in the optical, dominated by TiO and VO lines, and the decreasing H<sub>2</sub> opacity in the K-band with increasing metallicity shift the peak of the flux toward larger wavelengths. Thus, for fixed $`T_{\mathrm{eff}}`$ the V-flux decreases and the K-flux increases with increasing $`[\mathrm{M}/\mathrm{H}]`$. On the other hand, for a given mass, the total flux (and $`T_{\mathrm{eff}}`$) decreases with increasing metallicity, as mentioned in §4.1. These two effects add up in the V-band and yield an important variation of the V-flux with metallicity. In the K-band, they cancel and yield similar fluxes for a given mass below $`0.4M_{}`$ ($`T_{\mathrm{eff}}3500`$-$`4300`$ K, depending on $`[\mathrm{M}/\mathrm{H}]`$) at different metallicities. These arguments remain valid as long as H<sub>2</sub> CIA does not significantly depress the K-band, as it does for very metal-depleted objects. ### 4.3 Mass-Spectral type The knowledge of spectral type ($`Sp`$) is extremely useful for analyzing objects with unknown distance or with colors altered by reddening, as in young clusters, for example. The determination of a mass-$`Sp`$ relationship provides a powerful complement to the mass-magnitude relationship for assigning a mass to an object. Based on spectroscopic observations of low-mass nearby composite systems, Kirkpatrick & McCarthy (1994) have determined an empirical $`m`$-$`Sp`$ relationship for M-dwarfs, restricted to M2-M6 spectral types. A theoretical relation has been derived by Baraffe and Chabrier (1996) from M0 to M10, emphasizing the non-linear behavior of this relationship near the bottom of the MS, for $`>`$M6. Figure 10 shows the $`m`$-$`Sp`$ relationship for M-, K- and G-dwarfs based on the Baraffe et al (1998) models. The spectral type for K- and G-dwarfs is derived from the model synthetic color ($`IK`$) through the empirical $`Sp(IK)`$ relation recently derived by Beuermann et al (1998). As stressed by Baraffe and Chabrier (1996), the $`m`$-$`Sp`$ relationship depends crucially on age and metallicity. At $``$ 1 Gyr, objects with $`Sp`$ later than M10 are below the HBMM and can be considered as bona-fide BDs. This limit decreases to $``$ M7 at 100 Myr and $``$ M6 at 10 Myr. A difficulty arises when dealing with very young objects ($`t<`$ 10 Myr) because of gravity effects. No reliable $`Sp`$-color relationship is presently available for these objects for which spectroscopic and photometric properties are found to be intermediate between giants and dwarfs (Luhman 1999). The present $`m`$-$`Sp`$ relationships should be extended in the future to dwarfs cooler than M9.5-M10, the so-called L-dwarfs and methane-dwarfs. Note however that SSOs evolve at different rates through a series of spectral types as they cool, so we cannot associate a given spectral type with a specific mass for these objects. ### 4.4 Irradiated planets The mass-radius relation and evolutionary sequences described above reflect the relations for isolated objects. After a rapid initial accretion phase and subsequent hydrodynamical collapse, planets orbiting stars will evolve differently. As shown by Hubbard (1977), illumination from a parent star will yield thermal expansion of the less massive (low gravity) objects, toward an asymptotic temperature $`T_{eq}`$ set only by the thermalized photons from the parent star and toward a (larger) asymptotic equilibrium radius $`R_p`$, with $`R_p=2a(L_{eq}/L_{})^{1/2}[1/(1A)]^{1/2}`$ and $`T_{eq}=(1A)^{1/4}(R_{}/2a)^{1/2}T_{eff_{}}`$, where $`a`$ is the orbital distance, $`A`$ the Bond albedo, $`L_{}=4\pi \sigma R_{}^2T_{eff_{}}`$ the parent star luminosity, effective temperature and radius and $`L_{eq}=4\pi \sigma R_p^2T_{eq}^4`$ denotes the equilibrium luminosity (see Saumon et al 1994). As noted in Guillot et al (1996) and Guillot (1999), EGPs in close orbits are heated substantially by their parent star and their atmosphere cannot cool substantially. They develop an inner radiative region, as a result of this stellar heating and contract at almost constant $`T_{\mathrm{eff}}`$, at a much smaller rate than if they were not heated by the star. Planets whose luminosity is larger than the absorbed stellar flux of the parent star just evolve along a fully convective Hayashi track for $`10^6`$ yr before reaching the afore-mentioned equilibrium state. The flux from an irradiated planet includes two contributions: the intrinsic thermal emission and the reflected starlight (the Albedo contribution) (see e.g. Seager & Sasselov 1998, Marley et al 1999): $`_\nu =({\displaystyle \frac{R_p}{d}})^2_\nu ^p+({\displaystyle \frac{A}{4}})P(\varphi )({\displaystyle \frac{R_{}}{d}})^2({\displaystyle \frac{R_p}{a}})^2_\nu ^{}`$ (5) where $`d`$ is the distance of the system to Earth and $`P(\varphi )`$ is the dependence of the reflected light upon the phase angle between the star, the planet and Earth ($`P=1`$ if the light reflected by the planet is redistributed uniformly over 4$`\pi `$ steradians). Figure 11 illustrates the flux from a young EGP (assuming the irradiation is isotropic) with $`T_{\mathrm{eff}}=1000`$ K orbiting a Sun-like star at 15 pc for various orbital distances. ### 4.5 Color-magnitude diagrams Thanks to the recent progress both on the observational side, with the development of several ground-based and space-based observational surveys of unprecedented sensitivity in various optical and infrared passbands, and on the theoretical side, with the derivation of synthetic spectra and consistent evolutionary calculations based on non-grey atmosphere models, it is now possible to confront theory with observation directly in the observational planes, i.e. in color-color and color-magnitude diagrams (CMD). This avoids dubious color-$`T_{\mathrm{eff}}`$ or color-$`M_{bol}`$ transformations and allows more accurate determinations of the intrinsic properties ($`m,L,T_{\mathrm{eff}},R`$) of an object from its observed magnitude and/or color. In this section we examine the behaviour of LMS and SSOs in various CMDs characteristic of different populations in terms of age and metallicity. These diagrams capture the essence of the observational signatures of the very mechanical and thermal properties of these objects, described in the previous sections. #### 4.5.1 Pre-main sequence and young clusters Numerous surveys devoted to the detection of SSOs have been conducted in young clusters with ages spanning from $``$ 1-10 Myr to $`>\mathrm{\hspace{0.17em}10}^8`$ yr for the Pleiades or the Hyades (see Martín 1999 and Basri this volume for reviews). Observations of young clusters present two important advantages, namely (i) all objects in the cluster are likely to be coeval within a reasonable range, except possibly in star forming regions, where the spread in ages for cluster members can be comparable with the age of the cluster, (ii) young objects are brighter for a given mass (see Figure 6), which makes the detection of very-low mass objects easier. Conversely, they present four major difficulties, (i) extinction caused by the surrounding dust modifies both the intrinsic magnitude and the colors of the object, (ii) accurate proper motion measurements are necessary to assess whether the object belongs to the cluster, (iii) gravity affects both the spectrum and the evolution, (iv) the evolution and spectrum of very young objects ($`t<`$ 1 Myr) may still be affected by the presence of an accretion disk or circumstellar material residual from the protostellar stage. Young multiple systems remove some of these difficulties and provide excellent tests for models at young ages. Like for example the quadruple system GG TAU (White et al 1999), with components covering the whole mass-range of LMS and BDs from 1 $`M_{}`$ to $``$ 0.02 $`M_{}`$. Models based on non-grey atmospheres (Baraffe et al 1998) are the only ones consistent with these observations (White et al 1999, Luhman 1999). The comparison of $`10^6`$ yr isochrones for SSOs with observations in a near-IR CMD is shown in Zapatero Osorio et al (1999) for the young cluster $`\sigma `$-Orionis (their figure 1). #### 4.5.2 Disk field stars Figures 12 and 13 display the CMDs for LMS and SSOs in the optical and the infrared for various ages and metallicities (see also Burrows et al 1997). Various sets of models correspond to calculations where (i) grain formation is included in the atmosphere EOS but not in the transfer equations, mimicking a rapid settling below the photosphere (COND models), (ii) grain formation is included both in the EOS and the opacity (DUSTY models) (see Allard et al 1999), (iii) grainless models. Models (i) and (ii) represent extreme cases that bracket the complex grain processes at play in these atmospheres (see §2.2.1). The optical sequence shows a monotonic magnitude-color behaviour with 3 changes of slope around $`M_\mathrm{V}`$ 10, 13 and 19 for $`Z=Z_{}`$, respectively, which correspond to 0.5 $`M_{}`$, $`T_{\mathrm{eff}}`$=3600 K; 0.2 $`M_{}`$, $`T_{\mathrm{eff}}`$=3300 K and 0.08 $`M_{}`$, $`T_{\mathrm{eff}}`$ = 2330 K on the 5 Gyr isochrone. They reflect, respectively, the onset of convection in the atmosphere, degeneracy in the core and grain absorption near the photosphere, as described in the previous sections. The latter feature, which yields a steep increase of $`M_\mathrm{V}`$ at $`(VI)5`$ does not appear in dust-free models (cf. Baraffe et al 1998), where TiO is not depleted by grain formation in the atmosphere and thus absorbs the flux in the optical. Interestingly enough, this steepening is observed in the sample collected recently by Dahn et al (1999) which includes some DENIS and 2MASS L-dwarf parallaxes. This effect is also observed in the $`(RI)`$ CMD of the Pleiades (Bouvier et al 1998), in agreement with the theoretical predictions (Chabrier et al 2000). This clearly illustrates the effect of grain formation in the atmosphere which affects the spectra of LMS and SSOs below $`T_{\mathrm{eff}}<\mathrm{\hspace{0.17em}2400}`$ K. In the near-infrared color ($`J`$-$`K`$) (Figure 13), grain condensation yields very red colors for the DUSTY models, whereas COND and grainless models loop back to short wavelengths (blueward) below $`m0.11M_{}`$, $`M_K10`$ for $`[\mathrm{M}/\mathrm{H}]`$ = -2, and $`m0.072M_{}`$ (resp. $`0.04M_{}`$) at 5 Gyr (resp. 100 Myr) for solar metallicity, i.e. $`M_K12`$. As described in §2.2, this stems from (i) the onset of H<sub>2</sub> CIA for low-metallicity objects and (ii) the formation of CH<sub>4</sub> at the expense of CO when the peak of the Planck function moves in the wavelength range characteristic of the absorption bands of this species, for solar-like abundances. In terms of colors there is a competing effect between grain and molecular opacity sources for objects at the bottom and just below the MS. The backwarming effect resulting from the large grain opacity destroys molecules such as e.g H<sub>2</sub>O, one of the main sources of absorption in the near-IR (cf. §2.2.1), and yields the severe reddening near the bottom of the MS ($`T_{\mathrm{eff}}<\mathrm{\hspace{0.17em}2300}2400`$ K), as illustrated by the DENIS (Tinney et al 1998) and 2MASS (Kirkpatrick et al 1999a) objects and by GD165B (Kirkpatrick et al 1999b). As the temperature decreases, grains condense and settle at deeper layers below the photosphere while methane absorption in the IR increases (see §2.2) so that the peak of the flux will move back to shorter wavelengths and the color sequence will go from the DUSTY one to the COND one. The exact temperature at which this occurs is still uncertain because it involves the complex grain thermochemistry and dynamics outlined in §2.2, but it should lie between $`T_{\mathrm{eff}}1000`$ K, for Gl229B, and $`1700`$ K, corresponding to $`J`$-$`K2`$ for the presently reddest observed L-dwarf with determined parallax LHS102B (Goldman et al 1999, Basri et al 2000). In spite of the afore-mentioned strong absorption in the IR, BDs around 1500 K radiate nearly 90% (99% with dust) of their energy at wavelengths longward of $`1\mu `$m and infrared colors are still preferred to optical colors (at least for solar metal abundance), with $`J,Z,H`$ as the favored bands for detection, and M$`{}_{M}{}^{}`$M$`{}_{L}{}^{}1011`$, M$`{}_{K}{}^{}`$M$`{}_{J}{}^{}1112`$ at the H-burning limit, at 1 Gyr (see Chabrier et al 2000). #### 4.5.3 Halo stars. Globular clusters Observations of LMS belonging to the stellar halo population of our Galaxy (also called ”spheroid” to differentiate it from the $`\rho (r)1/r^2`$ dark halo) or to globular clusters down to the bottom of the MS have been rendered possible with the HST optical (WFPC2) and IR (NICMOS) cameras and with parallax surveys at very faint magnitudes (Leggett 1992, Monet et al 1992, Dahn et al 1995). This provides stringent constraints for our understanding of LMS structure and evolution for metal-depleted objects. The reliability of the LMS theory outlined in §2 has been assessed by the successful confrontation to various observed sequences of globular clusters ranging from $`[\mathrm{M}/\mathrm{H}]2.0`$ to -1.0<sup>2</sup><sup>2</sup>2Globular clusters with an observed $`[Fe/H]`$ must be compared with theoretical models with the corresponding metallicity $`[M/H]`$ = $`[Fe/H]`$ \+ $`[O/Fe]`$, in order to take into account the enrichment of $`\alpha `$-elements (see Baraffe et al 1997) (Baraffe et al 1997). Figure 8 of Baraffe et al (1997) portrays several tracks for different metallicities, corresponding to the MS of various globular clusters. The tracks are superposed with the subdwarf sequence of Monet et al (1992), which is identified as stellar halo objects from their kinematic properties. As shown in this figure, the pronounced variations of the slope in the CMD around $`0.5M_{}`$ and $`0.2M_{}`$, which stem from the very physical properties of the stars, namely the onset of molecular absorption in the atmosphere and of degeneracy in the core (see §2), are well reproduced by the theory at the correct magnitudes and colors. The predicted blue loop in IR colors caused by the ongoing CIA absorption of H<sub>2</sub> for the lowest-mass (coolest) objects (see §2.2.1 and Figure 13 for $`[\mathrm{M}/\mathrm{H}]=2`$), has also been confirmed by observations with the NICMOS camera (Pulone et al 1998). ### 4.6 The lithium and deuterium tests One ironclad certification of a BD is a demonstration that hydrogen fusion has not occurred in its core. As Figure 2 illustrates, lithium burning through the Li$`{}_{}{}^{7}(p,\alpha )`$He<sup>4</sup> reaction occurs at a lower temperature than is required for hydrogen fusion. The timescale for the destruction of lithium in the lowest-mass stars is about 100 Myr. Moreover the evolutionary timescale in LMS/BDs is many orders of magnitude larger than the convective timescale (see e.g. Bildsten et al 1997), so core abundances and atmospheric abundances can be assumed identical. Therefore, the retention of lithium in a fully mixed object older than 10<sup>8</sup> yr signifies the lack of hydrogen burning. This provides the basis of the so-called ”lithium-test” first proposed by Rebolo et al (1992) to identify bona-fide BDs. Evolutionary models based on non-grey atmospheres and the screening factors described in §2.1.2 yield a lithium-burning minimum mass $`m_{Li}0.06M_{}`$ (Chabrier et al 1996). Note from Figure 2 the strong age dependence of the lithium-test : young stars at $`t<\mathrm{\hspace{0.17em}10}^8`$ yr (depending on the mass) will exhibit lithium (which in passing precludes the application of the lithium-test for the identification of SSOs in star forming regions), whereas massive BDs within the mass range \[0.06-0.07 $`M_{}`$\] older than $`10^8`$ yr will have burned lithium. By observing a young cluster older than $`10^8`$ yr, one can look for the boundary-luminosity below which lithium has not yet been depleted. Objects fainter than this limit will definitely be in the substellar domain. Conversely, the determination of the lithium depletion edge yields the age of the cluster, as was first proposed by Basri et al (1996) (see also e.g. Stauffer et al 1998). Indeed, the lithium test provides a “nuclear” age that may be even more powerful a diagnostic than the conventional nuclear age from the upper main sequence turnoff, since the evolution of the present fully convective low-mass objects does not depend on ill-constrained parameters such as mixing length or overshooting. As noted by Stauffer et al (1998) the age-scale for open clusters based on the (more reliable) lithium depletion boundary has important implications for stellar evolution: the ages of several clusters (Pleiades, $`\alpha `$-Per, IC 2391) are all consistent with a small but non-zero amount of overshooting to be included in the evolutionary models at the turn-off mass in order to yield similar ages. See Basri (this issue) for a more detailed discussion of the lithium-test. The deuterium-test can be used in a similar manner, extending the lithium test to smaller masses and younger ages, typically $`t<\mathrm{\hspace{0.17em}10}^7`$ yr (Béjar et al 1999). The $`D`$-burning minimum mass is predicted as $`m_D0.013M_{}`$ (see Figure 2) and has been proposed (Shu et al 1987) as playing a key role in the formation of isolated star-like objects, in contrast to objects formed in a protoplanetary disk. Deuterium-depletion is illustrated in Figure 14. This figure displays the evolution of SSOs above $`m_D`$ in the I-band. The left and right diagonal solid lines correspond to 50% D-depletion ($`[D]/[D]_0=1/2`$) and 99% D-depletion, respectively ($`[D]_0=2\times 10^5`$ is the initial mass fraction). The inset displays the corresponding curves in $`T_{\mathrm{eff}}`$. Spectroscopic signatures of deuterium include absorption lines of deuterated water HDO between 1.2 and 2.1 $`\mu `$m (Toth 1997). The observation of deuterated methane CH<sub>3</sub>D can not be used for the deuterium test as an age-indicator: at the temperature methane forms ($`<\mathrm{\hspace{0.17em}1800}`$ K), SSOs above the D-burning minimum mass are old enough for all deuterium to have been burned (see Figure 6). ### 4.7 Low-mass objects in cataclysmic variables The study of cataclysmic variables (CVs) is closely related to the theory of LMS and SSOs regarding their inner structure and the magnetic field generation. CVs are semi-detached binaries with a white dwarf (WD) primary and a low-mass stellar or sub-stellar companion ($`m<\mathrm{\hspace{0.17em}\hspace{0.17em}1}M_{}`$) that transfers mass to the WD through Roche-lobe overflow (see e.g. King 1988 for a review). In the current standard model, mass transfer is driven by angular momentum loss caused by magnetic wind braking for predominantly radiative stars and by gravitational wave emission for fully convective objects. The Roche lobe filling secondary’s mean density $`\overline{\rho }`$ determines almost entirely the orbital period $`P`$, with $`P_h=k/\overline{\rho }^{1/2}`$ (where $`k8.85`$ is a weak function of the mass ratio, $`P_h`$ is the orbital period in hr, and $`\overline{\rho }=(m_2/R_2^3)/(M_{}/R_{}^3)`$ is the mean density of the secondary in solar units) (see King 1988). For objects obeying classical mechanics, i.e. in the stellar regime, $`mR`$ (see §3.1 and Figure 3) so that the mean density increases with decreasing mass ($`\overline{\rho }1/m^2`$) and the orbital period decreases along evolution as the secondary loses mass $`P_h1/\overline{\rho }^{1/2}m_2`$. The situation reverses near and below the HBMM, when electron degeneracy dominates. In that case, the radius increases very slightly with decreasing mass (Figure 3), so that the mean density is essentially proportional to the mass ($`\overline{\rho }m`$) and the orbital period increases along evolution $`P_hm_2^{1/2}`$. Then the secular evolution of the orbital period reverses when the donor becomes a BD. The analysis of CV orbital evolution thus provides important constraints on LMS and BD internal structure. The major puzzle of CVs is the observed distribution of orbital periods, namely the dearth of systems in the 2-3 h period range and the minimum period $`P_{\mathrm{min}}`$ at 80 min (see King 1988). The most popular explanation for the period gap is the disrupted magnetic braking scenario (Rappaport et al 1983, Spruit & Ritter 1983), and most evolutionary models including this process do reproduce the observed period distribution. However, the recent progress realized in the field of LMS and SSOs now allows the confrontation of the theoretical predictions with the observed atmospheric properties (colors, spectral types) of the secondaries (Beuermann et al 1998, Clemens et al 1998, Kolb & Baraffe 1999a). The afore-mentioned standard period-gap model may be in conflict with the observed spectral type of some CV secondaries (Beuermann et al 1998, Kolb & Baraffe 1999b). Although alternative explanations exist for the period gap, none of them has been proved as successful as the disrupted magnetic braking scenario. The most recent alternative suggestion, based on a characteristic feature of the mass-radius relationship of LMS (Clemens et al 1998), has been shown to fail to reproduce the observed period distribution around the period gap (Kolb et al 1998). On the other hand, as discussed in §2.3, no change in the level of activity is observed in isolated M-dwarfs along the transition between partially and fully convective structures at $`m0.35M_{}`$ and $`Sp`$ M2-M4, which indicates that a magnetic field is still generated in the mass-range corresponding to the CV-gap. In this case, the abrupt decrease of angular momentum losses by magnetic braking at the upper edge of the period gap could result from a rearrangement of the magnetic field when stars become fully convective, without necessarily implying a sudden decline in the magnetic activity, as suggested by Taam & Spruit (1989, see also Spruit 1994). Finally, a $``$ 10% discrepancy still remains between the theoretical and the observed minimum period $`P_{\mathrm{min}}=`$80 min, even when including the most recent improvements in stellar physics (see Kolb & Baraffe 1999a for details). Residual shortcomings in the theory, either in the EOS or in the atmosphere, cannot be ruled out as the cause of this discrepancy. Alternatively, an additional driving mechanism to gravitational radiation in fully convective objects can reconcile predicted and observed $`P_{\mathrm{min}}`$ (Kolb & Baraffe 1999a). Because magnetic activity is still observed in fully convective late spectral type M-dwarfs, magnetic braking could operate in CV secondaries even down to $`P_{\mathrm{min}}`$, but with a weaker efficiency than above the period gap. These open questions certainly require a better understanding of magnetic field generation and dissipation in LMS and BDs, a point already stressed in §2.3. ## 5 GALACTIC IMPLICATIONS ### 5.1 Stellar Luminosity Function and Mass Function It is now well established that visible stars are not numerous enough to account for the dynamics of our Galaxy, the so-called galactic dark matter problem. A precise determination of their density requires the correct knowledge of the luminosity function (LF) down to the H-burning limit, and a correct transformation into a mass-function (MF). The latter issue has been improved significantly with models that describe more accurately the color-color and color-magnitude diagrams of LMS and BDs for various metallicities, and most importantly that provide mass-magnitude relationships in good agreement with the observations (see §4). The former issue, however, is not completely settled at present and significant differences still exist between various determinations of the nearby LF and between the nearby LF and the photometric LF determined from the ground and with the HST, as shown later in this section. The main problem with the nearby sample is the limited number of stars at faint magnitudes. Although the sample is complete to 20 pc for stars with $`M_V<9.5`$, it is severely incomplete beyond 5 pc for $`M_V>12`$ (Henry et al 1997). Kroupa (unpublished) derived a nearby LF $`\mathrm{\Phi }_{near}`$ by combining Hipparcos parallax data, which is essentialy complete for $`M_V<12`$ at r=10 pc, and the sample of nearby stars (Dahn et al 1986) with ground-based parallaxes for $`V>12`$ to a completeness distance r=5.2 pc. The nearby LF determined by Reid and Gizis (1997) is based on a volume sample within 8 pc. Most of the stars in this survey have parallaxes. For all the late K and M dwarfs, however, trigonometric parallaxes are not available and these authors use a spectroscopic TiO-index vs $`M_V^{TiO}`$ relation to estimate the distance (Reid et al 1995). This sample was revised recently with Hipparcos measurements and new binary detections in the solar neighborhood (Delfosse et al 1999a) and leads to a revised northern 8-pc catalogue and nearby LF (Reid et al 1999). These authors argue that their sample should be essentially complete for $`M_V<14`$. However, the analysis of completeness limits by Henry et al (1997, their figure 1) shows that the known stellar census becomes substantially incomplete for distances larger than 5 pc. About 35% of the systems in the Reid et al sample are multiple and $`45\%`$ of all stars have a companion in binary or multiple systems. The photometric LFs, $`\mathrm{\Phi }_{phot}`$, based on large observed volumes via deep pencil-beam surveys, avoid the limitation due to small statistics but introduce other problems such as Malmquist bias and unresolved binaries (see e.g. Kroupa et al 1993, Kroupa 1995), yielding only the determination of the stellar system LF. For the HST LF (Gould et al 1998), with $`I<\mathrm{\hspace{0.17em}24}`$, the Malmquist bias is negligible because all stars down to $`0.1M_{}`$ are seen through to the edge of the thick disk. A major caveat of photometric LFs, however, is that the determination of the distance relies on a photometric determination from a color-magnitude diagram. In principle this requires the determination of the metallicity of the stars, since colors depend on metallicity (see Figure 12). The illustration of the disagreement between the three afore-mentioned LFs, $`\mathrm{\Phi }_{5pc}`$, $`\mathrm{\Phi }_{8pc}`$ and $`\mathrm{\Phi }_{HST}`$, at faint magnitude is apparent in Figure 15. The disagreement between the two nearby LFs, in particular, presently has no robust explanation. It might come from incompleteness of the 8 pc sample at faint magnitudes. Note also that in the spectroscopic relation used to estimate the distance, $`M_V`$ can be uncertain by $`1`$ mag (see Figure 3 of Reid et al 1995), certainly a source of Malmquist bias, even though the number of stars without parallax in this sample is small ($`10`$%) and a Malmquist bias on this part should not drastically affect the results. On the other hand, as noted in §4.1, the end of the 5-pc sample might be contaminated by a statistically significant number of young BDs, or by stars with slightly under-solar abundances. Figure 16 compares the MFs obtained from the two afore-mentioned nearby LFs, using the mass-magnitude relations described in §4.2. In practice, the observed sample constitutes a mixture of ages and metallicities. However the Hipparcos CMD indicates that $`90\%`$ of disk stars have $`0.3<[\mathrm{M}/\mathrm{H}]<0.1`$ (Reid 1999) so the spread of metallicity is unlikely to significantly affect the derivation of the MF through the mass-magnitude relation. Moreover the choice of $`m`$-M<sub>K</sub> minimizes metallicity effects (see Figure 9 and figure 3 of Baraffe et al 1998). The MS lifetimes for stars with $`m<\mathrm{\hspace{0.17em}1}M_{}`$ are longer than a Hubble time and age variations in the mass-magnitude relation will affect only objects $`<\mathrm{\hspace{0.17em}0.15}M_{}`$ younger than $`0.5`$ Gyr (see Figure 9). Superposed is a power-law MF normalized at 0.8 $`M_{}`$ on the Hipparcos sample $`dn/dm=0.031(m/0.8)^\alpha =0.024m^\alpha M_{}^1pc^3`$ with $`\alpha =1.2`$. As shown in Figure 16, this provides a very reasonable description of the low-mass part of the MF, within a variation of $`10\%`$ for $`\alpha `$, thus confirming the previous analysis of Kroupa et al (1993), except for the very last bin obtained from Kroupa’s LF, which predits about twice as many stars at 0.1 $`M_{}`$ (out of the figure). Similar results are obtained from the bolometric and $`M_K`$\- LFs for the 8-pc sample. Integration of this MF yields very-low-mass stars ($`m0.8M_{}`$) number- and mass-densities $`n_{VLMS}7.0\times 10^2`$ pc<sup>-3</sup> and $`\rho _{VLMS}2.0\times 10^2`$ $`\mathrm{M}_{}\mathrm{pc}^3`$, respectively. Adding up the contribution from more massive stars, $`\rho _{>0.8}1.4\times 10^2`$ $`\mathrm{M}_{}\mathrm{pc}^3`$ (Miller & Scalo 1978) and stellar remnants, $`\rho _{WD+NS}3.0\times 10^3`$ $`\mathrm{M}_{}\mathrm{pc}^3`$ yields a stellar density $`\rho _{}3.4\times 10^2`$ $`\mathrm{M}_{}\mathrm{pc}^3`$ in the Galactic disk, i.e. a surface density $`\mathrm{\Sigma }_{}22\pm 2`$ $`\mathrm{M}_{}\mathrm{pc}^2`$ (assuming a scale height $`h=320`$ pc). A Salpeter MF, $`dn/dmm^{2.35}`$, all the way down to the bottom of the MS would overestimate the density at 0.1 $`M_{}`$ by more than a factor of 10. For the Galactic spheroid the question is more settled. The LFs determined from nearby surveys (Dahn et al 1995) and from the HST (Gould et al 1998) are comparable within about 2 sigmas, and yield a MF with $`\alpha <\mathrm{\hspace{0.17em}1}`$ (Graff & Freese 1996, Chabrier & Méra 1997, Gould et al 1998) and a stellar number- and mass-density $`n_{}<10^3pc^3`$ and $`\rho _{}<4.0\times 10^5`$ $`\mathrm{M}_{}\mathrm{pc}^3`$, i.e. an optical depth $`\tau 10^9`$, about 1% of the value measured toward the LMC. At present, the dark halo MF is completely unknown but both the Hubble Deep Field observations and the narrow-range of the observed time-distribution of the microlensing events toward the LMC strongly suggest an IMF different from a Salpeter IMF below $`1M_{}`$ (Chabrier 1999). Assuming a homogeneous distribution, both constraints yield a mass-density for LMS in the dark halo about 2 orders of magnitude smaller than the afore-mentioned spheroid one (Chabrier & Méra 1997). This means that essentially no stars formed below about a solar mass in the dark halo. ### 5.2 Brown dwarf mass function Even though the possibility that BDs could make up for the Galactic missing mass is now clearly excluded, a proper census of the number of BDs has significant implications for our understanding of how stars and planets form. The determination of the BD MF is a complicated task. By definition, BDs never reach thermal equilibrium and most of the BDs formed at the early stages of the Galaxy will have dimmed to very low-luminosities ($`L1/t`$). Thus observations will be biased toward young and massive BDs. This age-undetermination is circumvented when studying the BD MF in clusters because objects in this case are likely to be coeval. The Pleiades cluster has been extensively surveyed and several BDs have been identified down to $`0.04M_{}`$ (Martín et al 1998, Bouvier et al 1998, Hambly et al 1999). A single power-law function from $`0.4`$ to $`0.04M_{}`$ seems to adequately reproduce the observations (not corrected for binaries) with some remaining uncertainties in the exponent: $`\alpha 0.61.0`$. However, stellar objects in very young clusters ($`<\mathrm{\hspace{0.17em}10}^6`$ yr) might still be accreting whereas older clusters may have already experienced significant dynamical evolution, mass segregation in the core and/or evaporation in the outer regions (see e.g. Raboud & Mermilliod 1998). In this case the present-day mass function does not reflect the initial mass function. For these reasons, it is probably premature to claim a robust determination of the MF in young clusters. The DENIS (Delfosse et al 1999b) and 2MASS (Kirkpatrick et al 1999b; Burgasser et al 1999) surveys, which have covered an area of several hundred square degrees and are complete to $`K=13.5`$ and 14.5, respectively, have revealed about 20 field L-dwarfs and 4 field methane-BDs. This yields an L-dwarf number-density $`n_L0.03\pm 0.01`$ sq deg<sup>-1</sup> for $`K<14.5`$ and a methane-BD number-density $`n_{CH_4}0.002`$ sq deg<sup>-1</sup> for $`J<16`$. Although these numbers correspond to small statistics and should be considered with caution, they provide the first observational constraints on the substellar MF in the Galactic disk and are consistent with a slowly rising BD MF with $`\alpha 1`$-2 (Reid et al 1999; Chabrier in preparation), assuming a constant formation rate. Independent, complementary information on the stellar and substellar MF comes from microlensing observations. Indeed, the time distribution of the events provides a (model-dependent) determination of the mass-distribution and thus of the minimum mass of the dark objects: $`dN_{ev}/dt_e=E\times ϵ(t_e)\times d\mathrm{\Gamma }/dt_eP(m)/\sqrt{m}`$, where $`E`$ is the observed exposure, i.e. the number of star$`\times `$years, $`ϵ`$ is the experimental efficiency, $`\mathrm{\Gamma }`$ is the event rate and $`P(m)`$ is the mass probability distribution. The analysis of the published 40 MACHO (Alcock et al 1997) + 9 OGLE (Udalsky et al 1994) events toward the bulge is consistent with a rising MF at the bottom of the MS whereas a decreasing MF below 0.2 $`M_{}`$ seems to be excluded at a high ($`>90\%`$) confidence level (Han & Gould 1996, Méra et al 1998). Although the time distribution might be affected by various biases (e.g. blending) and robust conclusions must wait for larger statistics, the present results suggest that, in order to explain both star counts and the microlensing experiments, a substantial number of BDs must be present in the Galactic disk. Extrapolation of the stellar MF determined in §5.1 into the BD domain down to 0.01 $`M_{}`$ yields for the Galactic disk a BD number-density comparable to the stellar one, $`n_{BD}0.1pc^3n_{}`$, and a mass-density $`\rho _{BD}3.0\times 10^3\mathrm{M}_{}\mathrm{pc}^3`$, i.e. $`\mathrm{\Sigma }_{BD}2\mathrm{M}_{}\mathrm{pc}^2`$. For the spheroid, extrapolation of the previously-determined stellar MF yields $`n_{BD}<\mathrm{\hspace{0.17em}10}^4pc^3`$, $`\rho _{BD}<10^5\mathrm{M}_{}\mathrm{pc}^3`$, less than 0.1% of the required dynamical density. For the dark halo the density is about 2 orders of magnitude smaller, as mentioned before. It is obviously premature to try to infer the mass distribution of exoplanets. This will first require a clear theoretical and observational distinction between planets and brown dwarfs. However, an interesting preliminary result comes from the observed mass distribution of the companions of $`G`$ and $`K`$ stars. As shown in figure 4 of Mayor et al (1998), there is a strong discontinuity in the mass distribution at $`m_2/\mathrm{sin}i5M_J`$, with a clear peak below this limit. This suggests that planet formation in a protoplanetary disk is a much more efficient mechanism than BD formation as star companions, at least around $`G`$ and $`K`$ stars (see e.g. Marcy & Butler 1998). It also suggests that the MF for BD stellar companions differs from the MF in the field. Ongoing observations around M-stars will tell us whether such a mass distribution still holds around low-mass stars. ## 6 Conclusions This review has summarized the significant progress achieved within the past few years in the theory of cool and dense objects at the bottom of and beyond the main sequence: low-mass stars, brown dwarfs and gaseous planets. The successful confrontation of the theory with the numerous detections of low-mass stellar and sub-stellar objects allows a better understanding of their structural and thermal properties, and allows reliable predictions about their evolution. This in turn brings confidence in the predicted characteristic properties of these objects, a major issue in terms of search strategies for future surveys. Important problems remain to be solved to improve the theory. A non-exhaustive list includes for example: (i) a better determination of the EOS in the pressure-ionization region, with the possibility of a first-order phase transition, (ii) the study of phase separation of elements in SSO interiors, (iii) an improved treatment of convection in optically-thin regions, (iv) a precise description of the dynamics of grain formation and sedimentation in SSO atmospheres, (v) the derivation of an accurate mass-$`T_{\mathrm{eff}}`$-age scale for SSOs and young objects, (vi) a correct understanding of magnetic-field generation and dissipation in active LMS and BDs. The increasing number of observed LMS and SSOs, together with the derivation of accurate models, eventually will allow a robust determination of the stellar and substellar mass functions, of the minimum mass for the formation of star-like objects, and thus of the exact density of these objects in the Galaxy. As we discussed in §5, present determinations in various Galactic regions point to a slowly rising MF near and below the H-burning limit, with a BD number-density comparable to the stellar one, and a MF truncated below $`1M_{}`$ in the dark halo. A more precise determination must await confirmation from future observations. At last, the amazingly rapid pace of exoplanet discoveries should yield the determination of the planetary MF and maximum mass, and eventually the direct detection of such objects, a future formidable test for the theory. Acnowledments: This review has benefited from various discussions with our colleagues F Allard, PH Hauschildt, D Alexander, K Lodders, A Burrows, J Lunine, B Fegley, T Forveille, X Delfosse, J Bouvier, P Kroupa, A Nordlund, T Guillot, D Saumon, U Kolb. Our profound gratitude goes to these individuals, who helped by improving the original manuscript. 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warning/0006/math0006104.html	ar5iv	text	# 1 Introduction ## 1 Introduction Vertex operator algebras, introduced in mathematics (\[B\], \[FLM\]), are known essentially to be chiral algebras, introduced in physics (\[BPZ\], \[MS\]). In terms of physical language, chiral algebras for bosonic field theories are vertex operator algebras while chiral algebras are vertex operator superalgebras for fermionic field theories. In physics, further generalizations of bosons and fermions are parafermions (\[ZF1-2\], \[G\]), where the chiral algebras were called parafermion algebras. Independently (and earlier), $`Z`$-operators and $`Z`$-algebras (\[LW1-6\], \[LP1-2\]) were introduced in mathematics to study standard modules for affine Lie algebras. In \[DL1-3\] and \[LP1-2\], the relations between $`Z`$-operators and parafermion operators were clarified. Furthermore, in \[DL2\], Lepowsky-Wilson’s $`Z`$-algebras were put into larger, more natural algebras where the notions of generalized vertex (operator) algebra and abelian intertwining algebras were introduced in \[DL2\]. Generalized algebraic structures associated with rational lattices were also studied in \[M\] and a notion called vertex operator para-algebra was independently introduced and studied in \[FFR\]. Roughly speaking, parafermion algebras are generalized vertex operator algebras. For bosonic or fermionic field operators $`a(z)`$ and $`b(z)`$, the locality amounts to $`(z_1z_2)^ka(z_1)b(z_2)=(1)^{\|a\|\|b\|}(z_1z_2)^kb(z_2)a(z_1)`$ (1.1) for some non-negative integer $`k`$, where $`\|a\|=0`$ if $`a(z)`$ is bosonic and $`1`$ if $`a(z)`$ is fermionic. Parafermions are always associated to an abelian group $`G`$ equipped with a $`𝐂^\times `$-valued alternating $`𝐙`$-bilinear form $`c(,)`$ and a $`𝐂/\mathrm{𝟐}𝐙`$-valued $`𝐙`$-bilinear form $`(,)`$ on $`G`$. For parafermion operators $`a(z),b(z)`$ with gradings $`g,hG`$, the following relation holds: $`(z_1z_2)^{k+(g,h)}a(z_1)b(z_2)=(1)^kc(g,h)(z_2z_1)^{k+(g,h)}b(z_2)a(z_1)`$ (1.2) for some nonnegative integer $`k`$. In physical literatures, a chiral algebra is often described by a set of generating field operators and a certain set of relations (such as operator product expansions). Now we know that a certain set of field operators on a vector space indeed gives rise to a vertex operator (super)algebra. A result proved in \[Li2\] is that any set of mutually local vertex operators on a vector space $`W`$ generates a canonical vertex superalgebra with $`W`$ as a module. (This result was extended in \[Li3\] for twisted modules.) This is an analogue of the simple fact in linear algebra that any set of mutually commutative endomorphisms on a vector space $`U`$ generates a commutative associative algebra with $`U`$ as a module. (See \[FKRW\], \[LZ\], \[MN\], \[MP\] and \[X\] for other related interesting results.) As the main result of this paper, we extend the results of \[Li2-3\] for the notion of generalized vertex algebra \[DL2\]. This paper is modeled on \[Li2\], however, the two key theorems (Propositions 3.8 and 3.13) require essentially new proofs. It seems that the most general and natural notion is the one of abelian intertwining algebra \[DL2\]. An abelian intertwining algebra by definition is associated to an abelian group $`G`$ and $`𝐂^\times `$-valued functions $`F`$ on $`G\times G\times G`$ and $`\mathrm{\Omega }`$ on $`G\times G`$ satisfying certain conditions. It would be nice to extend our result to the notion of abelian intertwining algebra. When we were trying, we found that the extension is almost straightforward except for extending the two key theorems (Propositions 3.8 and 3.13) we need certain identities purely about $`F`$ and $`\mathrm{\Omega }`$. We can prove that one of the identities follows from the assumptions on $`F`$ and $`\mathrm{\Omega }`$, but we are not be able to prove the others. We hope to discuss this issue in some other place. This paper consists of four sections including this introduction as the first section. In Section 2, we recall some basic definitions and results from \[DL2\]. The main result of this paper is given in Section 3. In Section 4, we construct canonical generalized vertex algebras from $`Z`$-algebras of any nonzero level. ## 2 Definition and duality for generalized vertex algebras This section is preliminary. In this section, we recall from \[DL2\] the basic definitions (of generalized vertex algebra and module) and basic duality properties. First, let us briefly review some formal variable calculus. (Best references are \[FLM\] and \[FHL\].) Throughout this paper, $`z,z_0,z_1,z_2`$ and $`x,y`$ will be mutually commuting (independent) formal variables. We shall use $`N`$ for the set of all nonnegative integers, $`𝐙_+`$ for the set of positive integers and $`𝐂`$ for the set of complex numbers. All vector spaces are assumed to be over $`𝐂`$. For a vector space $`U`$, set $`U\{z\}=\left\{{\displaystyle \underset{n𝐂}{}}u(n)z^n\right\|u(n)U\text{ for }n𝐂\}.`$ (2.1) Let $`D=d/dz`$ be the formal differential operator on $`U\{z\}`$: $`D\left({\displaystyle \underset{n𝐂}{}}u(n)z^n\right)={\displaystyle \underset{n𝐂}{}}nu(n)z^{n1}.`$ (2.2) The formal residue operator $`\mathrm{Res}_z`$ from $`U\{z\}`$ to $`U`$ is defined by $`\mathrm{Res}_zu(z)=u(1)`$ (2.3) for $`u(z)=_{n𝐂}u(n)z^nU\{z\}`$. The following are useful subspaces of $`U\{z\}`$: $`U[[z,z^1]]`$ $`=`$ $`\left\{{\displaystyle \underset{n𝐙}{}}u(n)z^n\right\|u(n)U\text{ for }n𝐙\},`$ (2.4) $`U((z))`$ $`=`$ $`\left\{{\displaystyle \underset{n𝐙}{}}u(n)z^nU[[z,z^1]]\right\|u(n)=0\text{ for }n\text{sufficiently small}\},`$ (2.5) $`U[[z]]`$ $`=`$ $`\left\{{\displaystyle \underset{n𝐙}{}}u(n)z^nU[[z,z^1]]\right\|u(n)=0\text{ for }n<0\}.`$ (2.6) A typical element of $`𝐂[[𝐳,𝐳^\mathrm{𝟏}]]`$ is the formal Fourier expansion of the delta-function at $`0`$: $`\delta (z)={\displaystyle \underset{n𝐙}{}}z^n.`$ (2.7) Its fundamental property is: $`f(z)\delta (z)=f(1)\delta (z)\text{ for }f(z)𝐂[𝐳,𝐳^\mathrm{𝟏}].`$ (2.8) For $`\alpha 𝐂`$, by definition, $`(z_1z_2)^\alpha ={\displaystyle \underset{i0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{\alpha }{i}}\right)(1)^iz_1^{\alpha i}z_2^i.`$ (2.9) Then $`\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)={\displaystyle \underset{n𝐙}{}}\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^n={\displaystyle \underset{n𝐙}{}}{\displaystyle \underset{i0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)(1)^iz_0^nz_1^{ni}z_2^i.`$ (2.10) We have the following fundamental properties of delta function (\[FLM\], \[FHL\], \[Le\], \[Zhu\]): ###### Lemma 2.1 For $`\alpha 𝐂`$, $`z_0^1\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^\alpha \delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)=z_1^1\left({\displaystyle \frac{z_0+z_2}{z_1}}\right)^\alpha \delta \left({\displaystyle \frac{z_0+z_2}{z_2}}\right);`$ (2.11) For $`r,s,k𝐙`$ and for $`p(z_1,z_2)𝐂[[𝐳_\mathrm{𝟏},𝐳_\mathrm{𝟐}]]`$, $`z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)z_1^rz_2^s(z_1z_2)^kp(z_1,z_2)z_0^1\delta \left({\displaystyle \frac{z_2z_0}{z_0}}\right)z_1^rz_2^s(z_2+z_1)^kp(z_1,z_2)`$ (2.12) $`=`$ $`z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)(z_2+z_0)^rz_2^sz_0^kp(z_2+z_0,z_2).`$ In particular, $`z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)z_0^1\delta \left({\displaystyle \frac{z_2z_0}{z_0}}\right)=z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right).`$ (2.13) Note that (2.11) is equivalent to $`(z_1z_2)^\alpha z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)=z_1^1z_0^\alpha \left({\displaystyle \frac{z_0+z_2}{z_1}}\right)^\alpha \delta \left({\displaystyle \frac{z_0+z_2}{z_2}}\right)\text{ for }\alpha 𝐂.`$ (2.14) A generalized vertex algebra by definition is associated to an abelian group $`G`$, a symmetric $`𝐂/\mathrm{𝟐}𝐙`$-valued $`𝐙`$-bilinear form (not necessarily nondegenerate) on $`G`$: $`(g,h)𝐂/\mathrm{𝟐}𝐙\text{ for }𝐠,𝐡𝐆`$ (2.15) and $`c(,)`$ is a $`𝐂^\times `$-valued alternating $`𝐙`$-bilinear form on $`G`$. A generalized vertex algebra associated with the group $`G`$ and the forms $`(,)`$ and $`c(,)`$ is a $`G`$-graded vector space $`V={\displaystyle \underset{gG}{}}V^g,`$ (2.16) equipped with a linear map $`Y:`$ $`V(\mathrm{End}V)\{z\}`$ (2.17) $`uY(u,z)={\displaystyle \underset{n𝐂}{}}u_nz^{n1}`$ and with a distinguished vector $`\mathrm{𝟏}V^0`$, called the vacuum vector, satisfying the following conditions for $`g,hG,u,vV`$ and $`l𝐂`$: $`u_lV^hV^{g+h}\text{ if }uV^g;`$ (2.18) $`u_lv=0\text{ if the real part of }l\text{ is sufficiently large};`$ (2.19) $`Y(\mathrm{𝟏},z)=1;`$ (2.20) $`Y(v,z)\mathrm{𝟏}V[[z]]\text{ and }v_1\mathrm{𝟏}(=\underset{z0}{lim}Y(v,z)\mathrm{𝟏})=v;`$ (2.21) $`Y(u,z)\|_{V^h}={\displaystyle \underset{n(g,h)\mathrm{mod}𝐙}{}}v_nz^{n1}\text{ if }uV^g`$ (2.22) (i.e., $`n+2𝐙(g,h)\mathrm{mod}𝐙/2𝐙`$); $`z_0^1\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y(u,z_1)Y(v,z_2)`$ (2.23) $`c(g,h)z_0^1\left({\displaystyle \frac{z_2z_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)Y(v,z_2)Y(u,z_1)`$ $`=`$ $`z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)Y(Y(u,z_0)v,z_2)\left({\displaystyle \frac{z_1z_0}{z_2}}\right)^g`$ (the generalized Jacobi identity) if $`uV^g,vV^h`$, where $`\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)\left({\displaystyle \frac{z_1z_0}{z_2}}\right)^gw=\left({\displaystyle \frac{z_1z_0}{z_2}}\right)^{(g,g^{})}\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)w`$ (2.24) for $`wV^g^{},g^{}G`$. This completes the definition. The map $`Y`$ is called the vertex operator map. The generalized vertex algebra is denoted by $$(V,Y,\mathrm{𝟏},G,c(,),(,))$$ or briefly, by $`V`$. ###### Remark 2.2 Note that we here slightly generalize the original definition in \[DL2\] where $`(,)`$ was assumed to be $`(\frac{1}{T}𝐙)/\mathrm{𝟐}𝐙`$-valued, where $`T`$ is a positive integer called the level. The main reason for this generalization is to include the generalized vertex algebras associated to affine algebras with a non-rational level. On the other hand, if $`G`$ is finite, there exists a positive integer $`T`$ such that $`(,)`$ ranges in $`\frac{1}{T}𝐙/2𝐙`$. We recall the following remarks from \[DL2\]: ###### Remark 2.3 If $`G=0`$, the notion of generalized vertex algebra reduces to the notion of vertex algebra. If $`G=𝐙/\mathrm{𝟐}𝐙`$ with $`(m+𝐙,𝐧+𝐙)=\mathrm{𝐦𝐧}+𝐙`$, the notion of generalized vertex algebra reduces to the notion of vertex superalgebra, noting that $`c(,)=1`$. ###### Remark 2.4 A generalized vertex operator algebra is a generalized vertex algebra $`V`$ associated to a finite group $`G`$ with $`c(,)=1`$ and $`(,)`$ being nondegenerate and furthermore, it is equipped with another distinguished vector $`\omega V_2^0`$, called the Virasoro vector, such that $`[L(m),L(n)]=(mn)L(m+n)+{\displaystyle \frac{1}{12}}(m^3m)\delta _{m+n,0}c\text{ for }m,n𝐙,`$ (2.25) $`[L(1),Y(v,z)]={\displaystyle \frac{d}{dz}}Y(v,z)\text{ for }vV,`$ (2.26) where $`Y(\omega ,z)=_{m𝐙}L(m)z^{m2}`$ and $`c`$ is a complex number, called the rank of $`V`$, and such that $`V={\displaystyle \underset{n𝐂}{}}V_n\text{ where }V_n=\{vV\|L(0)v=nv\}\text{ for }n𝐂;`$ (2.27) $`V^g={\displaystyle \underset{n𝐂}{}}V_n^g\text{(where }V_n^g=V_nV^g\text{)}\text{ for }gG;`$ (2.28) $`dimV_n<\mathrm{}\text{ for }n𝐂,`$ (2.29) $`V_n=0\text{ for }n\text{ whose real part is sufficiently small}.`$ (2.30) ###### Proposition 2.5 \[DL2\] In the presence of all the axioms except the generalized Jacobi identity in defining the notion of generalized vertex algebra, the generalized Jacobi identity is equivalent to the following generalized weak commutativity and associativity: (A) For $`g_1,g_2G`$ and $`v_1V^{g_1},v_2V^{g_2}`$, there exists a nonnegative integer $`k`$ such that $`(z_1z_2)^{k+(g_1,g_2)}Y(v_1,z_1)Y(v_2,z_2)`$ (2.31) $`=`$ $`(1)^kc(g_1,g_2)(z_2z_1)^{k+(g_1,g_2)}Y(v_2,z_2)Y(v_1,z_1).`$ (B) For $`g_1,g_2,hG`$ and $`v_1V^{g_1},v_2V^{g_2},wV^h`$, there exists a nonnegative integer $`l`$ such that $`(z_0+z_2)^{l+(g_1,h)}Y(v_1,z_0+z_2)Y(v_2,z_2)w=(z_2+z_0)^{l+(g_1,h)}Y(Y(v_1,z_0)v_2,z_2)w,`$ (2.32) where $`l`$ is independent of $`v_2`$. The following result was also due to \[DL2\]: ###### Proposition 2.6 In the presence of all the axioms except the generalized Jacobi identity in defining the notion of generalized vertex algebra, the generalized Jacobi identity follows from the generalized weak commutativity and the following property $`[D,Y(v,z)]={\displaystyle \frac{d}{dz}}Y(v,z)\text{ for }vV,`$ (2.33) where $`D`$ is an endomorphism of $`V`$ defined by $`D(v)=v_2\mathrm{𝟏}\text{ for }vV.`$ (2.34) Proof. We here give a slightly different proof by generalizing the proof given in \[Li2\] for the corresponding result for vertex superalgebras. First, just as in \[Li2\], from the vacuum property and (2.33) we have $`Y(v,z)\mathrm{𝟏}=e^{zD}v\text{ for }vV.`$ (2.35) Second, (2.33) is equivalent to the following conjugation formula: $`e^{z_0D}Y(v,z)e^{z_0D}=Y(v,z+z_0).`$ (2.36) Third, we shall derive a skew-symmetry. Let $`uV^g,vV^h,g,hG`$. Then there exists a nonnegative integer $`k`$ such that $`(z_1z_2)^{k+(g,h)}Y(u,z_1)Y(v,z_2)=(1)^kc(g,h)(z_2z_1)^{k+(g,h)}Y(v,z_2)Y(u,z_1)`$ (2.37) and such that $`z^{k+(g,h)}Y(v,z)uV[[z]].`$ (2.38) Then $`(z_1z_2)^{k+(g,h)}Y(u,z_1)Y(v,z_2)\mathrm{𝟏}`$ (2.39) $`=`$ $`(1)^kc(g,h)(z_2z_1)^{k+(g,h)}Y(v,z_2)Y(u,z_1)\mathrm{𝟏}`$ $`=`$ $`(1)^kc(g,h)(z_2z_1)^{k+(g,h)}Y(v,z_2)e^{z_1D}u`$ $`=`$ $`(1)^kc(g,h)e^{z_1D}(z_2z_1)^{k+(g,h)}Y(v,z_2z_1)u`$ $`=`$ $`(1)^kc(g,h)e^{z_1D}(e^{\pi i}z_1+z_2)^{k+(g,h)}Y(v,e^{\pi i}z_1+z_2)u.`$ We are using (2.38). Now, it is safe for us to replace $`z_2`$ with $`0`$. In this way we get $`z_1^{k+(g,h)}Y(u,z_1)v=c(g,h)e^{\pi i(g,h)}z_1^{k+(g,h)}e^{z_1D}Y(v,e^{\pi i}z_1)u.`$ (2.40) Then we obtain the following skew-symmetry: $`Y(u,z_1)v=c(g,h)e^{\pi i(g,h)}e^{z_1D}Y(v,e^{\pi i}z_1)u.`$ (2.41) Next, we prove the generalized weak associativity. Let $`uV^{g_1},vV^{g_2},wV^{g_3}`$. Let $`k`$ be a nonnegative integer (only depending on $`u,w`$) such that $`(z_1z_2)^{k+(g_1,g_3)}Y(u,z_1)Y(w,z_2)=(1)^kc(g_1,g_3)(z_2z_1)^{k+(g_1,g_3)}Y(w,z_2)Y(u,z_1).`$ (2.42) Then using the skew-symmetry (2.41) and the conjugation formula (2.36) we obtain the following generalized associativity relation $`(z_0+z_2)^{k+(g_1,g_3)}Y(u,z_0+z_2)Y(v,z_2)w`$ (2.43) $`=`$ $`c(g_2,g_3)e^{\pi i(g_2,g_3)}(z_0+z_2)^{k+(g_1,g_3)}Y(u,z_0+z_2)e^{z_2D}Y(w,e^{\pi i}z_2)v`$ $`=`$ $`c(g_2,g_3)e^{\pi i(g_2,g_3)}(z_0+z_2)^{k+(g_1,g_3)}e^{z_2D}Y(u,z_0)Y(w,e^{\pi i}z_2)v`$ $`=`$ $`(1)^kc(g_2,g_3)c(g_1,g_3)e^{\pi i(g_2,g_3)}(e^{\pi i}z_2z_0)^{k+(g_1,g_3)}e^{z_2D}Y(w,e^{\pi i}z_2)Y(u,z_0)v`$ $`=`$ $`(1)^ke^{\pi i(g_2,g_3)}(e^{\pi i}z_2z_0)^{k+(g_1,g_3)}c(g_1+g_2,g_3)e^{z_2D}Y(w,e^{\pi i}z_2)Y(u,z_0)v`$ $`=`$ $`(z_2+z_0)^{k+(g_1,g_3)}Y(Y(u,z_0)v,z_2)w.`$ Then it follows from Proposition 2.5. $`\mathrm{}`$ A $`V`$-module \[DL2\] is a vector space $`W=_{sS}W^s`$, where $`S`$ is a $`G`$-set equipped with a $`𝐂/\mathrm{𝟐}𝐙`$-valued function $`(,)`$ on $`G\times S`$ such that $`(g_1+g_2,g_3+s)=(g_1,g_3)+(g_2,g_3)+(g_1,s)+(g_2,s)`$ (2.44) for $`g_1,g_2,g_3G,sS`$, equipped with a vertex operator map $`Y`$ from $`V`$ to $`(\mathrm{End}W)\{z\}`$ such that the axioms (2.19), (2.20) and (2.23) hold with suitable changes. Sometimes, to distinguish the vertex operator map $`Y`$ for a module $`W`$ from that for the adjoint module $`V`$ we use notation $`Y_W`$. Using the proof of Lemma 2.2 of \[DLM\] with a slight modification we get: ###### Proposition 2.7 Let $`W`$ be a $`V`$-module. Then on $`W`$, $`Y(D(v),z)={\displaystyle \frac{d}{dz}}Y(v,z)\text{ for }vV.\mathrm{}`$ (2.45) At the end of this section we present the following simple generalization of Lemma 2.3.5 of \[Li2\]: ###### Lemma 2.8 Let $`(V,Y,\mathrm{𝟏},G,c(,),(,))`$ be a generalized vertex algebra and let $`W`$ be a $`V`$-module. Let $`uV^g,vV^h,n𝐙,𝐮_{(𝐢)}𝐕`$ for $`i=1,\mathrm{},k`$. If $`(z_1z_2)^{n+(g,h)}Y_V(u,z_1)Y_V(v,z_2)c(g,h)(1)^n(z_2z_1)^{n+(g,h)}Y_V(v,z_2)Y_V(u,z_1)`$ (2.46) $`=`$ $`{\displaystyle \underset{i=0}{\overset{k}{}}}Y_V(u_{(i)},z_2)\left({\displaystyle \frac{}{z_2}}\right)^i\left(z_1^1\delta (z_2/z_1)(z_2/z_1)^g\right)`$ then $`(z_1z_2)^{n+(g,h)}Y_W(u,z_1)Y_W(v,z_2)c(g,h)(1)^n(z_2z_1)^{n+(g,h)}Y_W(v,z_2)Y_W(u,z_1)`$ (2.47) $`=`$ $`{\displaystyle \underset{i=0}{\overset{k}{}}}Y_W(u_{(i)},z_2)\left({\displaystyle \frac{}{z_2}}\right)^i\left(z_1^1\delta (z_2/z_1)(z_2/z_1)^g\right).`$ Furthermore, if $`W`$ is a faithful module, the converse is also true. $`\mathrm{}`$ ###### Remark 2.9 With Lemma 2.8, we have the following loose statement: If a generalized vertex algebra $`V`$ is a module for a certain “algebra” with defining relations of type (2.46), then any $`V`$-module is also a module for this “algebra.” Conversely, if $`W`$ is a faithful $`V`$-module and it is a module for a certain “algebra,” then $`V`$ is also a module for the “algebra.” Examples for such “algebras” are affine Lie algebras, the Virasoro algebra, affine Griess algebra \[FLM\] and $`Z`$-algebras. ## 3 Generalized vertex algebras generated by parafermion operators Throughout this paper, $`G`$ is an abelian group, $`(,)`$ is a symmetric $`𝐙`$-bilinear $`𝐂/\mathrm{𝟐}𝐙`$-valued form on $`G`$, $`c(,)`$ is a $`𝐂^\times `$-valued alternating form on $`G`$, and $`S`$ is a $`G`$-set equipped with a $`𝐂/\mathrm{𝟐}𝐙`$-valued function denoted also by $`(,)`$ on $`G\times S`$ which satisfies (2.44). We fix a choice of representatives $`(g,h)`$ and $`(g,s)`$ in $`𝐂`$ for $`g,hG`$ and $`sS`$. However, it should be observed that the main notions and results do not depend on this choice. Let $`W=_{gG}W^g`$ be a $`G`$-graded vector space (over $`𝐂`$). By definition, $`(\mathrm{End}W)\{z\}=\left\{a(z)={\displaystyle \underset{n𝐂}{}}a_nz^{n1}\right\|a_n\mathrm{End}W\}.`$ (3.1) Following \[DL2\], for $`gG`$ we define an operator $`z^g`$ from $`W`$ to $`W\{z\}`$ by $`z^gw=z^{(g,s)}w\text{ for }wW^s,sS.`$ (3.2) Note: This operator of course depends on the choice of representatives of $`(g,s)`$. A formal series $`a(z)(\mathrm{End}W)\{z\}`$ is said to satisfy lower truncation condition if for every $`wW`$, there exist finitely many complex numbers $`\alpha _1,\mathrm{},\alpha _r`$ such that $$a(z)wz^{\alpha _1}W[[z]]+\mathrm{}+z^{\alpha _r}W[[z]].$$ Clearly, all such series form a subspace of $`(\mathrm{End}W)\{z\}`$. For $`gG`$, we define $`F(W)^g`$ to be the vector subspace of $`(\mathrm{End}W)\{z\}`$ consisting of $`a(z)`$ satisfying the lower truncation condition and $`z^{(g,s)}a(z)W^sW^{g+s}[[z,z^1]]\text{ for }sS.`$ (3.3) Note that the notion of $`F(W)^g`$ does not depend on the choice of representatives $`(g,s)`$. Set $`F(W)={\displaystyle \underset{gG}{}}F(W)^g.`$ (3.4) (It is a direct sum because $`W`$ is $`G`$-graded.) As indicated in the notions of $`F(W)^g`$ and $`F(W)`$, $`z`$ is treated as a dummy variable, i.e., $`a(z),a(z_1)`$ and $`a(z_2)`$ are considered as the same object. ###### Definition 3.1 Formal series $`a(z)F(W)^g,b(z)F(W)^h`$ are said to mutually satisfy the generalized weak commutativity if there exists a nonnegative integer $`k`$ such that $`(z_1z_2)^{k+(g,h)}a(z_1)b(z_2)=(1)^kc(g,h)(z_2z_1)^{k+(g,h)}b(z_2)a(z_1).`$ (3.5) Clearly, this definition is free of the choice of a representative of $`(g,h)`$. Note that (3.5) also holds if we replace $`k`$ by any integer greater than $`k`$. Recall that $`D`$ is the linear endomorphism of $`(\mathrm{End}W)\{z\}`$ such that $`D(a(z))=a^{}(z)`$ for $`a(z)(\mathrm{End}W)\{z\}`$, where $`a^{}(z)`$ is the formal derivative of $`a(z)`$. Clearly, $`D`$ preserves $`F(W)`$ and its $`G`$-grading. ###### Remark 3.2 It is easy to see that if $`a(z)`$ and $`b(z)`$ mutually satisfy the generalized weak commutativity, so do $`a^{}(z)`$ and $`b(z)`$. A homogeneous parafermion field operator on $`W`$ is a series $`a(z)`$ of $`F(W)^g`$ for some $`gG`$ that satisfies the generalized weak commutativity with itself. ###### Definition 3.3 Let $`a(z)F(W)^g,b(z)F(W)^h,g,hG`$. Suppose that $`a(z)`$ and $`b(z)`$ satisfy the generalized weak commutativity. Then for each $`n𝐂`$, we define an element $`a(z)_nb(z)`$ of $`(\mathrm{End}W)\{z\}`$ by $`(a(z)_nb(z))w=\mathrm{Res}_{z_0}\mathrm{Res}_{z_1}z_0^n\left({\displaystyle \frac{z+z_0}{z_1}}\right)^{(g,s)}X`$ (3.6) for $`wW^s,sS`$, where $`X`$ $`=`$ $`z_0^1\left({\displaystyle \frac{z_1z}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)a(z_1)b(z)w`$ (3.7) $`c(g,h)z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)b(z)a(z_1)w.`$ ###### Remark 3.4 Note that because of the generalized weak commutativity, $`z_0^{(g,h)}X`$, which involves only integral powers of $`z_0`$, contains only finitely many negative integral powers of $`z_0`$. Then $`a(z)_nb(z)`$ exists as a formal series for any choice of representatives $`(g,h),(g,s)`$. Furthermore, since $`\left({\displaystyle \frac{z_1z}{z_0}}\right)^m\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)=\delta \left({\displaystyle \frac{z_1z}{z_0}}\right),`$ (3.8) $`\left({\displaystyle \frac{zz_1}{z_0}}\right)^{2m}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)=\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)`$ (3.9) for $`m𝐙`$, the expression $`X`$ does not depend on the choice of a representative of $`(g,h)`$. For $`m0`$, we also have $`\left({\displaystyle \frac{z+z_0}{z_1}}\right)^mz_0^1\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)=z_0^1\delta \left({\displaystyle \frac{z_1z}{z_0}}\right),`$ $`\left({\displaystyle \frac{z+z_0}{z_1}}\right)^mz_0^1\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)=z_0^1\delta \left({\displaystyle \frac{zz_1}{z_0}}\right).`$ Then $`(a(z)_nb(z))w`$ does not depend on the choice of representatives of $`(g,h)`$ and $`(g,s)`$. If $`n(g,h)+𝐙`$, the right-hand side of (3.6) does not involve integral powers of $`z_0`$. Therefore, $`a(z)_nb(z)=0\text{ for }n(g,h)+𝐙.`$ (3.10) For $`n(g,h)+𝐙`$, we have $`(a(z)_nb(z))w`$ (3.11) $`=`$ $`{\displaystyle \underset{i0}{}}\mathrm{Res}_{z_1}\left({\displaystyle \genfrac{}{}{0pt}{}{(g,s)}{i}}\right)z_1^{(g,s)}z^{(g,s)i}`$ $`\left((z_1z)^{n+i}a(z_1)b(z)w(1)^{n(g,h)+i}c(g,h)(zz_1)^{n+i}b(z)a(z_1)w\right).`$ Again, because of the generalized weak commutativity, the sum is really a finite sum. Noticing that $`a(z_1)wW^{g+s}\{z_1\}`$ from (3.11) we have $$z^{(g,s)+(h,s)}(a(z)_nb(z))wW^{g+h+s}((z)).$$ Since $`(g+h,s)(g,s)(h,s)2𝐙`$, we have $`z^{(g+h,s)}(a(z)_nb(z))wW^{g+h+s}((z)).`$ (3.12) This shows that $`a(z)_nb(z)F(W)^{g+h}`$. To summarize we have: ###### Lemma 3.5 Suppose that $`a(z)F(W)^g,b(z)F(W)^h`$ satisfy the generalized weak commutativity. Then $`a(z)_nb(z)F(W)^{g+h}\text{ for }n𝐂.`$ (3.13) Furthermore, $`a(z)_nb(z)=0`$ if $`n(g,h)+𝐙`$, and $`a(z)_nb(z)=0`$ for $`n(g,h)+𝐙`$ with $`n(g,h)`$ being sufficiently large.$`\mathrm{}`$ ###### Remark 3.6 Note that in the ordinary untwisted case with $`G=0`$, or $`𝐙/\mathrm{𝟐}𝐙`$ (cf. \[Li2\]), $`a(z)_nb(z)`$ were well defined for all $`a(z),b(z)(\mathrm{End}W)[[z,z^1]]`$. In the generalized case, the definition of $`a(z)_nb(z)`$ requires the generalized weak commutativity. This is similar to the situation for twisted vertex operators in \[Li3\]. Write $`a(z)_nb(z)`$ in terms of generating series as $`Y(a(z),z_0)b(z)={\displaystyle \underset{n𝐂}{}}\left(a(z)_nb(z)\right)z_0^{n1},`$ (3.14) where $`z_0`$ is another formal variable. Then $`Y(a(z),z_0)b(z)=\mathrm{Res}_{z_1}\left({\displaystyle \frac{z+z_0}{z_1}}\right)^{(g,s)}X.`$ (3.15) Note that $`z_0^{(g,h)}Y(a(z),z_0)b(z)`$ involves only integral powers of $`z_0`$ and that the powers of $`z_0`$ are truncated from below. Thus, for $`\alpha 𝐂`$, $$(z+z_0)^\alpha Y(a(z),z_0)b(z)\text{exists }$$ in $`F(W)\{z_0,z\}`$, hence $$\left(\frac{z_1z_0}{z}\right)^\alpha z^1\delta \left(\frac{z_1z_0}{z}\right)Y(a(z),z_0)b(z)\text{exists }$$ in $`F(W)\{z_0,z_1,z\}`$. As expected we have: ###### Proposition 3.7 Let $`a(z)F(W)^g,b(z)F(W)^h,g,hG`$. Suppose that $`a(z),b(z)`$ mutually satisfy the generalized weak commutativity. Then for $`wW^s,sS`$, $`\left({\displaystyle \frac{z_1z_0}{z}}\right)^{(g,s)}z^1\delta \left({\displaystyle \frac{z_1z_0}{z}}\right)(Y(a(z),z_0)b(z))w`$ (3.16) $`=`$ $`z_0^1\left({\displaystyle \frac{z_1z}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)a(z_1)b(z)w`$ $`c(g,h)z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)b(z)a(z_1)w.`$ Proof. Let $`r`$ be a nonnegative integer such that $$z_1^{r+(g,s)}a(z_1)wW[[z_1]],$$ hence $`\mathrm{Res}_{z_1}z_1^{r+(g,s)}z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g,s)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)b(z)a(z_1)w=0.`$ (3.17) Then using (3.15) and the fundamental properties of the delta function, we have $`(z+z_0)^{r+(g,s)}(Y(a(z),z_0)b(z))w`$ (3.18) $`=`$ $`\mathrm{Res}_{z_1}(z+z_0)^rz_1^{(g,s)}X`$ $`=`$ $`\mathrm{Res}_{z_1}z_1^{r+(g,s)}z_0^1\left({\displaystyle \frac{z_1z}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)a(z_1)b(z)w`$ $`c(g,h)\mathrm{Res}_{z_1}z_1^{r+(g,s)}z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)b(z)a(z_1)w`$ $`=`$ $`\mathrm{Res}_{z_1}z_1^{r+(g,s)}z_1^1\left({\displaystyle \frac{z_0+z}{z_1}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_0+z}{z_1}}\right)a(z_1)b(z)w`$ $`=`$ $`\mathrm{Res}_{z_1}z_1^{r+(g,s)(g,h+s)+(g,h)}(z_0+z)^{(g,h)}z_1^1\delta \left({\displaystyle \frac{z_0+z}{z_1}}\right)\left(z_1^{(g,h+s)}a(z_1)b(z)w\right)`$ $`=`$ $`(z_0+z)^{r+(g,s)}a(z_0+z)b(z)w,`$ noting that $`(g,s)(g,h+s)+(g,h)2𝐙`$. Similar to the proof of Proposition 2.5, this generalized weak associativity relation together with the generalized weak commutativity relation implies the generalized Jacobi identity. $`\mathrm{}`$ ###### Proposition 3.8 Let $`a(z)F(W)^{g_1},b(z)F(W)^{g_2},c(z)F(W)^{g_3}`$ with $`g_1,g_2,g_3G`$. Suppose that $`a(z),b(z),c(z)`$ mutually satisfy the generalized weak commutativity. Then for $`n𝐂`$, $`a(z)_nb(z)`$ and $`c(z)`$ satisfy the generalized weak commutativity. Proof. Let $`r`$ be a positive integer such that the following identities hold: $`(z_1z_2)^{r+(g_1,g_3)}a(z_1)c(z_2)=(1)^rc(g_1,g_3)(z_2z_1)^{r+(g_1,g_3)}c(z_2)a(z_1),`$ $`(z_1z_2)^{r+(g_2,g_3)}b(z_1)c(z_2)=(1)^rc(g_2,g_3)(z_2z_1)^{r+(g_2,g_3)}c(z_2)b(z_1).`$ Let $`wW^s,sS`$. Using the Jacobi identity relation (3.16) and the above generalized weak commutativity relations we get $`(z_1z_2)^{r+(g_1,g_3)}(zz_2)^{r+(g_2,g_3)}`$ (3.19) $`\left({\displaystyle \frac{z_1z_0}{z}}\right)^{(g_1,g_3+s)}z^1\delta \left({\displaystyle \frac{z_1z_0}{z}}\right)(Y(a(z),z_0)b(z))c(z_2)w`$ $`=`$ $`z_0^1\left({\displaystyle \frac{z_1z}{z_0}}\right)^{(g_1,g_2)}\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)`$ $`(z_1z_2)^{r+(g_1,g_3)}(zz_2)^{r+(g_2,g_3)}a(z_1)b(z)c(z_2)w`$ $`c(g_1,g_2)z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g_1,g_2)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)`$ $`(z_1z_2)^{r+(g_1,g_3)}(zz_2)^{r+(g_2,g_3)}b(z)a(z_1)c(z_2)w`$ $`=`$ $`c(g_1,g_3)c(g_2,g_3)(z_2z_1)^{r+(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}`$ $`z_0^1\left({\displaystyle \frac{z_1z}{z_0}}\right)^{(g_1,g_2)}\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)c(z_2)a(z_1)b(z)w`$ $`c(g_1,g_3)c(g_2,g_3)(z_2z_1)^{r+(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}`$ $`c(g_1,g_2)z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g_1,g_2)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)c(z_2)b(z)a(z_1)w`$ $`=`$ $`c(g_1,g_3)c(g_2,g_3)(z_2z_1)^{r+(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}`$ $`\left({\displaystyle \frac{z_1z_0}{z}}\right)^{(g_1,s)}z^1\delta \left({\displaystyle \frac{z_1z_0}{z}}\right)c(z_2)(Y(a(z),z_0)b(z))w`$ $`=`$ $`c(g_1+g_2,g_3)(z_2z_1)^{r+(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}`$ $`\left({\displaystyle \frac{z_1z_0}{z}}\right)^{(g_1,s)}z^1\delta \left({\displaystyle \frac{z_1z_0}{z}}\right)c(z_2)(Y(a(z),z_0)b(z))w.`$ Thus $`(z_1z_2)^{r+(g_1,g_3)}(zz_2)^{r+(g_2,g_3)}`$ (3.20) $`\left({\displaystyle \frac{z_1z_0}{z}}\right)^{(g_1,g_3)}z^1\delta \left({\displaystyle \frac{z_1z_0}{z}}\right)(Y(a(z),z_0)b(z))c(z_2)w`$ $`=`$ $`c(g_1+g_2,g_3)(z_2z_1)^{r+(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}`$ $`z^1\delta \left({\displaystyle \frac{z_1z_0}{z}}\right)c(z_2)(Y(a(z),z_0)b(z))w,`$ noting that $`(g_1,g_3+s)(g_1,g_3)(g_1,s)2𝐙`$. Using the fundamental properties of delta-function we get $`(z+z_0z_2)^{r+(g_1,g_3)}(zz_2)^{r+(g_2,g_3)}`$ (3.21) $`z_1^1\delta \left({\displaystyle \frac{z+z_0}{z_1}}\right)(Y(a(z),z_0)b(z))c(z_2)w`$ $`=`$ $`c(g_1+g_2,g_3)(z_2zz_0)^{r+(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}`$ $`z_1^1\delta \left({\displaystyle \frac{z+z_0}{z_1}}\right)c(z_2)(Y(a(z),z_0)b(z))w.`$ Taking $`\mathrm{Res}_{z_1}`$ from (3.21) we obtain $`(z+z_0z_2)^{r+(g_1,g_3)}(zz_2)^{r+(g_2,g_3)}(Y(a(z),z_0)b(z))c(z_2)w`$ (3.22) $`=`$ $`c(g_1+g_2,g_3)(z_2zz_0)^{r+(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}c(z_2)(Y(a(z),z_0)b(z))w.`$ Let $`n𝐂`$ be arbitrarily fixed. Since $`a(z)_nb(z)=0`$ for $`n(g_1,g_2)+𝐙`$, we only need to consider $`n(g_1,g_2)+𝐙`$. Let $`N`$ be a fixed nonnegative integer such that $`a(z)_mb(z)=0`$ for $`mN+n`$, so that $`\mathrm{Res}_{z_0}z_0^{N+n+i}Y(a(z),z_0)b(z)=0`$ (3.23) for $`iN`$. We may replace $`r`$ by $`r+N`$ so that we may assume that $`rN`$ and that $`r+n(g_1,g_2)>0`$. Then using (3.22) we obtain $`(zz_2)^{3r+(g_1,g_3)+(g_2,g_3)}(a(z)_nb(z))c(z_2)w`$ (3.24) $`=`$ $`\mathrm{Res}_{z_0}z_0^n(zz_2)^{3r+(g_1,g_3)+(g_2,g_3)}(Y(a(z),z_0)b(z))c(z_2)w`$ $`=`$ $`\mathrm{Res}_{z_0}{\displaystyle \underset{i0}{}}(1)^i\left({\displaystyle \genfrac{}{}{0pt}{}{2r+(g_1,g_3)}{i}}\right)z_0^{n+i}(zz_2+z_0)^{2r+(g_1,g_3)i}`$ $`(zz_2)^{r+(g_2,g_3)}(Y(a(z),z_0)b(z))c(z_2)w`$ $`=`$ $`\mathrm{Res}_{z_0}{\displaystyle \underset{i=0}{\overset{N}{}}}(1)^i\left({\displaystyle \genfrac{}{}{0pt}{}{2r+(g_1,g_3)}{i}}\right)z_0^{n+i}`$ $`(zz_2+z_0)^{2r+(g_1,g_3)i}(zz_2)^{r+(g_2,g_3)}(Y(a(z),z_0)b(z))c(z_2)w`$ $`=`$ $`\mathrm{Res}_{z_0}{\displaystyle \underset{i=0}{\overset{N}{}}}(1)^i\left({\displaystyle \genfrac{}{}{0pt}{}{2r+(g_1,g_3)}{i}}\right)z_0^{n+i}c(g_1+g_2,g_3)(1)^{ri}`$ $`(z_2zz_0)^{2r+(g_1,g_3)i}(z+z_0)^{(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}c(z_2)(Y(a(z),z_0)b(z))w`$ $`=`$ $`\mathrm{Res}_{z_0}{\displaystyle \underset{i0}{}}(1)^r\left({\displaystyle \genfrac{}{}{0pt}{}{2r+(g_1,g_3)}{i}}\right)z_0^{n+i}c(g_1+g_2,g_3)`$ $`(z_2zz_0)^{2r+(g_1,g_3)i}(z+z_0)^{(g_1,g_3)}(z_2z)^{r+(g_2,g_3)}c(z_2)(Y(a(z),z_0)b(z))w`$ $`=`$ $`(1)^r\mathrm{Res}_{z_0}z_0^nc(g_1+g_2,g_3)(z_2z)^{2r+(g_1,g_3)+(g_2,g_3)}c(z_2)(Y(a(z),z_0)b(z))w`$ $`=`$ $`(1)^rc(g_1+g_2,g_3)(z_2z)^{3r+(g_1,g_3)+(g_2,g_3)}c(z_2)(a(z)_nb(z))w.`$ Since $`(g_1+g_2,g_3)(g_1,g_3)(g_2,g_3)2𝐙`$, there exist $`k,k^{}N`$ such that $$2k+(g_1,g_3)+(g_2,g_3)=(g_1+g_2,g_3)+2k^{}.$$ Then $`(zz_2)^{3r+2k^{}+(g_1+g_2,g_3)}(a(z)_nb(z))c(z_2)w`$ (3.25) $`=`$ $`(1)^rc(g_1+g_2,g_3)(z_2z)^{3r+2k^{}+(g_1,+g_2,g_3)}c(z_2)(a(z)_nb(z))w.`$ This proves that $`a(z)_nb(z)`$ and $`c(z)`$ mutually satisfy the generalized weak commutativity. $`\mathrm{}`$ ###### Definition 3.9 A $`G`$-graded subspace $`A`$ of $`F(W)`$ is called a generalized vertex pre-algebra if every pair of homogeneous elements of $`A`$ satisfy the generalized weak commutativity. For homogeneous $`a(z),b(z)A`$, $`a(z)_nb(z)`$ was defined for $`n𝐂`$. Using linearity, we define $`a(z)_nb(z)`$ for all $`a(z),b(z)A`$, so that $`Y(a(z),z_0)b(z)`$ is defined for all $`a(z),b(z)A`$. Furthermore, $`A`$ is said to be closed if $`a(z)_nb(z)A`$ for $`a(z),b(z)A,n𝐂`$. Since the identity operator $`I(z)=\text{id}_W`$ (independent of $`z`$) and any element of $`F(W)`$ mutually satisfy the generalized weak commutativity, any maximal generalized vertex pre-algebra contains $`I(z)`$. In view of Proposition 3.8 we immediately have: ###### Corollary 3.10 Any maximal generalized vertex pre-algebra $`V`$ contains the identity operator $`I(z)`$ and it is closed and $`D`$-stable.$`\mathrm{}`$ For the rest of this section, $`V`$ will be a fixed closed generalized vertex pre-algebra containing the identity operator $`I(z)`$ on $`W`$. The same proof of Lemma 3.10 of \[Li3\] gives the following results: ###### Lemma 3.11 For $`a(z)V`$, we have $`Y(I(z),z_0)a(z)=a(z),`$ (3.26) $`Y(a(z),z_0)I(z)=e^{z_0\frac{}{z}}a(z)=a(z+z_0).\mathrm{}`$ (3.27) Furthermore, we have: ###### Lemma 3.12 Suppose that $`a(z)F(W)^g,b(z)F(W)^h`$ mutually satisfy the generalized weak commutativity. Then $`{\displaystyle \frac{}{z_0}}Y(a(z),z_0)b(z)=Y(a^{}(z),z_0)b(z)=Y(D(a(z)),z_0)b(z),`$ (3.28) $`[D,Y(a(z),z_0)]b(z)={\displaystyle \frac{}{z_0}}Y(a(z),z_0)b(z).`$ (3.29) Proof. Let $`wW^s,sS`$. Recall that the generalized Jacobi relation (3.16) holds. We also have $`\left({\displaystyle \frac{z_1z_0}{z}}\right)^{(g,s)}z^1\delta \left({\displaystyle \frac{z_1z_0}{z}}\right)(Y(a^{}(z),z_0)b(z))w`$ (3.30) $`=`$ $`z_0^1\left({\displaystyle \frac{z_1z}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)a^{}(z_1)b(z)w`$ $`c(g,h)z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)b(z)a^{}(z_1)w.`$ Let $`X_L`$ and $`X_R`$ be the term on the left-hand side and the term on the right-hand side of (3.16). Using (3.30) we get $`\left({\displaystyle \frac{z+z_0}{z_1}}\right)^{(g,s)}z_1^1\delta \left({\displaystyle \frac{z+z_0}{z_1}}\right)(Y(a^{}(z),z_0)b(z))w`$ (3.31) $`=`$ $`\left({\displaystyle \frac{z_1z_0}{z}}\right)^{(g,s)}z^1\delta \left({\displaystyle \frac{z_1z_0}{z}}\right)(Y(a^{}(z),z_0)b(z))w`$ $`=`$ $`{\displaystyle \frac{}{z_1}}X_Ra(z_1)b(z)w{\displaystyle \frac{}{z_1}}\left(z_0^1\left({\displaystyle \frac{z_1z}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)\right)`$ $`+c(g,h)b(z)a(z_1)w{\displaystyle \frac{}{z_1}}\left(z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)\right)`$ $`=`$ $`{\displaystyle \frac{}{z_1}}X_L+a(z_1)b(z)w{\displaystyle \frac{}{z_0}}\left(z_0^1\left({\displaystyle \frac{z_1z}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)\right)`$ $`c(g,h)b(z)a(z_1)w{\displaystyle \frac{}{z_0}}\left(z_0^1\left({\displaystyle \frac{zz_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)\right)`$ $`=`$ $`{\displaystyle \frac{}{z_1}}X_L+{\displaystyle \frac{}{z_0}}X_R`$ $`=`$ $`{\displaystyle \frac{}{z_1}}X_L+{\displaystyle \frac{}{z_0}}X_L.`$ We are using the fact $`{\displaystyle \frac{}{z_1}}\left(\left({\displaystyle \frac{zz_1}{z_0}}\right)^\alpha z_0^1\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)\right)={\displaystyle \frac{}{z_0}}\left(\left({\displaystyle \frac{zz_1}{z_0}}\right)^\alpha z_0^1\delta \left({\displaystyle \frac{zz_1}{z_0}}\right)\right).`$ (3.32) for $`\alpha 𝐂`$. Multiplying by $`z_1^{(g,s)}`$, then taking $`\mathrm{Res}_{z_1}`$ and using a variant of (3.32) we get $`(z+z_0)^{(g,s)}(Y(a^{}(z),z_0)b(z))w`$ (3.33) $`=`$ $`\mathrm{Res}_{z_1}z_1^{(g,s)}{\displaystyle \frac{}{z_1}}X_L+\mathrm{Res}_{z_1}z_1^{(g,s)}{\displaystyle \frac{}{z_0}}X_L`$ $`=`$ $`\mathrm{Res}_{z_1}(Y(a(z),z_0)b(z))wz_1^{(g,s)}{\displaystyle \frac{}{z_1}}\left(\left({\displaystyle \frac{z+z_0}{z_1}}\right)^{(g,s)}z_1^1\delta \left({\displaystyle \frac{z+z_0}{z_1}}\right)\right)`$ $`+\mathrm{Res}_{z_1}(Y(a(z),z_0)b(z))wz_1^{(g,s)}{\displaystyle \frac{}{z_0}}\left(\left({\displaystyle \frac{z+z_0}{z_1}}\right)^{(g,s)}z_1^1\delta \left({\displaystyle \frac{z+z_0}{z_1}}\right)\right)`$ $`+\mathrm{Res}_{z_1}z_1^{(g,s)}\left({\displaystyle \frac{z+z_0}{z_1}}\right)^{(g,s)}z_1^1\delta \left({\displaystyle \frac{z+z_0}{z_1}}\right){\displaystyle \frac{}{z_0}}(Y(a(z),z_0)b(z))w`$ $`=`$ $`\mathrm{Res}_{z_1}z_1^{(g,s)}\left({\displaystyle \frac{z+z_0}{z_1}}\right)^{(g,s)}z_1^1\delta \left({\displaystyle \frac{z+z_0}{z_1}}\right){\displaystyle \frac{}{z_0}}(Y(a(z),z_0)b(z))w`$ $`=`$ $`(z+z_0)^{(g,s)}{\displaystyle \frac{}{z_0}}(Y(a(z),z_0)b(z))w.`$ Multiplying by $`(z+z_0)^{(g,s)}`$ from left we get the first identity. The second identity follows from a similar argument and the fact: $`{\displaystyle \frac{}{z_1}}\left(\left({\displaystyle \frac{z_1z}{z_0}}\right)^\alpha z_0^1\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)\right)={\displaystyle \frac{}{z}}\left(\left({\displaystyle \frac{z_1z}{z_0}}\right)^\alpha z_0^1\delta \left({\displaystyle \frac{z_1z}{z_0}}\right)\right)`$ (3.34) for $`\alpha 𝐂`$.$`\mathrm{}`$ ###### Proposition 3.13 Let $`V`$ be a closed generalized vertex pre-algebra of parafermion operators on $`W`$. Let $`a(z)V^{g_1},b(z)V^{g_2}`$ and let $`rN`$ be such that $`(z_1z_2)^{r+(g_1,g_2)}a(z_1)b(z_2)=(1)^rc(g_1,g_2)(z_2z_1)^{r+(g_1,g_2)}b(z_2)a(z_1).`$ (3.35) Then $`(x_1x_2)^{r+(g_1,g_2)}Y(a(z),x_1)Y(b(z),x_2)`$ (3.36) $`=`$ $`(1)^rc(g_1,g_2)(x_2x_1)^{r+(g_1,g_2)}Y(b(z),x_2)Y(a(z),x_1),`$ acting on $`V`$. Proof. Let $`c(z)V^{g_3},wW^s,g_3G,sS`$. Using the generalized Jacobi relation (3.16) and the fundamental properties of the delta-function we get $`\left({\displaystyle \frac{z+x_1}{z_1}}\right)^{(g_1,s)}z_1^1\delta \left({\displaystyle \frac{z+x_1}{z_1}}\right)\left({\displaystyle \frac{z+x_2}{z_2}}\right)^{(g_2,s)}z_2^1\delta \left({\displaystyle \frac{z+x_2}{z_2}}\right)`$ (3.37) $`(Y(a(z),x_1)Y(b(z),x_2)c(z))w`$ $`=`$ $`\left({\displaystyle \frac{z_1x_1}{z}}\right)^{(g_1,s)}z^1\delta \left({\displaystyle \frac{z_1x_1}{z}}\right)\left({\displaystyle \frac{z_2x_2}{z}}\right)^{(g_2,s)}z^1\delta \left({\displaystyle \frac{z_2x_2}{z}}\right)`$ $`(Y(a(z),x_1)Y(b(z),x_2)c(z))w`$ $`=`$ $`\left({\displaystyle \frac{z_2x_2}{z}}\right)^{(g_2,s)}z^1\delta \left({\displaystyle \frac{z_2x_2}{z}}\right)`$ $`x_1^1\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)a(z_1)(Y(b(z),x_2)c(z))w`$ $`c(g_1,g_2+g_3)\left({\displaystyle \frac{z_2x_2}{z}}\right)^{(g_2,s)}z^1\delta \left({\displaystyle \frac{z_2x_2}{z}}\right)`$ $`x_1^1\left({\displaystyle \frac{zz_1}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{zz_1}{x_1}}\right)(Y(b(z),x_2)c(z))a(z_1)w`$ $`=`$ $`x_1^1\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)`$ $`x_2^1\left({\displaystyle \frac{z_2z}{x_2}}\right)^{(g_2,g_3)}\delta \left({\displaystyle \frac{z_2z}{x_2}}\right)a(z_1)b(z_2)c(z)w`$ $`x_1^1\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)`$ $`c(g_2,g_3)x_2^1\left({\displaystyle \frac{zz_2}{x_2}}\right)^{(g_2,g_3)}\delta \left({\displaystyle \frac{zz_2}{x_2}}\right)a(z_1)c(z)b(z_2)w`$ $`c(g_1,g_2+g_3)x_1^1\left({\displaystyle \frac{zz_1}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{zz_1}{x_1}}\right)`$ $`\left({\displaystyle \frac{z_2x_2}{z}}\right)^{(g_2,s)}z^1\delta \left({\displaystyle \frac{z_2x_2}{z}}\right)(Y(b(z),x_2)c(z))a(z_1)w.`$ Let $`p,qN`$ be such that $`z_1^{p+(g_1,s)}a(z_1)wW[[z_1]],z_2^{q+(g_2,s)}b(z_2)wW[[z_2]].`$ (3.38) Notice that for $`jN`$ we have $`(z+x_3)^jA=z_3^jA,`$ (3.39) where $`A`$ is one of the three delta-functions $`x_3^1\delta \left(\frac{z_3z}{x_3}\right),x_3^1\delta \left(\frac{zz_3}{x_3}\right)`$ and $`z^1\delta \left(\frac{z_3x_3}{z}\right)`$. Applying $`\mathrm{Res}_{z_1}\mathrm{Res}_{z_2}(z+x_1)^s(z+x_2)^qz_1^{(g_1,s)}z_2^{(g_2,s)}`$ to (3.37), then using (3.39) and (3.38) we get $`(z+x_1)^{p+(g_1,h)}(z+x_2)^{q+(g_2,s)}(Y(a(z),x_1)Y(b(z),x_2)c(z))w`$ (3.40) $`=`$ $`\mathrm{Res}_{z_1}\mathrm{Res}_{z_2}z_1^{p+(g_1,h)}z_2^{q+(g_2,s)}x_1^1\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)`$ $`x_2^1\left({\displaystyle \frac{z_2z}{x_2}}\right)^{(g_2,g_3)}\delta \left({\displaystyle \frac{z_2z}{x_2}}\right)a(z_1)b(z_2)c(z)w`$ $`\mathrm{Res}_{z_1}\mathrm{Res}_{z_2}z_1^{p+(g_1,s)}x_1^1\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)`$ $`c(g_2,g_3)x_2^1\left({\displaystyle \frac{zz_2}{x_2}}\right)^{(g_2,g_3)}\delta \left({\displaystyle \frac{zz_2}{x_2}}\right)a(z_1)c(z)\left(z_2^{q+(g_2,s)}b(z_2)w\right)`$ $`\mathrm{Res}_{z_1}\mathrm{Res}_{z_2}z_2^{q+(g_2,s)}c(g_1,g_2+g_3)x_1^1\left({\displaystyle \frac{zz_1}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{zz_1}{x_1}}\right)`$ $`\left({\displaystyle \frac{z_2x_2}{z}}\right)^{(g_2,s)}z^1\delta \left({\displaystyle \frac{z_2x_2}{z}}\right)(Y(b(z),x_2)c(z))\left(z_1^{p+(g_1,s)}a(z_1)w\right)`$ $`=`$ $`\mathrm{Res}_{z_1}\mathrm{Res}_{z_2}z_1^{p+(g_1,s)}z_2^{q+(g_2,s)}x_1^1\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_2+g_3)}\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)`$ $`x_2^1\left({\displaystyle \frac{z_2z}{x_2}}\right)^{(g_2,g_3)}\delta \left({\displaystyle \frac{z_2z}{x_2}}\right)a(z_1)b(z_2)c(z)w.`$ Notice that $`(x_1x_2)^{r+(g_1,g_2)}x_1^1\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)x_2^1\delta \left({\displaystyle \frac{z_2z}{x_2}}\right)`$ (3.41) $`=`$ $`(x_1z_2+z)^{r+(g_1,g_2)}x_1^1\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)x_2^1\delta \left({\displaystyle \frac{z_2z}{x_2}}\right)`$ $`=`$ $`\left({\displaystyle \frac{z_1z}{x_1}}\right)^{r(g_1,g_2)}(z_1z_2)^{r+(g_1,g_2)}x_1^1\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)x_2^1\delta \left({\displaystyle \frac{z_2z}{x_2}}\right)`$ $`=`$ $`\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_2)}(z_1z_2)^{r+(g_1,g_2)}x_1^1\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)x_2^1\delta \left({\displaystyle \frac{z_2z}{x_2}}\right).`$ Then $`(z+x_1)^{p+(g_1,s)}(z+x_2)^{q+(g_2,s)}(x_1x_2)^{r+(g_1,g_2)}(Y(a(z),x_1)Y(b(z),x_2)c(z))w`$ (3.42) $`=`$ $`\mathrm{Res}_{z_1}\mathrm{Res}_{z_2}z_1^{p+(g_1,s)}z_2^{q+(g_2,s)}x_1^1\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_3)}\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)`$ $`x_2^1\left({\displaystyle \frac{z_2z}{x_2}}\right)^{(g_2,g_3)}\delta \left({\displaystyle \frac{z_2z}{x_2}}\right)(z_1z_2)^{r+(g_1,g_2)}a(z_1)b(z_2)c(z)w.`$ Using the obvious symmetry, we have $`(z+x_1)^{p+(g_1,s)}(z+x_2)^{q+(g_2,s)}(x_2x_1)^{r+(g_1,g_2)}(Y(b(z),x_2)Y(a(z),x_1)c(z))w`$ (3.43) $`=`$ $`\mathrm{Res}_{z_1}\mathrm{Res}_{z_2}z_1^{p+(g_1,s)}z_2^{q+(g_2,s)}x_2^1\left({\displaystyle \frac{z_2z}{x_2}}\right)^{(g_2,g_3)}\delta \left({\displaystyle \frac{z_2z}{x_2}}\right)`$ $`x_1^1\left({\displaystyle \frac{z_1z}{x_1}}\right)^{(g_1,g_3)}\delta \left({\displaystyle \frac{z_1z}{x_1}}\right)(z_2z_1)^{r+(g_1,g_2)}b(z_2)a(z_1)c(z)w.`$ Therefore $`(z+x_1)^{p+(g_1,s)}(z+x_2)^{q+(g_2,s)}(x_1x_2)^{r+(g_1,g_2)}(Y(a(z),x_1)Y(b(z),x_2)c(z))w`$ (3.44) $`=`$ $`(1)^rc(g_1,g_2)(z+x_1)^{p+(g_1,s)}(z+x_2)^{q+(g_2,s)}(x_2x_1)^{r+(g_1,g_2)}`$ $`(Y(b(z),x_2)Y(a(z),x_1)c(z))w.`$ Multiplying both sides by $`(z+x_1)^{p(g_1,s)}(z+x_2)^{q(g_2,s)}`$ we obtain $`(x_1x_2)^{r+(g_1,g_2)}(Y(a(z),x_1)Y(b(z),x_2)c(z))w`$ (3.45) $`=`$ $`(1)^rc(g_1,g_2)(x_2x_1)^{r+(g_1,g_2)}(Y(b(z),x_2)Y(a(z),x_1)c(z))w.`$ Then the generalized weak commutativity (3.36) follows immediately.$`\mathrm{}`$ Now, we are in a position to present our main theorem: ###### Theorem 3.14 Let $`V`$ be a closed generalized vertex pre-algebra of parafermionic field operators on $`W`$, containing the identity operator $`I(z)`$. Then $`V`$ is a generalized vertex algebra and $`W`$ is a canonical $`V`$-module with $`Y(a(z),z_1)=a(z_1)`$ for $`a(z)V`$. Proof. It follows from Propositions 2.6, 3.13 and Lemmas 3.5, 3.11 and 3.12 that $`(V,Y)`$ is a generalized vertex algebra. It follows from Proposition 3.7 that $`W`$ is a $`V`$-module under the natural action. $`\mathrm{}`$ The following is a very useful consequence: ###### Theorem 3.15 Let $`\mathrm{\Gamma }`$ be a set of homogeneous parafermion field operators on $`W`$ that satisfy the generalized weak commutativity. Then the subspace $`\mathrm{\Gamma }`$ of $`F(W)`$, linearly spanned by $`a^{(1)}(z)_{n_1}\mathrm{}a^{(r)}(z)_{n_r}I(z)`$ (3.46) for $`rN,a^{(i)}(z)\mathrm{\Gamma },n_1,\mathrm{},n_r𝐂`$, equipped with the vertex operator map $`Y`$ is a generalized vertex algebra with $`W`$ as a natural module. Proof. By Zorn’s lemma, there exists a maximal generalized vertex pre-algebra $`V`$ containing $`\mathrm{\Gamma }`$. By Corollary 3.10, $`V`$ is closed and contains $`I(z)`$. By Theorem 3.14, $`V`$ is a generalized vertex algebra. It follows from Proposition 14.8 of \[DL2\] that $`\mathrm{\Gamma }`$ is a generalized vertex subalgebra of $`V`$. Clearly, $`\mathrm{\Gamma }`$ does not depend on the choice of $`V`$ and it is the intersection of all closed generalized vertex pre-algebra containing $`\mathrm{\Gamma }`$ and the identity operator $`I(z)`$. $`\mathrm{}`$ ###### Lemma 3.16 Let $`V`$ be a generalized vertex algebra and let $`U`$ be a graded subspace which generates $`V`$. Let $`W`$ be a $`V`$-module and let $`eW^0`$ such that $`Y(u,z)eV[[z]]`$ for $`uU`$. Then $`Y(v,z)eW[[z]]\text{ for all }vV.`$ (3.47) Proof. Let $`V^{}`$ be the collection of all $`vV`$ such that $`Y(v,z)eV[[z]]`$. Clearly, $`V^{}`$ is a graded subspace of $`V`$ containing $`\mathrm{𝟏}`$ and $`U`$. Let $`u(V^{})^g,v(V^{})^h,g,hG`$. Then $`Y(u,z)e,Y(v,z)eV[[z]]`$. Applying the generalized Jacobi identity to $`e`$, then taking $`\mathrm{Res}_{z_1}`$ we get $`Y(Y(u,z_0)v,z_2)e`$ (3.48) $`=`$ $`\mathrm{Res}_{z_1}z_0^1\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y(u,z_1)Y(v,z_2)e`$ $`\mathrm{Res}_{z_1}c(g,h)z_0^1\left({\displaystyle \frac{z_2z_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)Y(v,z_2)Y(u,z_1)e`$ $`=`$ $`\mathrm{Res}_{z_1}z_0^1\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y(u,z_1)Y(v,z_2)e`$ $``$ $`(V((z_0)))[[z_2]].`$ Then $`Y(u_nv,z_2)eV[[z_2]]\text{ for }n𝐂.`$ (3.49) Thus $`u_nvV^{}`$. Thus, $`V^{}`$ is a subalgebra of $`V`$, containing $`\mathrm{𝟏}`$ and $`U`$. Consequently, $`V^{}=V`$. The proof is complete.$`\mathrm{}`$ The same proof of Proposition 3.4 of \[Li1\] gives: ###### Proposition 3.17 Let $`V`$ be a generalized vertex algebra, let $`W`$ be a $`V`$-module and let $`eW^0`$ such that $`Y(v,z)eV[[z]]`$ for $`vV`$. Then $`Y(v,z)e=e^{zL(1)}v_1e\text{ for }vV,`$ (3.50) and the linear map $`f:`$ $`VW`$ (3.51) $`vv_1e`$ is a $`V`$-homomorphism. Furthermore, if $`W`$ is faithful and generated by $`e`$, then $`f`$ is an isomorphism.$`\mathrm{}`$ Now we have the following result the first part of which is a generalization of a result obtained in \[FKRW\] and \[PM\]: ###### Theorem 3.18 Let $`V`$ be a $`G`$-graded vector space and $`\mathrm{𝟏}V^0`$ and let $`U`$ a generalized vertex pre-algebra of parafermions on $`V`$ such that $`\psi (z)\mathrm{𝟏}V[[z]]`$ and such that $`V`$ is generated from $`\mathrm{𝟏}`$ by all component operators of $`\psi (z)z^g`$ for $`\psi (z)U^g,gG`$. Then there exists a unique generalized vertex algebra structure on $`V`$ such that $`\mathrm{𝟏}`$ is the vacuum vector and that $`Y(a,z)=a(z)`$ for $`aU`$. Furthermore, this generalized vertex algebra $`V`$ is isomorphic to the generalized vertex algebra generated by $`U`$ inside $`(\mathrm{End}V)\{z\}`$. Proof. Let $`\overline{U}`$ be the generalized vertex algebra generated by $`U`$ inside $`F(V)`$. Then $`V`$ is a $`\overline{U}`$-module and $`Y(u,z)\mathrm{𝟏}V[[z]]`$ for $`uU`$. Since $`V`$ is a faithful $`\overline{U}`$-module, it follows from Lemma 3.16 and Proposition 3.17 that the linear map $`f`$ from $`\overline{U}`$ to $`V`$ such that $`f(u)=u_1\mathrm{𝟏}`$ for $`u\overline{U}`$ is an one-to-one $`\overline{U}`$-homomorphism. Clearly, $`f`$ is onto. Then $`V=\overline{U}`$ has a natural generalized vertex algebra structure. The other assertions are clear.$`\mathrm{}`$ Similar to Lemma 3.16 we have the following result: ###### Lemma 3.19 Let $`V`$ be a generalized vertex algebra with a graded generating subspace $`A`$, $`W`$ be a $`V`$-module and $`D_W`$ be a grading-preserving endomorphism of $`W`$ such that $`[D_W,Y(v,z)]={\displaystyle \frac{d}{dz}}Y(v,z)(=Y(D(v),z))\text{ for }vA.`$ (3.52) Then (3.52) for all $`vV`$. Proof. Recall that $`Y(D(v),z)=\frac{d}{dz}Y(v,z)`$ for all $`vV`$. Set $`K=\{vV\|[D_W,Y(v,z)]=Y(D(v),z)={\displaystyle \frac{d}{dz}}Y(v,z)\}.`$ (3.53) Then $`AK`$. Let $`uKV^g,vKV^h,g,hG`$ and let $`wW^s,sS`$. Then we have $`z_2^1\left({\displaystyle \frac{z_1z_0}{z_2}}\right)^{(g,s)}\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)[D_W,Y(Y(u,z_0)v,z_2)]w`$ (3.54) $`=`$ $`z_0^1\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)[D_W,Y(u,z_1)]Y(v,z_2)w`$ $`+z_0^1\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y(u,z_1)[D_W,Y(v,z_2)]w`$ $`c(g,h)z_0^1\left({\displaystyle \frac{z_2z_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)[D_W,Y(v,z_2)]Y(u,z_1)w`$ $`c(g,h)z_0^1\left({\displaystyle \frac{z_2z_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)Y(v,z_2)[D_W,Y(u,z_1)]w`$ $`=`$ $`z_0^1\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y(D(u),z_1)Y(v,z_2)w`$ $`+z_0^1\left({\displaystyle \frac{z_1z_2}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y(u,z_1)Y(D(v),z_2)w`$ $`c(g,h)z_0^1\left({\displaystyle \frac{z_2z_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)Y(D(v),z_2)Y(u,z_1)w`$ $`c(g,h)z_0^1\left({\displaystyle \frac{z_2z_1}{z_0}}\right)^{(g,h)}\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)Y(v,z_2)Y(D(u),z_1)w`$ $`=`$ $`z_2^1\left({\displaystyle \frac{z_1z_0}{z_2}}\right)^{(g,s)}\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)(Y(Y(D(u),z_0)v,z_2)+Y(Y(u,z_0)D(v),z_2))w`$ $`=`$ $`z_2^1\left({\displaystyle \frac{z_1z_0}{z_2}}\right)^{(g,s)}\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)Y(DY(u,z_0)v,z_2)w.`$ This gives $`[D_W,Y(Y(u,z_0)v,z_2)]w=Y(DY(u,z_0)v,z_2)w.`$ (3.55) Then $`Y(u,z_0)vK\{z_0\}`$. Thus $`K`$ is a subalgebra of $`V`$ containing $`A`$ and $`\mathrm{𝟏}`$. Since $`A`$ generates $`V`$, we must have $`K=V`$. This concludes the proof. $`\mathrm{}`$ With Lemma 3.19 we immediately have: ###### Proposition 3.20 Let $`A`$ be a generalized vertex pre-algebra of parafermions on $`W`$ and let $`D_W`$ be a grading-preserving endomorphism of $`W`$ such that $`[D_W,a(z)]={\displaystyle \frac{d}{dz}}a(z)\text{ for }a(z)A.`$ (3.56) Denote by $`V`$ the generalized vertex algebra generated by $`A`$. Then on $`W`$, $`[D_W,Y(v,z)]=Y(D(v),z)={\displaystyle \frac{d}{dz}}Y(v,z)\text{ for }vV.\mathrm{}`$ (3.57) ## 4 Generalized vertex algebras associated with $`Z`$-algebras In this section we briefly recall from \[LW1-3\] and \[LP1-2\] the fundamental results about $`Z`$-operators and we then show that the vacuum space $`\mathrm{\Omega }_{M(\mathrm{},0)}`$ of the generalized Verma module $`M(\mathrm{},0)`$ for an affine Lie algebra $`\widehat{g}`$ with respect to the homogeneous Heisenberg subalgebra has a canonical generalized vertex algebra structure. We also show that for any highest weight $`\widehat{g}`$-module $`W`$ of level $`\mathrm{}`$, $`\mathrm{\Omega }_W`$ is an $`\mathrm{\Omega }_{M(\mathrm{},0)}`$-module. Let $`g`$ be a finite-dimensional simple Lie algebra, $`𝐡`$ be a Cartan subalgebra, $`\mathrm{\Phi }`$ be the set of roots and $`Q`$ be the root lattice. Let $`,`$ be the normalized killing form such that $`\alpha ,\alpha =2`$ for a long root $`\alpha `$. Using $`,`$ we identify $`𝐡^{}`$ with $`𝐡`$. Let $`\widehat{g}`$ be the affine Lie algebra: $`\widehat{g}=g𝐂[𝐭,𝐭^\mathrm{𝟏}]\mathrm{𝐂𝐜},`$ (4.1) where $`[at^m,bt^n]=[a,b]t^{m+n}+ma,b\delta _{m+n,0}c,`$ (4.2) $`[\widehat{g},c]=0`$ (4.3) for $`a,bg,m,n𝐙`$. Following the tradition, we also use $`a(n)`$ for $`at^n`$. For $`n𝐙`$, we denote $`g(n)=\{a(n)\|ag\}.`$ (4.4) For $`ag`$, define the generating series $`a(z)={\displaystyle \underset{n𝐙}{}}(at^n)z^{n1}\widehat{g}[[z,z^1]].`$ (4.5) Set $`\widehat{𝐡}^+=𝐡t𝐂[𝐭],\widehat{𝐡}^{}=𝐡𝐭^\mathrm{𝟏}𝐂[𝐭^\mathrm{𝟏}],`$ (4.6) subalgebras of $`\widehat{𝐡}=𝐡𝐂[𝐭,𝐭^\mathrm{𝟏}]+\mathrm{𝐂𝐜}`$. Then $`\widehat{𝐡}_𝐙=\widehat{𝐡}^++\widehat{𝐡}^{}+\mathrm{𝐂𝐜}`$ (4.7) is a Heisenberg subalgebra of $`\widehat{g}`$. For $`\mathrm{}𝐂^\times `$, let $`M(\mathrm{})`$ be the standard irreducible $`\widehat{𝐡}_𝐙`$-module with $`c`$ acting as scalar $`\mathrm{}`$. We may also consider $`M(\mathrm{})`$ as an $`\widehat{𝐡}`$-module with the action of $`𝐡`$ being zero. ###### Definition 4.1 For $`\mathrm{}𝐂`$, we define a category $`𝒞_{\mathrm{}}`$ whose objects are level-$`\mathrm{}`$ $`\widehat{g}`$-modules $`W`$ which are $`𝐡`$-weight modules satisfying the condition that for every $`wW`$, $`g(n)w=0`$ for $`n`$ sufficiently large and $`dimU(\widehat{𝐡}^+)w<\mathrm{}`$. It follows from \[LW3\] and \[K\] that each $`W`$ from $`𝒞_{\mathrm{}}`$ with $`\mathrm{}0`$ is a completely reducible $`\widehat{𝐡}`$-module. For $`W𝒞_{\mathrm{}}`$, set $`\mathrm{\Omega }_W=\{wW\|h(n)w=0\text{ for }h𝐡,n𝐙_+\}.`$ (4.8) Since $`[h(0),h^{}(m)]=0`$ for $`h,h^{}𝐡,m𝐙`$, $`h(0)`$ preserves $`\mathrm{\Omega }_W`$. With $`\widehat{g}`$, as an $`𝐡`$-module, being naturally $`𝐡(=𝐡^{})`$-graded we have $`\mathrm{\Omega }_W={\displaystyle \underset{\alpha 𝐡}{}}\mathrm{\Omega }_W^\alpha ,`$ (4.9) where $`\mathrm{\Omega }_W^\alpha =\{w\mathrm{\Omega }_W\|h(0)w=\alpha ,hw\text{ for }h𝐡\}`$. For $`h𝐡`$, set (\[LW1\]) $`E^\pm (h,z)=\mathrm{exp}\left({\displaystyle \underset{\pm n𝐙_+}{}}{\displaystyle \frac{h(n)}{n}}z^n\right)U(\widehat{g})[[z^1]].`$ (4.10) For $`ag_\alpha ,\alpha \mathrm{\Phi }`$, we define $`Z(a,z)=E^{}({\displaystyle \frac{1}{\mathrm{}}}\alpha ,z)a(z)E^+({\displaystyle \frac{1}{\mathrm{}}}\alpha ,z),`$ (4.11) a formal object. (It is an element of $`\overline{U(\widehat{g})}[[z,z^1]]`$, where $`\overline{U(\widehat{g})}`$ is a certain formal completion of $`U(\widehat{g})`$.) For every $`W𝒞_{\mathrm{}}`$, $`Z(v,z)`$ is a well defined element of $`(\mathrm{End}W)[[z,z^1]]`$. Then we have (\[LW4\], \[LP2\]): ###### Proposition 4.2 Let $`\mathrm{}𝐂^\times ,𝐖𝒞_{\mathrm{}}`$. For $`h𝐡,ug_\alpha ,vg_\beta ,\alpha ,\beta \mathrm{\Phi }`$, on $`W`$, $`[h(0),Z(u,z)]=\alpha ,hZ(u,z),`$ (4.12) $`[h(n),Z(u,z)]=0\text{ for }n0,`$ (4.13) $$(1z_2/z_1)^{\alpha ,\beta /\mathrm{}}Z(u,z_1)Z(v,z_2)(1z_1/z_2)^{\alpha ,\beta /\mathrm{}}Z(v,z_2)Z(u,z_1)=$$ $$=\{\begin{array}{c}z_1^1\delta (z_2/z_1)Z([u,v],z_2)\text{ if }\alpha +\beta 0,\\ z_1^1\delta (z_2/z_1)[u,v]z_2^1+\mathrm{}u,v\frac{}{z_2}z_1^1\delta (z_2/z_1)\text{ if }\alpha +\beta =0.\end{array}$$ (4.14) It follows immediately from (4.13) that $`Z(v,z)`$ maps $`\mathrm{\Omega }_W`$ to $`\mathrm{\Omega }_W[[z,z^1]]`$. Set $`Z(v,z)={\displaystyle \underset{n𝐙}{}}Z(v,n)z^n.`$ (4.15) ###### Definition 4.3 Following \[LP2\] we define a category $`𝒵_{\mathrm{}}`$ (which was denoted by $`P_{\mathrm{}}`$ in \[LP2\]) whose objects are $`𝐡`$-weight modules $`U`$ equipped with a family of operators $`Z_U(a,m)`$ (linear in $`a`$) on $`U`$ for $`ag_\alpha ,\alpha \mathrm{\Phi },m𝐙`$ such that $`Z_U(a,z)wU((z))`$ and such that (4.12) and (4.14) hold for $`Z_W(u,z)`$ in place of $`Z(u,z)`$. Clearly, we have a functor $`\mathrm{\Omega }`$ from $`𝒞_{\mathrm{}}`$ to $`𝒵_{\mathrm{}}`$. Conversely, given $`U𝒵_{\mathrm{}}`$, we set $`E(U)=M(\mathrm{})U\text{ as a vector space.}`$ (4.16) View $`E(U)`$ as a natural $`\widehat{𝐡}_𝐙`$-module with $`h(n)`$ (for $`n0`$) acting on the first factor. Then define an action of $`\widehat{g}`$ by $`c\mathrm{},`$ (4.17) $`h1h\text{ for }h𝐡,`$ (4.18) $`a(z)E^{}({\displaystyle \frac{1}{\mathrm{}}}\alpha ,z)E^+({\displaystyle \frac{1}{\mathrm{}}}\alpha ,z)Z_U(a,z)\text{ for }ag_\alpha ,\alpha \mathrm{\Phi }.`$ (4.19) It was proved in \[LP2\] that $`E(U)`$ is a $`\widehat{g}`$-module in $`𝒞_{\mathrm{}}`$. Furthermore, we have (\[LW4\], \[LP2\]): ###### Proposition 4.4 Let $`\mathrm{}𝐂^\times `$. Then the functors $`\mathrm{\Omega }:W\mathrm{\Omega }_W\text{and }E:UE(U)`$ (4.20) are exact and they define equivalences between the categories $`𝒞_{\mathrm{}}`$ and $`𝒵_{\mathrm{}}`$. In particular, $`W`$ is irreducible in $`𝒞_{\mathrm{}}`$ if and only if $`\mathrm{\Omega }_W`$ is irreducible in $`𝒵_{\mathrm{}}`$. For $`\lambda 𝐡`$, let $`M(\mathrm{},\lambda )`$ be the Verma $`\widehat{g}`$-module. In view of the universal property for $`M(\mathrm{},\lambda )`$, with Proposition 4.4 we immediately have: ###### Corollary 4.5 Let $`\mathrm{}𝐂,\lambda 𝐡`$ and let $`v`$ be a (nonzero) highest weight vector of $`M(\mathrm{},\lambda )`$. Let $`U𝒵_{\mathrm{}}`$ and let $`eU`$ satisfying the following conditions: $`he=h,\lambda e\text{ for }h𝐡,`$ (4.21) $`Z_U(u,z)ez^1U[[z]]\text{ for }ug_\alpha ,\alpha \mathrm{\Phi },`$ (4.22) $`Z_U(v,z)eU[[z]]\text{ for }vg_\beta ,\beta \mathrm{\Phi }_+.`$ (4.23) Then there exists a unique morphism in $`𝒵_{\mathrm{}}`$ from $`\mathrm{\Omega }_{M(\mathrm{},\lambda )}`$ to $`U`$ sending $`v`$ to $`e`$. $`\mathrm{}`$ ###### Definition 4.6 For $`ag_\alpha ,\alpha \mathrm{\Phi }`$, we define $`\psi (a,z)=Z(a,z)z^{\frac{1}{\mathrm{}}\alpha (0)}=E^{}({\displaystyle \frac{1}{\mathrm{}}}\alpha ,z)a(z)E^+({\displaystyle \frac{1}{\mathrm{}}}\alpha ,z)z^{\frac{1}{\mathrm{}}\alpha (0)}.`$ (4.24) Then $`[h(n),\psi (a,z)]=0\text{ for }h𝐡,n0,`$ (4.25) hence $`\psi (a,z)`$ maps $`\mathrm{\Omega }_W`$ to $`\mathrm{\Omega }_W\{z\}`$ for $`W𝒞_{\mathrm{}}`$. Note that (4.12) amounts to $`z^{h(0)}Z(u,z_1)=Z(u,z_1)z^{\alpha ,h+h(0)}`$ (4.26) for $`h𝐡,ug_\alpha ,\alpha \mathrm{\Phi }`$. We have the following reformulation of Proposition 4.2 in terms of $`\psi `$-operators: ###### Proposition 4.7 Let $`ug_\alpha ,vg_\beta ,\alpha ,\beta \mathrm{\Phi },\mathrm{}𝐂^\times `$ and $`W𝒵_{\mathrm{}}`$. On $`W`$, $`[h(0),\psi (u,z)]=\alpha ,h\psi (u,z),`$ (4.27) $`[h(n),\psi (u,z)]=0\text{ for }h𝐡,n0,`$ (4.28) $$(z_1z_2)^{\alpha ,\beta /\mathrm{}}\psi (u,z_1)\psi (v,z_2)(z_2z_1)^{\alpha ,\beta /\mathrm{}}\psi (v,z_2)\psi (u,z_1)=$$ $`=\{\begin{array}{c}z_1^1\delta (z_2/z_1)\psi ([u,v],z_2)(z_2/z_1)^{\frac{1}{\mathrm{}}\alpha (0)}\text{ if }\alpha +\beta 0,\\ \mathrm{}u,v\frac{}{z_2}\left(z_1^1\delta (z_2/z_1)(z_2/z_1)^{\frac{1}{\mathrm{}}\alpha (0)}\right)\text{ if }\alpha +\beta =0.\end{array}`$ (4.31) Proof. The first two identities are obvious. Using (4.26) we obtain $$(z_1z_2)^{\alpha ,\beta /\mathrm{}}\psi (u,z_1)\psi (v,z_2)(z_2z_1)^{\alpha ,\beta /\mathrm{}}\psi (v,z_2)\psi (u,z_1)=$$ $`=\{\begin{array}{c}z_1^1\delta (z_2/z_1)\psi ([u,v],z_2)(z_2/z_1)^{\frac{1}{\mathrm{}}\alpha (0)}\text{ if }\alpha +\beta 0,\\ \left(z_1^1\delta (z_2/z_1)[u,v]z_2^1+\mathrm{}u,v\frac{}{z_2}z_1^1\delta (z_2/z_1)\right)(z_2/z_1)^{\frac{1}{\mathrm{}}\alpha (0)}\text{ if }\alpha +\beta =0.\end{array}`$ (4.34) It remains to consider the case $`\alpha +\beta =0`$. Using the fact that $`[u,v]=u,v\alpha `$ and $`\delta (z)z^m=\delta (z)`$ for $`m𝐙`$ we obtain $`\left(z_1^1\delta (z_2/z_1)[u,v]z_2^1+\mathrm{}u,v{\displaystyle \frac{}{z_2}}z_1^1\delta (z_2/z_1)\right)(z_2/z_1)^{\frac{1}{\mathrm{}}\alpha (0)}`$ (4.35) $`=`$ $`\mathrm{}u,v{\displaystyle \frac{}{z_2}}\left(z_1^1\delta (z_2/z_1)(z_2/z_1)^{\frac{1}{\mathrm{}}\alpha (0)}\right).`$ This completes the proof.$`\mathrm{}`$ It is a simple fact (see for example \[Li2\]) that $`(z_1z_2)^m\left({\displaystyle \frac{}{z_2}}\right)^nz_2^1\delta (z_1/z_2)=0`$ (4.36) for $`m,n𝐙`$ with $`m>n0`$. Then using Proposition 4.7 we get $`(z_1z_2)^2\left((z_1z_2)^{\alpha ,\beta /\mathrm{}}\psi (u,z_1)\psi (v,z_2)(z_2z_1)^{\alpha ,\beta /\mathrm{}}\psi (v,z_2)\psi (u,z_1)\right)=0,`$ (4.37) where $`u,v`$ are as in Proposition 4.7. Then for any $`U𝒵_{\mathrm{}}`$, e.g., $`U=\mathrm{\Omega }_W`$ for some $`W𝒞_{\mathrm{}}`$, $`\psi (u,z)`$ for $`ug_\alpha ,\alpha \mathrm{\Phi }`$ linearly span a generalized vertex pre-algebra of parafermion operators on $`U`$, which by Theorem 3.16 generates a canonical generalized vertex algebra inside $`(\mathrm{End}U)\{z\}`$ with $`G=𝐡,c(,)=1,(,)=\frac{1}{\mathrm{}},`$. ###### Lemma 4.8 Let $`U𝒵_{\mathrm{}}`$ and let $`V`$ be the generalized vertex algebra generated by $`\psi (u,z)`$ for $`ug_\alpha ,\alpha \mathrm{\Phi }`$ inside $`(\mathrm{End}U)\{z\}`$. Then $`V`$ is a natural object of $`𝒵_{\mathrm{}}`$ where $`h\varphi (z)=[h,\varphi (z)],`$ (4.38) $`Z_V(u,z_0)=Y_V(\psi (u,z),z_0)z_0^{\frac{1}{\mathrm{}}\alpha (0)}`$ (4.39) for $`h𝐡,\varphi (z)V,ug_\alpha ,\alpha \mathrm{\Phi }`$. Proof. Since $`U`$ is an $`𝐡`$-weight module, $`\mathrm{End}U`$ is a natural $`𝐡`$-module where $`hf=[h,f](=hffh)\text{ for }h𝐡,f\mathrm{End}U.`$ (4.40) Then $`(\mathrm{End}U)\{z\}`$ is a natural $`𝐡`$-module. Since the generators $`\psi (a,z)`$ for $`ag_\alpha `$ of $`V`$ are $`𝐡`$-eigenvectors (recall (4.27)), using the proof of Lemma 3.19 we can easily show that $`V`$ is an $`𝐡`$-weight module and (4.27) holds on $`V`$. Note that $`U`$ is a faithful $`V`$-module. Then it follows immediately from Lemma 2.8 that (4.31) holds on $`V`$. This shows that $`V`$ is a natural object of $`𝒵_{\mathrm{}}`$. $`\mathrm{}`$ Let $`\mathrm{}𝐂^\times `$. Consider the generalized Verma $`\widehat{g}`$-module $`M(\mathrm{},0)=U(\widehat{g})_{U(g𝐂[𝐭]+\mathrm{𝐂𝐜})}𝐂_{\mathrm{}},`$ (4.41) where $`𝐂_{\mathrm{}}=𝐂`$ as a vector space and $`g𝐂[𝐭]`$ acts as zero on $`𝐂_{\mathrm{}}`$ and $`c`$ acts as $`\mathrm{}`$. Denote by $`\mathrm{𝟏}`$ the highest weight vector $`11`$ of $`M(\mathrm{},0)`$. Let $`L(\mathrm{},0)`$ be the (unique) irreducible quotient module with $`\mathrm{𝟏}`$ as a fixed highest weight vector. Identify $`g`$ as a subspace of $`M(\mathrm{},0)`$ and $`L(\mathrm{},0)`$ through $`aa(1)\mathrm{𝟏}`$. Then we have $`g_\alpha \mathrm{\Omega }_{M(\mathrm{},0)}^\alpha \mathrm{\Omega }_{M(\mathrm{},0)}\text{ for }\alpha \mathrm{\Phi }.`$ (4.42) ###### Theorem 4.9 Let $`\mathrm{}𝐂^\times `$ and $`V=M(\mathrm{},0)`$ or $`L(\mathrm{},0)`$. Then there exists a unique generalized vertex algebra structure $`Y_\mathrm{\Omega }`$ on $`\mathrm{\Omega }_V`$ with $`G=Q,c(,)=1`$ and $`(,)=,/\mathrm{}`$ such that $`Y_\mathrm{\Omega }(\mathrm{𝟏},z)=1`$ and $`Y_\mathrm{\Omega }(a,z)=\psi (a,z)`$ for $`ug_\alpha ,\alpha \mathrm{\Phi }`$. Furthermore, $`\mathrm{\Omega }_V`$ is generated by $`g_\alpha `$ ($`\alpha \mathrm{\Phi }`$). Proof. Clearly, $`V`$ is $`Q`$-graded. Then we take $`G=Q`$, $`(,)=,/\mathrm{}`$ and $`c(,)=1`$. Let $`A`$ be the linear span of $`\psi (a,z)`$ for $`ag_\alpha ,\alpha \mathrm{\Phi }`$. It follows from Proposition 4.7 and (4.36) that $`A`$ is a generalized vertex pre-algebra. Since $`\mathrm{\Omega }_V`$ is generated from $`\mathrm{𝟏}`$ by all the components of $`Z(a,z)=\psi (a,z)z^{\alpha (0)/\mathrm{}}`$, there exists a unique generalized vertex algebra structure $`Y_\mathrm{\Omega }`$ on $`\mathrm{\Omega }_V`$ with the required conditions.$`\mathrm{}`$ ###### Proposition 4.10 Let $`\mathrm{}𝐂^\times `$ and let $`U𝒵_{\mathrm{}}`$. Then $`U`$ is a natural $`\mathrm{\Omega }_{M(\mathrm{},0)}`$-module. In particular, $`\mathrm{\Omega }_W`$ is an $`\mathrm{\Omega }_{M(\mathrm{},0)}`$-module for $`W𝒞_{\mathrm{}}`$. Proof. Set $`M=\mathrm{\Omega }_{M(\mathrm{},0)}U`$, an object of $`𝒵_{\mathrm{}}`$. Let $`V`$ be the generalized vertex algebra generated by $`\psi (v,z)`$ for $`vg_\alpha ,\alpha \mathrm{\Phi }`$ inside $`(\mathrm{End}\mathrm{\Omega }_{M(\mathrm{},0)})\{z\}(\mathrm{End}M)\{z\}.`$ Then $`M`$ is a $`V`$-module with $`\mathrm{\Omega }_{M(\mathrm{},0)}`$ and $`U`$ as submodules. From Proposition 3.19, there is a $`V`$-homomorphism $`f`$ from $`V`$ onto $`\mathrm{\Omega }_{M(\mathrm{},0)}`$, which maps $`I(z)`$ to $`\mathrm{𝟏}`$. In view of Lemma 4.8, $`V`$ is a natural $`𝐡`$-module in $`𝒵_{\mathrm{}}`$. It follows from Corollary 4.5 that $`f`$ is a linear isomorphism, hence $`V=\mathrm{\Omega }_{M(\mathrm{},0)}`$. Thus $`M`$ is an $`\mathrm{\Omega }_{M(\mathrm{},0)}`$-module. Therefore, $`U`$ is an $`\mathrm{\Omega }_{M(\mathrm{},0)}`$-module.$`\mathrm{}`$ We define an $`\mathrm{\Omega }_{M(\mathrm{},0)}`$-$`𝐡`$-module to be an $`\mathrm{\Omega }_{M(\mathrm{},0)}`$-module and an $`𝐡`$-weight module such that (4.27) holds. In view of Proposition 4.10 and Lemma 4.8 we immediately have: ###### Corollary 4.11 The category $`𝒵_{\mathrm{}}`$ is canonically isomorphic to the category of $`\mathrm{\Omega }_{M(\mathrm{},0)}`$-$`𝐡`$-modules. $`\mathrm{}`$ ###### Remark 4.12 Note that for a positive integer $`\mathrm{}`$, the generalized vertex operator algebra $`\mathrm{\Omega }_{L(\mathrm{},0)}^B`$ constructed in \[DL2\] is a quotient algebra of $`\mathrm{\Omega }_{M(\mathrm{},0)}`$. This is recently studied in \[Li4\] from a different point of view.
warning/0006/hep-lat0006029.html	ar5iv	text	# 1 Introduction ## 1 Introduction The determination of non-leptonic kaon-decay amplitudes from Lattice QCD remains a challenging task. However, recently there have been several developments which may lead to substantial progress in this field, ranging from new ideas on how to cope with chiral symmetry on the lattice to a sizable increase of computer power available for the necessary numerical computations. Still, there are important theoretical difficulties afflicting the determination of the relevant weak matrix elements, which are a consequence of the fact that the final states contain more than one strongly interacting particle (the pions). This is formalized in what is sometimes called the Maiani–Testa theorem , which says that it is not possible to extract the physical matrix elements with the correct kinematics from the asymptotic behavior in time of the euclidean correlation functions accessible to numerical computation. There are several old and new ideas on the market on how to deal with this situation, which divide into three groups. First, one may determine the matrix elements with a kaon in the initial state and two (or more) pions in the final state for an unphysical choice of external momenta. From the large-time behavior of the euclidean correlation function $$C(t_2,t_1)=0\|\pi _1(t_2)\pi _2(t_2)𝒪_{\mathrm{weak}}(t_1)K(0)\|0,$$ (1.1) with the kaon at rest, one obtains the matrix element for $`\stackrel{}{p}\stackrel{}{p}_{\pi _1}=\stackrel{}{p}_{\pi _2}=0`$ instead of the physical $`\|\stackrel{}{p}\|=(1/2)\sqrt{m_K^24m_\pi ^2}`$, i.e. energy is not conserved . The idea is then to use chiral perturbation theory (ChPT) in order to correct for this unphysical choice of momenta . The most recent computation of the $`\mathrm{\Delta }I=3/2`$ $`K2\pi `$ matrix element using this method can be found in Ref. . In this computation, all meson masses were taken degenerate and the quenched approximation was used. Adjustments for all these unphysical effects, which also include power-like finite-volume effects coming from pion rescattering diagrams, were made using one-loop ChPT . For recent ideas on choosing $`m_K=2m_\pi `$ for the lattice computation (for which $`\stackrel{}{p}=0`$ does conserve energy), see Refs. . For the $`\mathrm{\Delta }I=1/2`$ case, the situation is more complicated for a number of reasons. Here we only mention (since this is less well known) that a quenched or partially quenched computation appears to be afflicted by “enhanced finite-volume effects,” which do not occur for the $`\mathrm{\Delta }I=3/2`$ case. This problem appears at one–loop in quenched ChPT , but not much is known about this effect beyond one loop. We are investigating this issue. For a nice review of many other issues, see Ref. . A second idea was very recently proposed in Ref. , where it was shown how the matrix elements of interest can in principle be determined from finite-volume correlation functions without analytic continuation. The energy-conserving amplitude is obtained by tuning the spatial volume such that the first excited level of the two-pion final state has an energy equal to the kaon mass. With sufficient accuracy to determine the lowest excited levels, the finite-volume matrix element may then be computed on the lattice, and subsequently be converted into the physical infinite-volume amplitude. For this, it is obviously necessary to choose meson masses such that $`2m_\pi <m_K`$ (as well as $`m_K<4m_\pi `$, so that the final-state pions are in the elastic regime). Again, if such a computation is done in a (partially) quenched setting one might expect that enhanced finite-volume effects could also occur with this method. A third idea is based on the observation that, if one needs ChPT anyway in order to convert an unphysical matrix element into a physical one, one might as well choose the unphysical matrix element as simple as possible. Chiral symmetry relates the $`K2\pi `$ matrix elements of interest to the simpler $`K\pi `$ and $`K`$ vacuum ($`K0`$) matrix elements of the same weak operators which mediate non-leptonic kaon decay . Advantages of this approach are that there are no strongly interacting particles in the final state, and that lattice computations of these simpler matrix elements may be less difficult. The first advantage is, in a sense, not really an advantage if one wishes to convert the results of a lattice computation into a calculation of the non-leptonic kaon decay rates, because final-state interactions will still have to be taken into account. However, this method does avoid all the unphysical effects, such as power-like or even enhanced finite-volume effects, associated with the multi-pion final state. Formulated in another way, with this method the simplest possible matrix elements (in this case $`K\pi `$ and $`K0`$) are used to obtain the relevant weak low-energy constants (LECs) of the weak effective lagrangian. Using ChPT, these can then be converted into estimates of the kaon-decay rates. A key question is which order in ChPT will be needed in order to carry out such a program. At tree level (i.e. $`O(p^2)`$), only three LECs come into play, but at one loop (i.e. $`O(p^4)`$) many more LECs contribute to all relevant matrix elements . In fact, from the analysis reported in this paper as well as from previous work it is clear that tree-level ChPT is not enough . In addition (as we will demonstrate), not all $`O(p^4)`$ LECs needed for $`K2\pi `$ decays can be obtained from $`K\pi `$ and $`K0`$ matrix elements. However, as we will advocate in this paper, it may be possible to determine at least the $`O(p^2)`$ LECs from a lattice computation, taking one-loop ChPT effects into account. A reliable, first-principle determination of the $`O(p^2)`$ octet and 27-plet LECs would clearly be interesting by itself. Moreover, phenomenological estimates of these LECs, based on a one-loop ChPT analysis of experimental data are available , making a direct comparison possible. In this paper we present an analysis of $`K\pi `$ and $`K0`$ amplitudes in one-loop ChPT, with the above described philosophy in mind. For $`K\pi `$ we choose our valence quark masses to be degenerate, thus conserving energy for this case. The analysis is performed in partially quenched ChPT , with an arbitrary number of degenerate sea quarks. We present the results for these matrix elements in terms of the quark masses. In Sect. 2, we list and discuss all $`O(p^2)`$ and $`O(p^4)`$ operators needed for our calculation, including those containing the $`\eta ^{}`$ meson. In Sect. 3, we discuss the role of the $`\eta ^{}`$ in partially quenched QCD in some more detail than has been done so far in the literature. This section can be skipped if one is only interested in results. In Sect. 4, we give complete one-loop expressions for the octet and 27-plet $`K\pi `$ and $`K0`$ matrix elements, including contributions from $`O(p^4)`$ operators, organized by subsection. In Subsect. 4.1 partially quenched results for $`N`$ degenerate sea quarks are presented, which are valid also in the case that the meson made out of sea quarks (the “sea meson”) is not light compared to the $`\eta ^{}`$ (a realistic situation in actual lattice computations). In Subsect. 4.2 we specialize to the case that the sea meson is light compared to the $`\eta ^{}`$. Subsect. 4.3 contains the completely quenched results, obtained by setting $`N=0`$ and keeping the $`\eta ^{}`$. For completeness, we include the fully unquenched results, with non-degenerate sea-quark masses for the $`K0`$ matrix elements, in Subsect. 4.4. In Sect. 5 we present a detailed discussion of the results, including the role of $`O(p^4)`$ operators and numerical examples for typical choices of the parameters. The last section contains our conclusions. ## 2 Definition of operators Partially quenched QCD may be defined by separately introducing valence- and sea-quark fields, each with their own mass. The valence quarks are quenched by introducing for each valence quark a “ghost” quark, which has the same mass and quantum numbers, but opposite statistics . This, in effect, removes the valence-quark determinant from the QCD partition function. We will consider a theory with $`n`$ quarks, of which $`N`$ are sea quarks, and $`nN3`$ valence quarks. This requires $`nN`$ ghost quarks, with masses equal to those of the valence quarks. We will consider valence quarks with arbitrary masses $`m_1,\mathrm{},m_{nN}`$, and degenerate sea quarks, all with mass $`m_S`$. The relevant chiral symmetry group is the graded group SU($`n\|nN`$)$`{}_{L}{}^{}`$SU($`n\|nN`$)<sub>R</sub> . Fully quenched QCD arises as a special case of this construction by taking $`N=0`$ , or equivalently, when the sea quarks are decoupled by taking $`m_S\mathrm{}`$. The euclidean low-energy effective lagrangian which mediates non-leptonic weak transitions with $`\mathrm{\Delta }S=1`$ is given by $$_{\mathrm{\Delta }S=1}=_2+_2^\eta ^{}+_4+\mathrm{},$$ (2.1) where the dots denote higher order terms in the chiral expansion. The $`O(p^2)`$ lagrangian $`_2`$ contains three terms <sup>1</sup><sup>1</sup>1In the analysis of CP conserving weak amplitudes one can safely disregard $`O(e^2p^0)`$ terms induced by electroweak interactions. $$_2=\alpha _1^8\text{str}(\mathrm{\Lambda }L_\mu L_\mu )+\alpha _2^8\text{str}(\mathrm{\Lambda }X_+)+\alpha ^{27}T_{kl}^{ij}(L_\mu )_i^k(L_\mu )_j^l+\text{h.c.},$$ (2.2) where str denotes the supertrace in flavor space. Note that the supertraces become normal traces, $`\text{str}\text{tr}`$, in the case $`N=n`$. The terms with couplings $`\alpha _1^8,\alpha _2^8`$ transform as $`(8_L,1_R)`$, while the term with coupling $`\alpha ^{27}`$ transforms as $`(27_L,1_R)`$. The order $`p^2`$ lagrangian $`_2^\eta ^{}`$ will be discussed toward the end of this section. The $`O(p^4)`$ lagrangian can be written as $$_4=\frac{1}{(4\pi f)^2}\left(\underset{i}{}\beta _i^8𝒪_i^8+\underset{i}{}\stackrel{~}{\beta }_i^8\stackrel{~}{𝒪}_i^8+\underset{i}{}\beta _i^{27}𝒪_i^{27}\right),$$ (2.3) with $`(8_L,1_R)`$ $`O(p^4)`$ operators $`𝒪_i^8`$, $`\stackrel{~}{𝒪}_i^8`$ and $`(27_L,1_R)`$ $`O(p^4)`$ operators $`𝒪_i^{27}`$ . The operators $`\stackrel{~}{𝒪}_i^8`$ denote total-derivative operators; they do not contribute to energy-momentum conserving matrix elements. However, they do contribute to the $`K0`$ matrix element, which does not conserve energy for non-degenerate quark masses. Note that there are no 27-plet total-derivative operators that contribute to the matrix elements considered in this paper, to $`O(p^4)`$. The fields entering the weak operators are defined as follows $`L_\mu `$ $`=`$ $`i\mathrm{\Sigma }_\mu \mathrm{\Sigma }^{},`$ (2.4) $`X_\pm `$ $`=`$ $`2B_0(\mathrm{\Sigma }M^{}\pm M\mathrm{\Sigma }^{}),`$ where $`B_0`$ is the parameter $`B_0`$ of Ref. and $`B_0=4v/f^2`$ in the notation of Ref. . The unitary field $`\mathrm{\Sigma }`$ is defined in terms of the hermitian field $`\mathrm{\Phi }`$ describing the Goldstone meson multiplet as $$\mathrm{\Sigma }=\text{exp}(2i\mathrm{\Phi }/f),$$ (2.5) where $`f`$ is the bare pion-decay constant, normalized such that $`f_\pi =132`$ MeV. Note that $`L_\mu `$ and $`X_\pm `$ transform as $`(8_L,1_R)`$ under (valence-flavor) SU(3)$`{}_{L}{}^{}`$SU(3)<sub>R</sub>. The matrix $`\mathrm{\Lambda }`$ in the lagrangian (2.2) picks out the $`\mathrm{\Delta }S=1`$, $`\mathrm{\Delta }D=1`$ part of the octet operators, all with $`\mathrm{\Delta }I=1/2`$: $$\mathrm{\Lambda }_j^i=\delta ^{i3}\delta _{j2},$$ (2.6) where $`i,j=1,2,3\mathrm{}`$ (or $`u,d,s\mathrm{}`$) denote valence flavors. The tensor $`T_{kl}^{ij}`$ projects onto the $`\mathrm{\Delta }I=1/2`$ part of the $`(27_L,1_R)`$ operator, with nonzero components $`T_{12}^{13}`$ $`=`$ $`T_{12}^{31}=T_{21}^{13}=T_{21}^{31}={\displaystyle \frac{1}{2}},`$ (2.7) $`T_{22}^{23}`$ $`=`$ $`T_{22}^{32}=1,`$ $`T_{32}^{33}`$ $`=`$ $`T_{23}^{33}={\displaystyle \frac{3}{2}},`$ or onto the $`\mathrm{\Delta }I=3/2`$ part, with nonzero components $`T_{12}^{13}`$ $`=`$ $`T_{12}^{31}=T_{21}^{13}=T_{21}^{31}={\displaystyle \frac{1}{2}},`$ (2.8) $`T_{22}^{23}`$ $`=`$ $`T_{22}^{32}={\displaystyle \frac{1}{2}}.`$ The term with coupling $`\alpha _2^8`$ is known as the “weak mass term,” and mediates the $`K0`$ transition at tree level. Its odd-parity part, which in principle can also contribute to the octet $`K\pi \pi `$ amplitude, is proportional to $`m_sm_d`$. For $`m_sm_d`$ the weak mass term is a total derivative , and therefore does not contribute to any energy-momentum-conserving physical matrix element, like $`K\pi \pi `$. Instead, the $`K0`$ and, for $`M_KM_\pi `$, $`K\pi `$ matrix elements do not conserve energy, and therefore the weak mass term does contribute to both of them. For $`m_s=m_d`$ the weak mass term is not a total derivative, so that it contributes also to the $`K\pi `$ matrix element with $`M_K=M_\pi `$. Hence, in order to determine the octet coupling $`\alpha _1^8`$ from a computation of the $`K\pi `$ matrix element, another quantity such as the $`K0`$ matrix element is always needed, in order to eliminate the dependence on $`\alpha _2^8`$. At order $`p^4`$ there are eight $`(8_L,1_R)`$ operators and six $`(27_L,1_R)`$ operators which can contribute to $`K0`$, $`K\pi `$ and $`K\pi \pi `$ matrix elements. The octet operators can be written as follows $`𝒪_1^8`$ $`=`$ $`\text{str}(\mathrm{\Lambda }X_+X_+),`$ (2.9) $`𝒪_2^8`$ $`=`$ $`\text{str}(\mathrm{\Lambda }X_+)\text{str}(X_+),`$ $`𝒪_3^8`$ $`=`$ $`\text{str}(\mathrm{\Lambda }X_{}X_{}),`$ $`𝒪_5^8`$ $`=`$ $`\text{str}(\mathrm{\Lambda }[X_+,X_{}]),`$ $`𝒪_{10}^8`$ $`=`$ $`\text{str}(\mathrm{\Lambda }\{X_+,L_\mu L_\mu \}),`$ $`𝒪_{11}^8`$ $`=`$ $`\text{str}(\mathrm{\Lambda }L_\mu X_+L_\mu ),`$ $`𝒪_{13}^8`$ $`=`$ $`\text{str}(\mathrm{\Lambda }X_+)\text{str}(L_\mu L_\mu ),`$ $`𝒪_{15}^8`$ $`=`$ $`\text{str}(\mathrm{\Lambda }[X_{},L_\mu L_\mu ]),`$ while the 27-plet operators are $`𝒪_1^{27}`$ $`=`$ $`T_{kl}^{ij}(X_+)_i^k(X_+)_j^l,`$ (2.10) $`𝒪_2^{27}`$ $`=`$ $`T_{kl}^{ij}(X_{})_i^k(X_{})_j^l,`$ $`𝒪_4^{27}`$ $`=`$ $`T_{kl}^{ij}(L_\mu )_i^k(\{L_\mu ,X_+\})_j^l,`$ $`𝒪_5^{27}`$ $`=`$ $`T_{kl}^{ij}(L_\mu )_i^k([L_\mu ,X_{}])_j^l,`$ $`𝒪_6^{27}`$ $`=`$ $`T_{kl}^{ij}(X_+)_i^k(L_\mu L_\mu )_j^l,`$ $`𝒪_7^{27}`$ $`=`$ $`T_{kl}^{ij}(L_\mu )_i^k(L_\mu )_j^l\text{str}(X_+).`$ These operators are the same as those in Ref. , apart from the replacement $`\text{tr}\text{str}`$. For the energy-momentum non-conserving matrix elements the only total-derivative term that is needed is $$\stackrel{~}{𝒪}_1^8=i_\mu \text{str}(\mathrm{\Lambda }[L_\mu ,X_+]).$$ (2.11) In general, in the (partially) quenched formulation of the effective theory one needs to keep the $`\eta ^{}`$ , defined as the SU($`n\|nN`$)$`{}_{L}{}^{}`$SU($`n\|nN`$)<sub>R</sub> invariant $$\eta ^{}=\text{str}(\mathrm{\Phi }).$$ (2.12) Note that this normalization differs from the one in Ref. , but is more convenient in keeping track of $`N`$ dependence. The presence of the $`\eta ^{}`$ leads to new operators in the strong and weak effective lagrangians. Under CPS symmetry all operators in $`_{2,4}`$ of Eq. (2.1) are even, and, since $`\eta ^{}`$ is CPS odd, new weak operators can be constructed by multiplying $`_{2,4}`$ by even powers of $`\eta ^{}`$. It turns out that such operators do not contribute to the quantities of interest in this paper. However, other new weak operators arise from multiplying CPS-odd weak operators by odd powers of $`\eta ^{}`$. There are two operators of interest at order $`p^2`$, so that $`_2^\eta ^{}`$ is given by $$_2^\eta ^{}=\gamma _1^8_\mu (\eta ^{}/f)\text{str}(\mathrm{\Lambda }L_\mu )+i\gamma _2^8(\eta ^{}/f)\text{str}(\mathrm{\Lambda }X_{}).$$ (2.13) Since we are not interested in processes with external $`\eta ^{}`$ lines, we do not consider new operators of order $`p^4`$ containing the $`\eta ^{}`$ field. For $`N=3`$ sea quarks the dynamics of the partially quenched theory is precisely that of unquenched QCD with degenerate quark masses . Since all the low-energy constants (LECs) $`\alpha _i^8,\alpha ^{27}`$, $`\gamma _i^8`$ and $`\beta _i^{8,27},\stackrel{~}{\beta }_i^8`$ (together with the strong counterterms) are independent of quark masses, it follows that their $`N=3`$ partially-quenched values are equal to those of the real world . However, for this equivalence to be valid, the $`\eta ^{}`$ should be treated in the same way in the partially quenched theory as in the real world, so that one has to consider the limit in which the $`\eta ^{}`$ decouples. ## 3 The role of the $`\eta ^{}`$ in partially quenched ChPT Before we present our results, we would like to discuss the role of the $`\eta ^{}`$ in a partially quenched theory in more detail. First, define the bare (or tree-level) meson masses $$M_{ij}^2=B_0(m_i+m_j),i,j=1,\mathrm{}2nN,$$ (3.1) for a light pseudoscalar (pseudo-Goldstone) meson made out of quarks or ghost quarks $`i`$ and $`j`$. For degenerate sea quarks, this simplifies to $`M_{SS}^2=2B_0m_S`$. The two-point function for neutral mesons $`\mathrm{\Phi }_{ii}`$ (in that basis) is given by $`\mathrm{\Phi }_{ii}(x)\mathrm{\Phi }_{jj}(0)`$ $`=`$ $`{\displaystyle \frac{d^4p}{(2\pi )^4}e^{ipx}G_{i,j}(p)},`$ $`G_{i,j}(p)`$ $`=`$ $`{\displaystyle \frac{\delta _{ij}ϵ_i}{p^2+M_{ii}^2}}X_{ij}(p),`$ (3.2) $`X_{ij}(p)`$ $`=`$ $`{\displaystyle \frac{1}{3+N\alpha }}{\displaystyle \frac{(m_0^2+\alpha p^2)(p^2+M_{SS}^2)}{(p^2+M_{ii}^2)(p^2+M_{jj}^2)(p^2+M_\eta ^{}^2)}},`$ where $$ϵ_i=\{\begin{array}{cc}+1,\hfill & \text{for }1in\hfill \\ 1,\hfill & \text{for }n+1i2nN\hfill \end{array}$$ and the $`\eta ^{}`$ mass is given by $$M_\eta ^{}^2=\frac{M_{SS}^2+Nm_0^2/3}{1+N\alpha /3}.$$ (3.3) The parameters $`m_0^2`$ and $`\alpha `$ (not to be confused with $`\alpha _i^{8,27}`$) come from the strong-lagrangian $`O(p^2)`$ operators quadratic in the $`\eta ^{}`$ field, $`(Nm_0^2/6)(\eta ^{})^2+(N\alpha /6)(_\mu \eta ^{})^2`$. The term $`X_{ij}(p)`$ has a double pole for $`M_{ii}=M_{jj}`$, unless $`M_{SS}=M_{ii}=M_{jj}`$. This implies that partially quenched theories suffer from the same “quenched infrared diseases” as the quenched theory unless all valence-quark masses are equal to the sea-quark mass . From Eq. (3.2) it is easily verified that the $`\eta ^{}`$ two-point function is just $`(N/(1+N\alpha /3))(p^2+M_\eta ^{}^2)^1`$. The quantity $`X_{ij}(p)`$ can also be written as $`X_{ij}(p)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{A_{ij}}{p^2+M_{ii}^2}}{\displaystyle \frac{B_{ij}}{p^2+M_\eta ^{}^2}}+{\displaystyle \frac{_{ij}^2A_{ij}M_{jj}^2}{(p^2+M_{ii}^2)(p^2+M_{jj}^2)}}\right),`$ (3.4) with $`_{ij}^2`$ $`=`$ $`{\displaystyle \frac{(N/3)(m_0^2\alpha M_{ii}^2)(m_0^2\alpha M_{jj}^2)M_{SS}^2+m_0^2(M_{SS}^2M_{ii}^2)(M_{SS}^2M_{jj}^2)}{[(N/3)(m_0^2\alpha M_{ii}^2)+M_{SS}^2M_{ii}^2][(N/3)(m_0^2\alpha M_{jj}^2)+M_{SS}^2M_{jj}^2]}},`$ $`A_{ij}`$ $`=`$ $`{\displaystyle \frac{(N/3)(m_0^2\alpha M_{ii}^2)(m_0^2\alpha M_{jj}^2)+\alpha (M_{SS}^2M_{ii}^2)(M_{SS}^2M_{jj}^2)}{[(N/3)(m_0^2\alpha M_{ii}^2)+M_{SS}^2M_{ii}^2][(N/3)(m_0^2\alpha M_{jj}^2)+M_{SS}^2M_{jj}^2]}},`$ $`B_{ij}`$ $`=`$ $`{\displaystyle \frac{(N/3)(m_0^2\alpha M_{SS}^2)^2/(1+\alpha N/3)}{[(N/3)(m_0^2\alpha M_{ii}^2)+M_{SS}^2M_{ii}^2][(N/3)(m_0^2\alpha M_{jj}^2)+M_{SS}^2M_{jj}^2]}}.`$ The coefficients $`A_{ij}`$, $`B_{ij}`$ and $`_{ij}^2`$ are complicated functions of the various mass scales in the partially quenched effective theory. We may consider various limits in which these expressions simplify considerably. First, one easily obtains the fully quenched expression by setting $`N=0`$, or equivalently taking $`M_{SS}\mathrm{}`$, finding for all $`ij`$ $$_{ij}^2m_0^2,A_{ij}\alpha ,B_{ij}0.$$ (3.6) It is clear from these expressions that in the quenched case, the $`\eta ^{}`$ should be kept in the effective theory. Another interesting limit is that in which the $`\eta ^{}`$ decouples (which, as inspection of Eq. (3.3) tells us, is only possible for $`N>0`$). In this limit, again for all $`ij`$, $$_{ij}^2\frac{3}{N}M_{SS}^2,A_{ij}\frac{3}{N},$$ (3.7) and we drop the $`\eta ^{}`$ pole in Eq. (3.4). The dependence of $`X_{ij}`$ on the $`\eta ^{}`$ parameters has disappeared in this limit. As argued in Ref. , in actual partially quenched Lattice QCD computations, the sea-meson mass $`M_{SS}`$ maybe comparable in size to $`m_0`$ so that the full dependence of $`_{ij}^2`$ and $`A_{ij}`$ on the parameters $`m_0`$, $`\alpha `$ and $`M_{SS}`$ should be kept. A third possibility is then given by the limit in which the valence-meson mass is small compared to the $`\eta ^{}`$ mass, i.e. $`M_{kk}M_\eta ^{}`$. The expressions for $`_{ij}^2`$, $`A_{ij}`$ and $`B_{ij}`$ given in Eq. (LABEL:MAB) reduce to those of Ref. if we expand in $`M_{kk}^2/M_\eta ^{}^2`$ but not in $`M_{SS}^2/M_\eta ^{}^2`$: $`_{ij}^2`$ $``$ $`^2={\displaystyle \frac{m_0^2M_{SS}^2}{(N/3)m_0^2+M_{SS}^2}}\left(1+O\left({\displaystyle \frac{M_{kk}^4}{M_\eta ^{}^4}}\right)\right),`$ (3.8) $`A_{ij}`$ $``$ $`A={\displaystyle \frac{(N/3)m_0^4+\alpha M_{SS}^4}{[(N/3)m_0^2+M_{SS}^2]^2}}+O\left({\displaystyle \frac{M_{kk}^2}{M_\eta ^{}^2}}\right),`$ $`B_{ij}`$ $``$ $`{\displaystyle \frac{(N/3)(m_0^2\alpha M_{SS}^2)^2/(1+\alpha N/3)}{[(N/3)m_0^2+M_{SS}^2]^2}}+O\left({\displaystyle \frac{M_{kk}^2}{M_\eta ^{}^2}}\right).`$ The (partially) quenched expansion we consider in this paper is systematic if we take $`^2`$ to be of order $`p^2`$, in other words, if we take the parameter $`m_0^2`$ to be of the same order as the quark mass, just as in the case of quenched ChPT . It was also shown in Ref. that, for the quantities considered there, simply ignoring one-loop contributions coming from the $`\eta ^{}`$ pole in Eq. (3.4) and then taking the limit $`M_\eta ^{}\mathrm{}`$ in the rest is the same as matching to the limit in which the $`\eta ^{}`$ decouples. In other words, if we ignore these contributions, the LECs appearing in those quantities are the same in the $`N=3`$ partially quenched world and the real world. In addition, when we take only $`M_\eta ^{}/M_{kk}`$ large, but not $`M_\eta ^{}/M_{SS}`$, these one-loop contributions are polynomial in the valence-meson masses, and can still be ignored if we are interested in the non-analytic dependence on the valence-meson masses. The same observations are also true here, even though there are diagrams with a more complicated topology than the simple tadpole diagrams needed in Ref. . For $`K\pi `$, there are contributions of the form depicted in Fig. 1. Taking $`M=M_{kk}`$ to be the degenerate valence-meson mass running in the loop and abbreviating $`B=B_{kk}`$, this diagram leads to one-loop integrals such as $$\frac{d^4p}{(2\pi )^4}\frac{1}{p^2+M^2}\frac{B}{p^2+M_\eta ^{}^2}=\frac{B}{M_\eta ^{}^2M^2}\frac{d^4p}{(2\pi )^4}\left(\frac{1}{p^2+M^2}\frac{1}{p^2+M_\eta ^{}^2}\right),$$ (3.9) multiplied by two powers of $`M^2`$ from the two $`O(p^2)`$ vertices in the diagram (there are no contributions from vertices proportional to $`M_{SS}^2`$, so the only dependence on $`M_{SS}`$ comes through the coefficient $`B`$). The integral contributes an $`\eta ^{}`$ chiral logarithm, which we drop as above, and a Goldstone chiral logarithm proportional to $`BM^6/(M_\eta ^{}^2M^2)\mathrm{log}M^2`$. If we do not assume that $`M`$ is small compared to $`M_\eta ^{}`$ this is of order $`M^4`$, but if we expand in $`M^2/M_\eta ^{}^2`$, this constitutes an $`O(p^6)`$ contribution. In all the following calculations, we will assume that valence-meson masses are sufficiently small compared to $`M_\eta ^{}`$ to justify the expansion in $`M^2/M_\eta ^{}^2`$. If one would not make this assumption, all contributions from integrals like Eq. (3.9) would have to be kept, since both the Goldstone-meson and the $`\eta ^{}`$ pole give rise to additional non-polynomial $`M`$ dependence. However, with this assumption, contributions coming from $`\eta ^{}`$ tadpoles or diagrams such as Fig. 1 containing an $`\eta ^{}`$ on the loop are analytic in the valence-meson masses to order $`p^4`$, and we will therefore drop them from consideration. Finally, we note that the coefficients of chiral logarithms will in general depend in a complicated non-polynomial way on the valence- and sea-meson masses, so that the $`O(p^4)`$ LECs cannot be defined in a mass-independent way. The coefficients of chiral logarithms have a polynomial dependence on meson masses only if we expand in both $`M^2/M_\eta ^{}^2`$ and $`M_{SS}^2/M_\eta ^{}^2`$. In that case the LECs of the partially quenched theory without the $`\eta ^{}`$ (and with $`N=3`$) are the same as those of the real world. ## 4 $`K\pi `$ and $`K0`$ matrix elements at $`O(p^4)`$ In this section we calculate $`K^00`$ and $`K^+\pi ^+`$ matrix elements at order $`p^4`$ in the effective theory. Four cases are considered: partially quenched with the valence-meson mass $`MM_\eta ^{}`$ but $`M_{SS}/M_\eta ^{}`$ arbitrary, partially quenched with also $`M_\eta ^{}/M_{SS}`$ large, quenched, and unquenched. ### 4.1 Partially quenched results We consider a partially quenched theory with three valence quarks, $`u`$, $`d`$ and $`s`$ and $`N`$ degenerate sea quarks. The matrix elements to be calculated are defined as $`[K^00]`$ $``$ $`0\|_{\mathrm{\Delta }S=1}\|K^0,`$ (4.1) $`[K^+\pi ^+]`$ $``$ $`\pi ^+\|_{\mathrm{\Delta }S=1}\|K^+.`$ The first process is calculated for $`m_sm_d`$ and $`m_u=m_d`$. In the last process we take the valence-quark masses all equal $`m_s=m_d=m_u`$, so that it conserves energy. In this SU(3) limit, the $`K^0\eta `$ matrix element does not contain any extra information. We have calculated these matrix elements to $`O(p^4)`$ in partially quenched ChPT, using dimensional regularization in the $`\overline{MS}`$ scheme, and including contributions from the $`O(p^4)`$ operators (2.9) and (2.10). In the case with degenerate valence quarks, let $`M=M_{ii}`$ be the physical mass of a meson made out of valence quarks, $`M_{SS}`$ that of a meson made out of sea quarks, and $`M_{VS}`$ that of a meson made out of a valence and a sea quark. At tree level in ChPT one has $$M_{VS}^{}{}_{}{}^{2}=\frac{1}{2}(M^2+M_{SS}^2).$$ (4.2) Define, for any mass M, $$L(M)=\mathrm{log}\frac{M^2}{\mathrm{\Lambda }^2},$$ (4.3) where $`\mathrm{\Lambda }`$ is the $`\overline{MS}`$ scale. In addition, define, for any two masses $`M_1`$ and $`M_2`$, $$L_n(M_1,M_2)=\frac{M_1^{2n}\mathrm{log}\frac{M_1^2}{\mathrm{\Lambda }^2}M_2^{2n}\mathrm{log}\frac{M_2^2}{\mathrm{\Lambda }^2}}{M_1^2M_2^2}.$$ (4.4) As discussed in the previous section, we will assume that the valence-meson masses are small compared to $`M_\eta ^{}`$ (cf. Eq. (3.3)), but not make the same assumption about $`M_{SS}`$. At one loop, we then have, for the octet $`K\pi `$ matrix element, using $`^2`$ and $`A`$ from Eq. (3.8), $`[K^+\pi ^+]_8={\displaystyle \frac{4\alpha _1^8M^2}{f^2}}(1{\displaystyle \frac{1}{(4\pi f)^2}}[NM_{VS}^{}{}_{}{}^{2}(L(M_{VS})1)`$ (4.5) $`+\mathrm{\hspace{0.17em}2}(^2{\displaystyle \frac{8}{3}}AM^2)L(M)+{\displaystyle \frac{2}{3}}(^2+2AM^2)])`$ $`{\displaystyle \frac{4\alpha _2^8M^2}{f^2}}\left(1{\displaystyle \frac{1}{(4\pi f)^2}}\left[{\displaystyle \frac{2}{3}}(^24AM^2)L(M)+{\displaystyle \frac{2}{3}}^2\right]\right)`$ $`+{\displaystyle \frac{4(\gamma _1^8+2\gamma _2^8)M^2}{f^2(4\pi f)^2}}[M^2(2L(M)1){\displaystyle \frac{1}{3}}NA(L_2(M_{SS},M)M_{SS}^2M^2)`$ $`{\displaystyle \frac{N}{3}}{\displaystyle \frac{^2AM^2}{M_{SS}^2M^2}}(2M^2L(M)L_2(M_{SS},M)+M_{SS}^2)],`$ and for the 27-plet $`K\pi `$ matrix element, $`[K^+\pi ^+]_{27}={\displaystyle \frac{4\alpha ^{27}M^2}{f^2}}(1{\displaystyle \frac{1}{(4\pi f)^2}}[2NM_{VS}^{}{}_{}{}^{2}(L(M_{VS})1)+2M^2(3L(M)2)`$ $`+{\displaystyle \frac{2}{3}}(^22AM^2)L(M)+{\displaystyle \frac{2}{3}}AM^2]).`$ (4.6) In the latter case, the matrix elements for $`\mathrm{\Delta }I=1/2`$ and $`\mathrm{\Delta }I=3/2`$ are the same, because of SU(3) symmetry. The $`\eta ^{}`$ does not contribute directly to the 27-plet matrix elements; the dependence on $`m_0^2`$ and $`\alpha `$ comes from the fact that we expressed all results in terms of bare meson masses, or, equivalently, in terms of quark masses (cf. Eq. (3.1)). The octet matrix element receives instead direct contributions from $`\eta ^{}`$ exchange. Note that the contributions from the two $`O(p^2)`$ $`\eta ^{}`$-operators of Eq. (2.13) have the same form. This is explained by the fact that, after a partial integration, the first term in Eq. (2.13), using the equation of motion for $`\mathrm{\Sigma }`$, is proportional to the second term. For the contributions from the $`O(p^4)`$ operators, we find $`(4\pi f)^2[K^+\pi ^+]_8^{(4)}`$ $`=`$ $`{\displaystyle \frac{8M^2}{f^2}}[(2\beta _1^8+2\beta _3^8+2\beta _{10}^8+\beta _{11}^8)M^2+\beta _2^8NM_{SS}^2`$ $`+16\alpha _1^8((\lambda _4\lambda _6)NM_{SS}^2+(\lambda _5\lambda _8)M^2)8\alpha _2^8(\lambda _4NM_{SS}^2+\lambda _5M^2)],`$ $`(4\pi f)^2[K^+\pi ^+]_{27}^{(4)}`$ $`=`$ $`{\displaystyle \frac{8M^2}{f^2}}[2(\beta _2^{27}+\beta _4^{27})M^2+\beta _7^{27}NM_{SS}^2`$ $`16\alpha ^{27}((\lambda _4\lambda _6)NM_{SS}^2+(\lambda _5\lambda _8)M^2)],`$ with again the $`\mathrm{\Delta }I=1/2`$ and $`\mathrm{\Delta }I=3/2`$ results the same for the 27-plet. The $`\lambda _i`$ are the strong $`O(p^4)`$ LECs; they are related to the Gasser–Leutwyler $`L_i`$ by $$\lambda _i=16\pi ^2L_i.$$ (4.9) We see an example here of the fact that a partially quenched simulation, with $`MM_{SS}`$, would in principle yield more information about the $`O(p^4)`$ LECs than an unquenched simulation in which $`M=M_{SS}`$. For the $`K0`$ matrix elements we take non-degenerate valence quarks with $`m_sm_d`$ and $`m_d=m_u`$, and define $$M_{33}^2=2M_K^2M_\pi ^2.$$ (4.10) The pion will be made out of two light valence quarks, and the kaon out of a light and a strange valence quark. We also define $`M_{iS}^2`$ to be the (tree-level) mass of a meson made out of the $`i`$-th valence quark and a sea quark, $$M_{iS}^2=B_0(m_i+m_S),i=u,d,s.$$ (4.11) Of course, $`M_{uS}=M_{dS}`$. For the octet matrix element we find, at one loop, $`[K0]_8={\displaystyle \frac{4i\alpha _1^8}{f(4\pi f)^2}}[N(M_{uS}^4(L(M_{uS})1)M_{sS}^4(L(M_{sS})1))`$ (4.12) $`+{\displaystyle \frac{2}{3}}^2\left(M_\pi ^2(L(M_\pi ){\displaystyle \frac{1}{2}})M_{33}^2(L(M_{33}){\displaystyle \frac{1}{2}})\right)`$ $`A(M_\pi ^4(L(M_\pi ){\displaystyle \frac{2}{3}})M_{33}^4(L(M_{33}){\displaystyle \frac{2}{3}}))]`$ $`+{\displaystyle \frac{4i\alpha _2^8(M_K^2M_\pi ^2)}{f}}(1+{\displaystyle \frac{1}{(4\pi f)^2}}[{\displaystyle \frac{N}{2}}(M_{uS}^2(L(M_{uS})1)+M_{sS}^2(L(M_{sS})1))`$ $`{\displaystyle \frac{1}{6}}^2\left(L(M_\pi )+L(M_{33})+2L_1(M_{33},M_\pi )2\right)`$ $`+{\displaystyle \frac{1}{6}}A(2M_\pi ^2L(M_\pi )+2M_{33}^2L(M_{33})+2L_2(M_{33},M_\pi )3M_\pi ^23M_{33}^2)])`$ $`+{\displaystyle \frac{2i\gamma _1^8}{f(4\pi f)^2}}[M_\pi ^4(L(M_\pi )1)M_{33}^4(L(M_{33})1)`$ $`+{\displaystyle \frac{N}{3}}(^2AM_\pi ^2)\left(L_2(M_{SS},M_\pi )M_{SS}^2M_\pi ^2\right)`$ $`{\displaystyle \frac{N}{3}}(^2AM_{33}^2)(L_2(M_{SS},M_{33})M_{SS}^2M_{33}^2)]`$ $`{\displaystyle \frac{4i\gamma _2^8(M_K^2M_\pi ^2)}{f(4\pi f)^2}}[M_\pi ^2(L(M_\pi )1)+M_{33}^2(L(M_{33})1){\displaystyle \frac{2N}{3}}AM_{SS}^2(L(M_{SS})1)`$ $`+{\displaystyle \frac{N}{3}}(^2AM_\pi ^2)(L_1(M_{SS},M_\pi )1)+{\displaystyle \frac{N}{3}}(^2AM_{33}^2)(L_1(M_{SS},M_{33})1)].`$ In this case, the leading order $`p^2`$ contribution comes only from the weak mass term proportional to $`\alpha _2^8`$. Unlike the case of $`K\pi `$, the contributions from the $`O(p^2)`$ $`\eta ^{}`$operators proportional to $`\gamma _1^8`$ and $`\gamma _2^8`$ are different in this case. The argument explaining the situation in the $`K\pi `$ case does not work here, because the total derivative in the partial integration cannot be dropped, as the process $`K0`$ does not conserve energy. For the 27-plet (both $`\mathrm{\Delta }I=1/2`$ and $`\mathrm{\Delta }I=3/2`$), we obtain $`[K0]_{27}={\displaystyle \frac{12i\alpha ^{27}}{f(4\pi f)^2}}[M_\pi ^4(L(M_\pi )1)+M_{33}^4(L(M_{33})1)2M_K^4(L(M_K)1)`$ $`+{\displaystyle \frac{2}{3}}^2\left(M_\pi ^2L(M_\pi )+M_{33}^2L(M_{33})L_2(M_{33},M_\pi )+{\displaystyle \frac{1}{2}}(M_\pi ^2+M_{33}^2)\right)`$ $`{\displaystyle \frac{1}{3}}A(L_3(M_{33},M_\pi )3M_\pi ^2M_{33}^2L_1(M_{33},M_\pi )+2M_\pi ^2M_{33}^2)].`$ (4.13) Finally, the $`O(p^4)`$ operators of Eqs. (2.9,2.10,2.11) give $`(4\pi f)^2[K0]_8^{(4)}`$ $`=`$ $`{\displaystyle \frac{8i(M_K^2M_\pi ^2)}{f}}[(2\beta _1^82\beta _5^8+\stackrel{~}{\beta }_1^8)M_K^2+\beta _2^8NM_{SS}^2`$ (4.14) $`4\alpha _2^8(\lambda _4NM_{SS}^2+\lambda _5M_K^2)],`$ $`(4\pi f)^2[K0]_{27}^{(4)}`$ $`=`$ $`{\displaystyle \frac{12i}{f}}\mathrm{\hspace{0.17em}4}\beta _1^{27}(M_K^2M_\pi ^2)^2.`$ (4.15) These results hold for an arbitrary number $`N`$ of degenerate sea quarks. Note the appearance of “quenched chiral logs,” contained in the one-loop logarithms proportional to $`^2`$. Since $`^2`$ does not depend on the valence masses, such terms decrease with decreasing valence quark masses at the same rate as tree-level terms (modulo the logarithms), i.e. typically as $`m\mathrm{log}m`$ instead of $`m^2\mathrm{log}m`$. Our results are presented here in a form somewhat different from that in Ref. . Here, we express the matrix elements in terms of the tree-level meson masses, or equivalently, the quark masses, while in Ref. they were expressed in terms of renormalized masses (also, only the chiral logarithms were given). These results can be converted into expressions for the matrix elements as a function of the actual meson masses computed on the lattice by using the one-loop expression for the mass of a meson made out of two non-degenerate valence quarks in terms of the tree-level masses, Eq. (3.1). This expression, in $`\overline{MS}`$, and including $`O(p^4)`$ contributions, is $`\left(M_K^2\right)^{1\mathrm{loop}}`$ $`=`$ $`M_K^2(1+{\displaystyle \frac{1}{(4\pi f)^2}}[{\displaystyle \frac{2}{3}}^2(L_1(M_{33},M_\pi )1)`$ $`+{\displaystyle \frac{2}{3}}A\left(L_2(M_{33},M_\pi )M_{33}^2M_\pi ^2\right)`$ $`+16(M_K^2(2\lambda _8\lambda _5)+NM_{SS}^2(2\lambda _6\lambda _4))]).`$ For degenerate valence-quark masses, this simplifies to $`\left(M_\pi ^2\right)^{1\mathrm{loop}}`$ $`=`$ $`M^2(1+{\displaystyle \frac{1}{(4\pi f)^2}}[{\displaystyle \frac{2}{3}}^2L(M)+{\displaystyle \frac{2}{3}}AM^2(2L(M)1)`$ $`+16(M^2(2\lambda _8\lambda _5)+NM_{SS}^2(2\lambda _6\lambda _4))]).`$ ### 4.2 Partially quenched results for large $`M_\eta ^{}`$ The results presented above simplify when $`M_\eta ^{}`$ is taken large compared to both the sea- and valence-meson masses. Taking $`M_\eta ^{}`$ large while keeping $`M_{ii,jj}`$ and $`M_{SS}`$ fixed in Eq. (3.4) gives $$X_{ij}(p)=\frac{1}{N}\left(\frac{1}{p^2+M_{ii}^2}+\frac{M_{SS}^2M_{jj}^2}{(p^2+M_{ii}^2)(p^2+M_{jj}^2)}\right),$$ (4.18) dropping the $`\eta ^{}`$ pole. This expression for $`X_{ij}(p)`$ leads to the simplified expressions in Eq. (3.7) for $`^2`$ and $`A`$, where all the $`\eta ^{}`$ parameters have disappeared, and the $`\eta ^{}`$ meson has been decoupled. Thus, for $`K\pi `$, we obtain $`[K^+\pi ^+]_8`$ $`=`$ $`{\displaystyle \frac{4\alpha _1^8M^2}{f^2}}(1{\displaystyle \frac{1}{(4\pi f)^2}}[NM_{VS}^{}{}_{}{}^{2}(L(M_{VS})1)`$ $`+{\displaystyle \frac{2}{N}}((3M_{SS}^28M^2)L(M)+M_{SS}^2+2M^2)])`$ $`{\displaystyle \frac{4\alpha _2^8M^2}{f^2}}\left(1{\displaystyle \frac{1}{(4\pi f)^2}}{\displaystyle \frac{2}{N}}\left[(M_{SS}^24M^2)L(M)+M_{SS}^2\right]\right),`$ and $`[K^+\pi ^+]_{27}`$ $`=`$ $`{\displaystyle \frac{4\alpha ^{27}M^2}{f^2}}(1{\displaystyle \frac{1}{(4\pi f)^2}}[2NM_{VS}^{}{}_{}{}^{2}(L(M_{VS})1)`$ $`+\mathrm{\hspace{0.17em}2}M^2(3L(M)2)+{\displaystyle \frac{2}{N}}((M_{SS}^22M^2)L(M)+M^2)]),`$ while for $`K0`$, we find $`[K0]_8={\displaystyle \frac{4i\alpha _1^8}{f(4\pi f)^2}}[N(M_{uS}^4(L(M_{uS})1)M_{sS}^4(L(M_{sS})1)`$ (4.21) $`+{\displaystyle \frac{2}{N}}M_{SS}^2\left(M_\pi ^2(L(M_\pi ){\displaystyle \frac{1}{2}})M_{33}^2(L(M_{33}){\displaystyle \frac{1}{2}})\right)`$ $`{\displaystyle \frac{3}{N}}(M_\pi ^4(L(M_\pi ){\displaystyle \frac{2}{3}})M_{33}^4(L(M_{33}){\displaystyle \frac{2}{3}}))]`$ $`+{\displaystyle \frac{4i\alpha _2^8(M_K^2M_\pi ^2)}{f}}(1+{\displaystyle \frac{1}{(4\pi f)^2}}[{\displaystyle \frac{N}{2}}(M_{uS}^2(L(M_{uS})1)+M_{sS}^2(L(M_{sS})1))`$ $`{\displaystyle \frac{1}{2N}}M_{SS}^2\left(L(M_\pi )+L(M_{33})+2L_1(M_{33},M_\pi )2\right)`$ $`+{\displaystyle \frac{1}{2N}}(2M_\pi ^2L(M_\pi )+2M_{33}^2L(M_{33})+2L_2(M_{33},M_\pi )3M_\pi ^23M_{33}^2)])`$ for the octet, and $`[K0]_{27}={\displaystyle \frac{12i\alpha ^{27}}{f(4\pi f)^2}}[M_\pi ^4(L(M_\pi )1)+M_{33}^4(L(M_{33})1)2M_K^4(L(M_K)1)`$ $`+{\displaystyle \frac{2}{N}}M_{SS}^2\left(M_\pi ^2L(M_\pi )+M_{33}^2L(M_{33})L_2(M_{33},M_\pi )+{\displaystyle \frac{1}{2}}(M_\pi ^2+M_{33}^2)\right)`$ $`{\displaystyle \frac{1}{N}}(L_3(M_{33},M_\pi )3M_\pi ^2M_{33}^2L_1(M_{33},M_\pi )+2M_\pi ^2M_{33}^2)]`$ (4.22) for the 27-plet. All the dependence on $`\eta ^{}`$ parameters ($`m_0^2`$, $`\alpha `$, $`\gamma _{1,2}^8`$) has disappeared from these expressions, as expected. The contributions from $`O(p^4)`$ operators, given in Eqs. (4.1,4.1) and (4.14,4.15), do not change. Note however, as mentioned before, that only in this limit the corresponding LECs are independent of the quark masses. ### 4.3 Quenched results A special case of practical interest is the completely quenched result. In the quenched approximation, there are no sea quarks, and hence, quenched expressions can be obtained by setting $`N=0`$ in Eqs. (4.54.1,4.12,4.13,4.14,4.15), or equivalently, by taking $`M_{SS}\mathrm{}`$. In this case, it is not possible to decouple the $`\eta ^{}`$ . The expressions given below maybe rewritten in terms of the parameter $`\delta `$ (introduced in Ref. ), by setting $$m_0^2=24\pi ^2f^2\delta .$$ (4.23) The quenched results are $`[K^+\pi ^+]_8={\displaystyle \frac{4\alpha _1^8M^2}{f^2}}\left(1{\displaystyle \frac{1}{(4\pi f)^2}}\left[2\left(m_0^2{\displaystyle \frac{8}{3}}\alpha M^2\right)L(M)+{\displaystyle \frac{2}{3}}(m_0^2+2\alpha M^2)\right]\right)`$ $`{\displaystyle \frac{4\alpha _2^8M^2}{f^2}}\left(1{\displaystyle \frac{1}{(4\pi f)^2}}\left[{\displaystyle \frac{2}{3}}(m_0^24\alpha M^2)L(M)+{\displaystyle \frac{2}{3}}m_0^2\right]\right)`$ (4.24) $`+{\displaystyle \frac{4(\gamma _1^8+2\gamma _2^8)M^2}{f^2(4\pi f)^2}}\left[2M^2L(M)M^2\right],`$ $`[K^+\pi ^+]_{27}={\displaystyle \frac{4\alpha ^{27}M^2}{f^2}}(1{\displaystyle \frac{1}{(4\pi f)^2}}[2M^2(3L(M)2)`$ (4.25) $`+{\displaystyle \frac{2}{3}}(m_0^22\alpha M^2)L(M)+{\displaystyle \frac{2}{3}}\alpha M^2]),`$ $`[K0]_8={\displaystyle \frac{4i\alpha _1^8}{f(4\pi f)^2}}[{\displaystyle \frac{2}{3}}m_0^2(M_\pi ^2(L(M_\pi ){\displaystyle \frac{1}{2}})M_{33}^2(L(M_{33}){\displaystyle \frac{1}{2}}))`$ (4.26) $`\alpha (M_\pi ^4(L(M_\pi ){\displaystyle \frac{2}{3}})M_{33}^4(L(M_{33}){\displaystyle \frac{2}{3}}))]`$ $`+{\displaystyle \frac{4i\alpha _2^8(M_K^2M_\pi ^2)}{f}}(1+{\displaystyle \frac{1}{(4\pi f)^2}}[{\displaystyle \frac{1}{6}}m_0^2(L(M_\pi )+L(M_{33})+2L_1(M_{33},M_\pi )2)`$ $`+{\displaystyle \frac{1}{6}}\alpha (2M_\pi ^2L(M_\pi )+2M_{33}^2L(M_{33})+2L_2(M_{33},M_\pi )3M_\pi ^23M_{33}^2)])`$ $`+{\displaystyle \frac{2i\gamma _1^8}{f(4\pi f)^2}}\left[M_\pi ^4(L(M_\pi )1)M_{33}^4(L(M_{33})1)\right]`$ $`{\displaystyle \frac{4i\gamma _2^8(M_K^2M_\pi ^2)}{f(4\pi f)^2}}\left[M_\pi ^2(L(M_\pi )1)+M_{33}^2(L(M_{33})1)\right],`$ $`[K0]_{27}={\displaystyle \frac{12i\alpha ^{27}}{f(4\pi f)^2}}[M_\pi ^4(L(M_\pi )1)+M_{33}^4(L(M_{33})1)2M_K^4(L(M_K)1)`$ $`+{\displaystyle \frac{2}{3}}m_0^2\left(M_\pi ^2L(M_\pi )+M_{33}^2L(M_{33})L_2(M_{33},M_\pi )+{\displaystyle \frac{1}{2}}(M_\pi ^2+M_{33}^2)\right)`$ $`{\displaystyle \frac{1}{3}}\alpha (L_3(M_{33},M_\pi )+3M_\pi ^2M_{33}^2L_1(M_{33},M_\pi )+2M_\pi ^2M_{33}^2)].`$ (4.27) Contributions from $`O(p^4)`$ operators are obtained by setting $`N=0`$ in Eqs. (4.1,4.1) and (4.14,4.15). However, we emphasize that all the information obtained from quenched lattice computations is about the $`N=0`$ values of the LECs appearing in these equations. These values are, in principle, different from their $`N=3`$ values. ### 4.4 Unquenched results For completeness, we also report the results for the unquenched theory with three light flavors. For $`K\pi `$ these can be simply obtained by setting $`N=3`$ and $`M_{SS}^2=M^2`$ in our (large-$`M_\eta ^{}`$) partially quenched expressions of Subsec. 4.2. For $`K0`$ we have to choose the sea-quark masses equal to the non-degenerate valence-quark masses. The results for $`K0`$ therefore cannot be derived from our partially quenched results, where we took all sea quarks to be degenerate in mass from the start. The results are $`[K^+\pi ^+]_8`$ $`=`$ $`{\displaystyle \frac{4\alpha _1^8M^2}{f^2}}\left(1+{\displaystyle \frac{1}{(4\pi f)^2}}\left[{\displaystyle \frac{1}{3}}M^2L(M)+M^2\right]\right)`$ $`{\displaystyle \frac{4\alpha _2^8M^2}{f^2}}\left(1+{\displaystyle \frac{1}{(4\pi f)^2}}\left[2M^2L(M){\displaystyle \frac{2}{3}}M^2\right]\right),`$ $`[K^+\pi ^+]_{27}`$ $`=`$ $`{\displaystyle \frac{4\alpha ^{27}M^2}{f^2}}\left(1{\displaystyle \frac{1}{(4\pi f)^2}}\left[{\displaystyle \frac{34}{3}}M^2L(M){\displaystyle \frac{28}{3}}M^2\right]\right)`$ (4.29) for $`K\pi `$, and $`[K0]_8`$ $`=`$ $`{\displaystyle \frac{4i\alpha _2^8}{f}}(M_K^2M_\pi ^2)(1+{\displaystyle \frac{1}{(4\pi f)^2}}[{\displaystyle \frac{3}{4}}M_\pi ^2L(M_\pi ){\displaystyle \frac{3}{2}}M_K^2L(M_K)`$ $`{\displaystyle \frac{1}{12}}M_\eta ^2L(M_\eta )+{\displaystyle \frac{29}{18}}M_K^2+{\displaystyle \frac{13}{18}}M_\pi ^2])`$ $`+{\displaystyle \frac{4i\alpha _1^8}{f}}{\displaystyle \frac{M_K^2M_\pi ^2}{(4\pi f)^2}}\left[{\displaystyle \frac{1}{3}}L_2(M_\eta ,M_K)2L_2(M_\eta ,M_\pi )+{\displaystyle \frac{17}{9}}M_K^2+{\displaystyle \frac{13}{9}}M_\pi ^2\right],`$ $`[K0]_{27}`$ $`=`$ $`{\displaystyle \frac{4i\alpha ^{27}}{f}}{\displaystyle \frac{M_K^2M_\pi ^2}{(4\pi f)^2}}\left[2L_2(M_\eta ,M_K)+2L_2(M_\eta ,M_\pi )+2M_K^22M_\pi ^2\right]`$ for $`K0`$. The contributions from $`O(p^4)`$ operators are obtained from the partially quenched expressions Eqs. (4.1,4.1,4.14,4.15), by replacing $`NM_{SS}^22M_K^2+M_\pi ^2`$ in those equations. ## 5 Relation to $`K\pi \pi `$ and numerical examples We now turn to a discussion on how our results can be used to extract physical information from lattice results. If tree-level ChPT were a good approximation, one could determine $`\alpha _{1,2}^8`$ and $`\alpha ^{27}`$ from a lattice computation of the $`K\pi `$ and $`K0`$ matrix elements, and then use ChPT to predict the $`\mathrm{\Delta }I=1/2`$ and $`\mathrm{\Delta }I=3/2`$ $`K\pi \pi `$ decay rates . For instance, for the $`\mathrm{\Delta }I=1/2`$ matrix element, one finds $`[K^0\pi ^+\pi ^{}]_{1/2}`$ $`=`$ $`{\displaystyle \frac{4i(m_K^2m_\pi ^2)}{f^3}}(\alpha _1^8\alpha ^{27})`$ $`=`$ $`{\displaystyle \frac{i}{f}}{\displaystyle \frac{m_K^2m_\pi ^2}{M^2}}\left([K^+\pi ^+]_{1/2}b[K^00]\right),`$ $`b`$ $``$ $`{\displaystyle \frac{iM^2}{f(M_K^2M_\pi ^2)}}={\displaystyle \frac{2im}{f(m_sm_d)}}.`$ (5.2) Here $`m_K`$ and $`m_\pi `$ are the physical kaon and pion masses, $`M`$ is the degenerate meson mass (corresponding to a degenerate quark mass $`m`$) used in the lattice computation of $`[K^+\pi ^+]_{1/2}`$, and $`M_K`$ and $`M_\pi `$ are the non-degenerate meson masses (corresponding to quark masses $`m_s`$ and $`m_d`$) used in the lattice computation of $`[K^00]`$. At tree level, the procedure is very simple, because the conversion involves only meson masses and the meson decay constant $`f`$, which can relatively easily be determined. To repeat a similar analysis at one loop, one would not only need to eliminate the $`O(p^2)`$ constants $`\alpha _{1,2}^8`$ and $`\alpha ^{27}`$, but also all the $`O(p^4)`$ LECs that can appear in the matrix elements for the kaon decays of interest. In other words, the $`O(p^4)`$ weak LECs $`\beta _{1,2,3,10,11,13,15}^8`$ and $`\beta _{1,2,4,5,6,7}^{27}`$ are needed , as well as the strong LECs $`\lambda _{4,5,6,8}`$. Only a few linear combinations of those can be determined from $`K\pi `$ and $`K0`$ matrix elements. For the $`K\pi \pi `$ matrix elements also only a few linear combinations are needed, but these are different linear combinations, involving also $`\beta _{13,15}^8`$ and $`\beta _{5,6}^{27}`$ which do not even appear in $`K\pi `$ and $`K0`$ at all. Also, we have seen that in unphysical matrix elements like $`K0`$ new LECs appear in these linear combinations, such as for instance $`\stackrel{~}{\beta }_1^8`$ in the linear combination $`\beta _1^82\beta _5^8+\stackrel{~}{\beta }_1^8`$ in Eq. (4.14). Likewise, one would be able to determine more LECs from $`K\pi `$ and $`K\eta `$ in the mass non-degenerate case, but also more unphysical (total-derivative) $`O(p^4)`$ operators would contribute, since these matrix elements would also not conserve energy for onshell external states, just as $`[K^00]`$. The only other way to determine more of the LECs from a lattice computation would be to consider more complicated correlation functions (such as $`K\pi \pi `$ itself). In that case, one necessarily has more than one strongly interacting particle in the initial or final states, and this leads to rather severe complications of its own (for $`K\pi \pi `$, see Refs. ). In general, on the lattice, one only has access to these matrix elements for unphysical choices of the kinematics . Again, not all relevant $`O(p^4)`$ LECs can be determined. We conclude that at one loop uncertainties are introduced in the determination of $`K\pi \pi `$ matrix elements from $`K\pi `$ and $`K0`$, which are not present at tree level. These uncertainties break down into two parts. One is the determination of $`\alpha _{1,2}^8`$, $`\alpha ^{27}`$ from $`K\pi `$ and $`K0`$, and the other is the conversion of results for these $`O(p^2)`$ LECs into $`K\pi \pi `$ decay rates. Here, we will only consider the first part, i.e. we will concentrate on the determination of $`\alpha _{1,2}^8`$ and $`\alpha ^{27}`$ from $`K\pi `$ and $`K0`$ matrix elements. If we are only interested in the determination of $`\alpha _{1,2}^8`$ and $`\alpha ^{27}`$, we need to know only the polynomial form (in the meson masses) of the contributions from $`O(p^4)`$ operators, as given in Eqs. (4.1,4.1,4.14,4.15), assuming that one-loop ChPT can be reliably applied to lattice computations of $`K\pi `$ and $`K0`$ matrix elements. The results of the previous section can be written in the generic form $`[K\pi ]_8`$ $`=`$ $`{\displaystyle \frac{4M^2}{f^2}}[\alpha _1^8(1+X_1^8)\alpha _2^8(1+X_2^8)+(\gamma _1^8+2\gamma _2^8)X^\gamma `$ $`+C_V^8{\displaystyle \frac{M^2}{(4\pi f)^2}}+C_S^8{\displaystyle \frac{NM_{SS}^2}{(4\pi f)^2}}],`$ $`[K\pi ]_{27}`$ $`=`$ $`{\displaystyle \frac{4M^2}{f^2}}\left[\alpha ^{27}\left(1+X^{27}\right)+C_V^{27}{\displaystyle \frac{M^2}{(4\pi f)^2}}+C_S^{27}{\displaystyle \frac{NM_{SS}^2}{(4\pi f)^2}}\right],`$ $`[K0]_8`$ $`=`$ $`{\displaystyle \frac{4i(M_K^2M_\pi ^2)}{f}}[\alpha _1^8Y_1^8+\alpha _2^8(1+Y_2^8)+\gamma _1^8Y_1^\gamma +\gamma _2^8Y_2^\gamma `$ $`+D_V^8{\displaystyle \frac{M_K^2}{(4\pi f)^2}}+D_S^8{\displaystyle \frac{NM_{SS}^2}{(4\pi f)^2}}],`$ $`[K0]_{27}`$ $`=`$ $`{\displaystyle \frac{4i(M_K^2M_\pi ^2)}{f}}\left[\alpha ^{27}Y^{27}+D^{27}{\displaystyle \frac{M_K^2M_\pi ^2}{(4\pi f)^2}}\right].`$ In these equations, $`X_{1,2}^8`$, $`X^{27}`$, $`Y_{1,2}^8`$, $`Y^{27}`$, $`X^\gamma `$ and $`Y_{1,2}^\gamma `$ stand for the one-loop contributions (“chiral logarithms”) given explicitly in Eqs. (4.5,4.6,4.12,4.13). The constants $`C_{V,S}^8`$ etc. are linear combinations of $`O(p^4)`$ LECs, which can be expressed in terms of the $`\beta `$’s and $`\lambda `$’s by comparison with Eqs. (4.1,4.1,4.14,4.15). From these relations $`\alpha _{1,2}^8`$ and $`\alpha _{27}`$ can be extracted by fitting these equations to lattice results for the matrix elements at various different values of the quark masses. With sufficient precision, also the $`O(p^4)`$ LECs could in principle be determined, but it is unlikely that this will work in practice with the currently available computational power. However, this does not imply that the $`O(p^2)`$ LECs cannot be determined with a reasonable accuracy. In order to get an idea about the size of one-loop effects, we will set the $`O(p^4)`$ coefficients $`C_{V,S}^{8,27}`$, $`D_{V,S}^8`$ and $`D^{27}`$ to zero, and evaluate the chiral logarithms at typical lattice values of the parameters and at $`\mathrm{\Lambda }=1`$ GeV, $`\mathrm{\Lambda }=m_\rho =770`$ MeV, and $`\mathrm{\Lambda }=m_\eta =550`$ MeV. We will consider three different “theories,” partially quenched with $`N=2`$ or $`3`$ and $`M_\eta ^{}`$ large, and quenched ($`N=0`$) with arbitrary $`\delta `$. We will also set $`f=f_\pi =132`$ MeV. We take $`M_{SS}=500`$ MeV, which corresponds to a sea-quark mass of about half the strange quark mass, vary the degenerate “lattice” meson mass $`M`$ at which $`[K\pi ]`$ is determined, and take $`2M_K^2/3=M_\pi ^2=M^2`$ for $`[K0]`$, which corresponds to $`m_s=2m_d=2m`$. It turns out that the one-loop correction $`X^{27}`$ for $`[K\pi ]_{27}`$ is very large. However, just as one defines the kaon $`B`$ parameter $`B_K`$ in the case of $`K^0\overline{K}^0`$ mixing, it makes sense to consider the ratio of this matrix element with its value evaluated by vacuum saturation or for large $`N_c`$, which is proportional to $`(M^2f^2)_{\mathrm{phys}}`$. In one-loop partially quenched ChPT, this can be expressed in terms of $`M`$ and $`f`$ as $`(M^2f^2)_{\mathrm{phys}}`$ $`=`$ $`M^2f^2\left(1+X_{\mathrm{vs}}+{\displaystyle \frac{2}{(\pi f)^2}}\left(M^2\lambda _8+NM_{SS}^2\lambda _6\right)\right),`$ (5.4) $`X_{\mathrm{vs}}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi f)^2}}[{\displaystyle \frac{2}{3}}(^2L(M)AM^2(2L(M)1))`$ $`+2NM_{VS}^{}{}_{}{}^{2}(L(M_{VS})1)].`$ For the ratio $$B_{27}=\frac{[K\pi ]_{27}}{(M^2f^2)_{\mathrm{phys}}}$$ (5.5) the relevant one-loop correction is $$\widehat{X}^{27}=X^{27}X_{\mathrm{vs}}=\frac{1}{(4\pi f)^2}\left[2M^2(3L(M)2)\right].$$ (5.6) In our examples, we will always consider the quantity $`\widehat{X}^{27}`$ instead of $`X^{27}`$. From Eqs. (4.6) and (5.4) we see that $`\widehat{X}^{27}`$ is independent of $`N`$, $`M_{SS}`$ and the $`\eta ^{}`$ parameters. In Tables 1 to 6 and Table 8 we computed the one-loop corrections for various choices of the valence-meson mass $`M`$ and three values of the $`\overline{MS}`$ scale $`\mathrm{\Lambda }=1`$ GeV, 770 MeV and 550 MeV. For $`N=0`$, we have set the parameters $`\alpha `$ and $`\gamma _{12,}^8`$ equal to zero, for simplicity. In Figs. 2, 4 and 6, we show the dependence of the three largest octet one-loop corrections, $`X_1^8`$ and $`Y_{1,2}^8`$ on $`M`$ and $`M_{SS}`$, for $`N=2`$ and large $`M_\eta ^{}`$. We do not show similar plots for $`X_2^8`$ and $`Y^{27}`$ because these one-loop corrections are typically much smaller. We also do not show a plot for $`\widehat{X}^{27}`$, because it only depends on $`M`$ and not on $`M_{SS}`$, nor on any of the $`\eta ^{}`$ parameters. These examples illustrate various points: * One-loop corrections can be substantial, and will have to be taken into account in order to obtain a reliable estimate for $`\alpha _{1,2}^8`$ and $`\alpha ^{27}`$ from the lattice. It may happen that the $`O(p^4)`$ LECs have values such that the combined, scale independent contribution of $`O(p^4)`$ non-analytic terms and counterterms are smaller than the non-analytic terms alone at a given scale $`\mathrm{\Lambda }1`$ GeV, improving the convergence of ChPT. One would hope this to be the case, especially for $`\widehat{X}^{27}`$, which is very large at the larger values of $`M`$ and $`\mathrm{\Lambda }`$, and for $`Y_1^8`$ at larger $`\mathrm{\Lambda }`$. This issue can be investigated on the lattice. * The size of the one-loop corrections grows with increasing $`\mathrm{\Lambda }`$. In particular, they are relatively small for $`\mathrm{\Lambda }=m_\eta 550`$ MeV. However, without further knowledge of $`O(p^4)`$ LECs, it is unnatural to use $`m_\eta `$ as the scale, because the $`\eta `$ is itself a Goldstone boson, with mass very close to the meson masses we are considering here. In the absence of information on $`O(p^4)`$ LECs, we believe that using higher values for $`\mathrm{\Lambda }`$ gives a better a priori estimate of the size of $`O(p^4)`$ effects. If, however, the values of $`O(p^4)`$ LECs turn out to be such that $`\mathrm{\Lambda }=550`$ MeV gives the better estimate, one-loop corrections would be reasonably small, and one-loop ChPT should be applicable in the computation of $`\alpha _{1,2}^8`$ and $`\alpha ^{27}`$ from the lattice at realistic values of the quark masses. * The fact that, in a number of cases, a one-loop correction becomes larger (and, in fact, diverges) with decreasing $`M`$ is a consequence of the fact that we do not vary $`M_{SS}`$ at the same time, i.e. of (partial) quenching. These “enhanced” chiral logarithms can be seen very clearly in the figures in the region $`M<M_{SS}`$, toward smaller $`M`$. Chiral logarithms are typically smaller near the line $`M=M_{SS}`$, for fixed $`M+M_{SS}`$. Enhanced chiral logarithms appear in all quantities except $`\widehat{X}^{27}`$. * So far, we have used the large-$`M_\eta ^{}`$ results, i.e. those of Subsect. 4.2, for all our examples. It is interesting to compare the $`N0`$ values obtained at a finite $`\eta ^{}`$ mass with those obtained for large $`M_\eta ^{}`$. In order to do this we chose $`\delta =0.18`$ for the finite-$`M_\eta ^{}`$ case. (In the real world $`\delta 0.18`$; for quenched QCD $`\delta 0.1`$ with a large error ). We set the other $`\eta ^{}`$-related parameters $`\alpha `$ and $`\gamma _{1,2}^8`$ equal to zero, simply because not much is known about their values (but see below for some remarks on contributions proportional to these parameters). For $`N=2`$, we show the difference between large and finite $`M_\eta ^{}`$ in Figs. 3, 5 and 7, for $`X_1^8`$ and $`Y_{1,2}^8`$. We show the ratios $`[1+X(\mathrm{finite}M_\eta ^{})]/[1+X(\mathrm{large}M_\eta ^{})]`$, with $`X=X_1^8,Y_2^8`$, and the difference $`Y_1^8(\mathrm{finite}M_\eta ^{})Y_1^8(\mathrm{large}M_\eta ^{})`$ (there is no tree-level contribution proportional to $`\alpha _1^8`$). The first thing to be noticed is that there are still enhanced chiral logarithms in the small $`M`$, large $`M_{SS}`$ region. This is because the coefficients of the enhanced chiral logarithms depend on $`M_\eta ^{}`$. In the small $`M`$, $`M_{SS}`$ region, the ratios are close to one (or, the difference is close to zero), as one would expect if both $`M`$ and $`M_{SS}`$ are small compared to $`M_\eta ^{}`$. For $`N=2`$, $`M_\eta ^{}=863`$ MeV, and $`M_{SS}^2/M_\eta ^{}^20.34`$ for $`M_{SS}=500`$ MeV, while for $`N=3`$, $`M_\eta ^{}=996`$ MeV, so that $`M_{SS}^2/M_\eta ^{}^20.25`$ for $`M_{SS}=500`$ MeV. We point out that for larger meson masses, the plots are less meaningful in the region $`M>M_{SS}`$. This is a consequence of the fact that we expanded the finite-$`M_\eta ^{}`$ results in $`M^2/M_\eta ^{}^2`$, but not in $`M_{SS}^2/M_\eta ^{}^2`$. This makes sense if $`M<M_{SS}`$, but not for $`MM_{SS}`$. Therefore, for valence-quark masses close to the sea-quark mass, one either has to work consistently with the large-$`M_\eta ^{}`$ results, or a more general analysis using Eq. (3.4) instead of Eq. (3.8) is needed. The latter is outside the scope of the present paper, but for $`M_{SS}^2/M_\eta ^{}^20.250.34`$ it appears reasonable to use the large-$`M_\eta ^{}`$ results of Subsect. 4.2. As a further illustration of this point, Table 7 shows the values of $`^2`$ and $`A`$ calculated from Eq. (3.8) (3rd, resp. 6th columns), the “exact” expression Eq. (3.4) with degenerate valence quark masses (4th and 7th columns), and in the limit of large $`M_\eta ^{}`$, Eq. (3.7) (5th and 8th columns). We see that the variation of values for $`^2`$ and $`A`$ is at most about 20, resp. 50%. This is not very large from a practical point of view: the current best determination of $`m_0^2=^2(N=0)`$ in the quenched theory has an error of the same order. We note that for $`M=350`$ MeV the values of $``$ and $`A`$ are already much closer to their “exact” values $`(M)`$ and $`A(M)`$ than $`_{\mathrm{}}`$ and $`A_{\mathrm{}}`$. All this means that, for the value of parameters chosen in our examples, one could use the simplest possible form of the chiral logarithms, given in Subsect. 4.2, to fit numerical results. Tables 1 to 6 then give the relevant examples of the size of non-analytic one-loop corrections. Only with numerical results so precise that one would be able to determine $`^2`$ and $`A`$ with better precision would it be important to take the dependence on $`\eta ^{}`$ parameters into account. Note again that these conclusions do depend on the values of the meson masses (and hence quark masses) we considered in our examples. * We also considered the effect of the $`\eta ^{}`$ couplings $`\gamma _{1,2}^8`$. They contribute to the octet matrix elements at one loop through the chiral logarithms $`X^\gamma `$ and $`Y_{1,2}^\gamma `$. In Table 8 we show the size of these one-loop corrections for the choice $`\mathrm{\Lambda }=1`$ GeV. We see that these quantities are not very small, even though they vanish in the limit in which the $`\eta ^{}`$ decouples (for $`N0`$). They are typically smaller than $`X_{1,2}^8`$ and $`Y_{1,2}^8`$. If the couplings $`\gamma _{1,2}^8`$ themselves are small as well (compared to $`\alpha _{1,2}^8`$), it may be possible to ignore these terms. But if these couplings are not small, this would mean that the $`\eta ^{}`$ does play a role, even if we can use the large-$`M_\eta ^{}`$ partially-quenched expressions of Subsect. 4.2 for the other chiral logarithms. The results are not very sensitive to $`\alpha `$, except in the quenched case ($`N=0`$). We expect that, within the precision of current lattice computations, the effect of (arbitrarily) setting $`\alpha `$, as well as $`\gamma _{1,2}^8`$, to zero, can be accommodated in fits by shifts in the $`O(p^4)`$ LECs $`C_V^8`$ etc. in Eq. (5), without significant impact on the values of the $`O(p^2)`$ LECs $`\alpha _{1,2}^8`$ and $`\alpha ^{27}`$. This can be checked by fitting lattice results while constraining these parameters to a few different values. ## 6 Conclusion We presented a complete analysis of $`K\pi `$ and $`K0`$ weak matrix elements in one-loop ChPT for partially quenched QCD with $`N`$ degenerate sea quarks. For $`K\pi `$ we took the valence quarks degenerate in mass, while for $`K0`$ they are kept non-degenerate in order to get a non-trivial result. Three cases have been considered. The first is a partially quenched theory with a valence-meson mass much smaller than the $`\eta ^{}`$ mass, but with arbitrary value of $`M_{SS}/M_\eta ^{}`$, so that an expansion in powers of $`M_{kk}^2/M_\eta ^{}^2`$ is allowed. The second choice corresponds to the large $`M_\eta ^{}`$ limit where the $`\eta ^{}`$ decouples, so that we can also expand in $`M_{SS}^2/M_\eta ^{}^2`$. The third case is the quenched theory obtained for $`N=0`$ or $`M_{SS}\mathrm{}`$. For completeness, we also included results for the unquenched theory. These results should be useful for extracting the values of the LECs $`\alpha _{1,2}^8`$ and $`\alpha ^{27}`$ from Lattice QCD. As we emphasized already in the Introduction, these are interesting quantities in their own right. Estimates extracted from experiment exist, with which lattice results can be compared. The matrix elements considered here are the simplest weak matrix elements which can be used for this goal. The expressions to be used in fits to lattice data are given in Eq. (5), in which $`X_{1,2}^8`$, $`X^{27}`$, $`Y_{1,2}^8`$, $`Y^{27}`$, $`X^\gamma `$ and $`Y_{1,2}^\gamma `$ represent one-loop corrections (chiral logarithms), and $`C_{V,S}^{8,27}`$, $`D_{V,S}^8`$ and $`D^{27}`$ are linear combinations of $`O(p^4)`$ LECs. They can be read off from the explicit one-loop results in Sects. (4.1-4). From our numerical examples in Sect. 5 it is clear that one-loop expressions will be needed for typical values of quark masses used in present lattice computations. Contributions from $`O(p^4)`$ operators, represented by the LECs $`C_{V,S}^{8,27}`$, $`D_{V,S}^8`$ and $`D^{27}`$, will also need to be included. Theoretically, values for these $`O(p^4)`$ LECs can also be extracted from lattice computations. However, realistically, we expect that such estimates would have large uncertainties, both as a consequence of the typical statistics of present lattice computations, as well as because of uncertainties related to the role of the $`\eta ^{}`$ discussed in more detail in Sect. 5. We note that the $`O(p^4)`$ LEC $`\beta _1^{27}`$ can be accessed directly by a computation of the ratio of the $`[K0]_{27}`$ and the $`K^0\overline{K}^0`$-mixing matrix elements (cf. Eq. (5)). In practice, for reasonably small values of the sea-quark mass (of order less than one-half times the strange-quark mass), it may be possible to use the partially quenched results for large $`M_\eta ^{}`$, given in Subsect. 4.2. In this case the $`\eta ^{}`$ decouples, and therefore all dependence on $`\eta ^{}`$ parameters, $`m_0`$, $`\alpha `$ and $`\gamma _{1,2}^8`$ is removed, making the analysis simpler. In addition, it is only in this limit that estimates obtained for $`O(p^4)`$ LECs can be directly compared to those of the real world, provided that one chooses $`N=3`$ sea quarks. For a more detailed discussion, see Sects. 3 and 5. The more general results for the case that the sea-quark mass is larger, comparable to the $`\eta ^{}`$ mass, but the valence-quark mass is still small enough, are given in Subsect. 4.1. A completely quenched lattice computation (for which the relevant results are in Subsect. 4.3) should be useful for assessing the feasibility of this approach. An $`N=2`$ computation, in combination with a quenched computation, could give insight into the dependence of the LECs on the number of light flavors. However, since we do not know the functional form of the $`N`$ dependence of the (finite part of the) LECs, an $`N=3`$ computation will be needed to obtain estimates without an uncontrolled systematic error. The emphasis of this paper is on the extraction of reliable numbers for $`\alpha _1^8`$ and $`\alpha ^{27}`$ from lattice computations. While these $`O(p^2)`$ LECs are interesting, because of the availability of phenomenological estimates, the final aim of such Lattice QCD computations would be to convert these numbers into quantitative estimates of the $`\mathrm{\Delta }I=1/2`$ and $`\mathrm{\Delta }I=3/2`$ $`K2\pi `$ decay amplitudes. This can be done using ChPT, and complete $`O(p^4)`$ formulae for doing so are given in Ref. . Assuming that $`O(p^4)`$ ChPT is precise enough, the largest uncertainty arises because of the poor knowledge of all needed $`O(p^4)`$ LECs. Many of these, as we discussed in Sect. 5, cannot even in principle be determined from the $`K\pi `$ and $`K0`$ matrix elements. (More $`O(p^4)`$ LECs are accessible through the $`K^0\overline{K}^0`$ and $`K2\pi `$ matrix elements with both pions at rest , the computation of which can also serve as a check on lattice results for $`\alpha _1^8`$ and $`\alpha ^{27}`$.) One would have to resort either to the use of available phenomenological information , or to theoretical estimates based on arguments such as large $`N_c`$ or models (for recent discussions see Refs. ). In addition, it is well known that final-state interactions are responsible for a large enhancement of the $`I=0`$ $`K\pi \pi `$ amplitude . In that case it could be necessary to resum those effects instead of relying on an $`O(p^4)`$ ChPT calculation. A possible way of resumming final-state interactions has recently been proposed in Ref. . #### Acknowledgements We would like to thank Claude Bernard, Steve Sharpe as well as Akira Ukawa and other members of the CP-PACS collaboration for very useful discussions. MG would like to thank the Physics Departments of the Università “Tor Vergata,” Rome, the Universitat Autonoma, Barcelona, and the University of Washington, Seattle, for hospitality. MG is supported in part by the US Department of Energy, and EP by the Ministerio de Educaçion y Cultura of Spain.
warning/0006/quant-ph0006107.html	ar5iv	text	# Untitled Document Group-theoretical Structure of the Entangled States of N Identical Particles Suranjana Rai and Jagdish Rai RAITECH M - 12/28 DLF CITY GURGAON 122002 INDIA 2000 June 23 Abstract We provide a group-theoretical classification of the entangled states of N identical particles. The connection between quantum entanglement and the exchange symmetry of the states of N identical particles is made explicit using the duality between the permutation group and the simple unitary group. Each particle has n-levels and spans the n-dimensional Hilbert space. We shall call the general state of the particle as a qunit. The direct product of the N qunit space is given a decomposition in terms of states with definite permutation symmetry. The nature of fundamental entanglement of a state can be related to the classes of partitions of the integer N. The maximally entangled states are generated from linear combinations of the less entangled states of the direct product space. We also discuss the nature of maximal entanglement and its measures. 1 Introduction We give a fundamental group theoretical result, which explains quantum entanglement in terms of the exchange symmetry of any number N of identical particles, as well as the exchange symmetry of the levels of those particles. We claim that there is a requirement of symmetry for entanglement to occur. We give symmetrization procedures to generate entangled states starting from non-entangled states. We see that maximum symmetry leads to maximum entanglement, i.e. states which are maximally symmetric with respect to the number of particles as well as to the number of levels, also have the maximum possible entanglement . We proceed in two steps. Firstly, the particle symmetrization procedure gives rise to ordered subspaces of various types of entanglement. In the second step, level symmetrization leads to combinations of these ordered states which are conjugate to each other in level excitations. We thus get the maximally entangled states we seek, the states with maximum symmetry. The power of this procedure is that all possible types of entanglement, i.e. all N particle entanglements, and their various combinations, can be generated and classified. The number of maximally entangled states is equal to the dimension of the Hilbert space, and they are orthonormal. We are thus able to span the entire Hilbert space of the particles using these ordered maximally entangled states as basis states. Indeed, the ordering can be taken to be a measure of the entanglement. We discuss the meaning and nature of maximum entanglement. In the ordered maximally entangled basis, by maximal entanglement we mean maximum possible entanglement for each different type of entanglement. A useful measure of maximal entanglement is the criterion of concurrence . We illustrate these concepts with the help of examples. Another quantity of interest, is the N particle maximum entanglement which is measured by the maximum one-particle entropy obtained by tracing over the rest of the particles. We have also generated N-particle maximally entangled states by taking proper linear combinations of the different types of entanglement. It is also possible to use these states as a basis to span the Hilbert space. We also obtain the N-1, N-2, etc. particle entanglements in terms of entropy criteria. 2 Entanglement of N-identical particles We now discuss the general case of a system of N identical particles. Each particle has n levels. By analogy with a qubit, we call it a qunit. In this work we consider the case of pure states only. A qunit is written as $`i_1,i_2,i_3,\mathrm{},i_N`$ where each $`i_1,i_2,i_3,\mathrm{},i_N`$ take values from 1 to n. This state spans an n dimensional Hilbert space H for one particle. For N particles, the space H<sup>N</sup> spanned has a dimension of n<sup>N</sup>. We are interested in obtaining the structure of the direct product space of N particles in terms of entanglement. We use the methods of group theory, to decompose the n<sup>N</sup> dimensional space of the direct product of the of the spaces of each particle into a direct sum of spaces . Each constituent of the direct sum has a definite permutational symmetry. This is quite naturally expected since the system contains identical particles and permutations of the particles leave the system invariant. Each of these constituents of the direct sum forms a representation of the permutation group of N particles, also known as the symmetric group S<sub>N</sub>. For the case of N, n-level atoms, the state space can be written as $`SU(n)\times SU(n)\times \mathrm{}`$ SU(n) (N copies),(2.1) The group SU(n) consists of all n x n unitary matrices. Each SU(n) group describes the states of a single n-level atom. As there are N copies in the direct product, it is possible to decompose the above direct product as $`S_N\times SU(n)`$ ,(2.2) We wish to generate states from this space that have a definite permutational symmetry and are maximally entangled for that set. The full space can be given a decomposition using the various representations of the permutation group as described above. For this purpose it is necessary to consider the representations of the symmetric group $`S_N`$ of all N! permutations, given in detail in References 3 and 4. We choose a basis of H<sup>N</sup> whose elements transform simultaneously as would a basis for an irreducible representation of SU(n) and a basis for an irreducible representation of $`S_N`$. The decomposition is possible using the well known result $`H^N=S^\lambda \times T^\lambda `$, where $`\lambda `$ is a partition of N treated as a non-increasing ordered k-tuple of positive integers with sum N. Also $`S^\lambda `$ is a representation of the group $`S_N`$ and $`T^\lambda `$is a representation of H. The states of the direct product space are taken to be the basis states of the various irreducible representations occurring in the direct product space $`H^N`$. It is interesting to note that the frequency of the $`S^\lambda `$ is equal to the dimensionality of the $`T^\lambda `$ and the frequency of the $`T^\lambda `$ is equal to the dimensionality of $`S^\lambda `$. A simple formula for obtaining the dimensionality of the $`T^\lambda `$ can be obtained from the possible ways of filling the Young’s diagram. Now we classify the n<sup>N</sup> states based on the symmetry which corresponds to a definite Young’s diagram. We call these the collective states of the system of N identical particles. This procedure produces a hierarchy of states that are ordered according to the degree of entanglement in subspaces of definite symmetry. We explain this using several examples for two and three particles in Section 4. For a collection of two-level atoms the states can be ordered with the quantum numbers j and m. The state with m = j, is unentangled and the degree of entanglement increases as the m value decreases from j to zero or half, and then starts decreasing again corresponding to a bell shaped curve. The procedure for generating maximally entangled states from states of lower entanglement is to combine conjugate states i.e., states with fixed Casimir operator value but opposite sign for values of the diagonal generators . The conjugate states are states related by spin flip operators and are connected by local operations . In this way, one can generate all possible states of the system with a definite degree of entanglement (Table given below). We find that the classification obtained on symmetry group lines is much more powerful than other methods as we are able to generate all possible types of entanglement for a system of many particles. It is also possible to generalize our methods to non-identical particles. Our procedure generates states that are connected by non-local operations. There has been extensive work done to understand entanglement by local unitary operations . However, non-locality is at the heart of entanglement. We feel one should not be limited to local operations only. In fact, non-local operations have already been found to be useful e.g. in bound entanglement. Our procedure in fact is powerful enough to introduce both local and non-local operations. The general case, involves non-local or joint unitary operations where the unitary operator cannot be decomposed in the form given by Eq. 2.1. It is necessary to have joint unitary operations as we are increasing entanglement by our procedures. Entanglement does not increase in local operations. However, in the specific case, where the particles are distinguishable, like the case where they are at different locations, it is not necessary to symmetrize. Then, there is no need for the permutation operator, the total Hilbert space can be decomposed into the space of local unitary operations each represented by SU(n). The structure of H<sup>N</sup> has been considered by several authors. Werner has used the idea of symmetric subspaces of H<sup>N</sup> and applied it to the optimal cloning of pure states. J. Eisert et. al., have found a class of mixed states with known distillable entanglement , using techniques that employ decomposition of the product space into direct sum space. Cirac et. al., also consider similar decompositions for the SU(2) group in the context of the two-level atom space . In our paper, we provide an understanding of the necessity of symmetrization, and a group theoretical justification of the mathematical basis of the above works. 3 Nature of Entanglement The methods of group theory give a very simple and elegant description of the nature of entanglement for a system of N identical particles. The different types of entanglements can be related to the various partitions of the integer N. These partitions form a class . For example in the case of the group S<sub>2</sub>, there are two classes, identity, (e) and the two-cycle permutation (12). Essentially for two particles, there are two possibilities \- either the particles are entangled or not. There is a very easy geometric way to describe this as in the top line of the figure. For three particles, we have three classes for the group S<sub>3</sub> : (e), {(12),(23),(13)} and {(123), (132)}. Geometrically , this can be shown as In the above diagrams, the line attaching the two particles shows entanglement while the two particles unconnected show unentangled particles. The second line of the figure shows the fundamental entanglements which represent the classes. The figures at the bottom show the possible combinations of the fundamental entanglements which have also been discussed in Section 4. This description can easily be generalized to a system of N particles. It is well known that the number of representations is equal to the number of classes (Schur’s lemma). According to this description, each class can be put in correspondence with a unique representation of the group. In this way, the various types of entanglement which are related to the representations, are connected to the various classes. 4 Examples We now consider some specific examples of the bipartite and tripartite system to illustrate our procedure. Example 1. Bipartite two-level system - 2 qubits For a bipartite two-level system, the state space consists of $`SU(2)\times SU(2)`$. The four possible states of the system are $`11,12,21and22`$. This space is four-dimensional (2<sup>N</sup> = 4). The well-known symmetrization procedure which gives the singlet and triplet states is as follows. $`j,m`$ Symmetric $`1,1=11`$ $`1,0=1/\sqrt{2}(12+21)`$ $`1,1=22`$ $`j,m`$ Antisymmetric $`0,0=1/\sqrt{2}(1221)`$ This decomposition is on the lines of $`S_2\times SU(2)`$. The group S<sub>2</sub> has two irreducible representations. The three symmetric states form the basis of one irreducible representations. The antisymmetric states correspond to the other representation. With this break up, the decomposition is complete. The four-dimensional product space of $`SU(2)\times SU(2)`$ has been broken into the direct sum of $`S_2\times SU(2)`$ as $`2\times 2`$= 3 + 1. In this classification, there are two unentangled states and two maximally entangled states. If we also symmetrize with respect to the levels then we need to add/subtract to the $`11`$ state the $`22`$ state and normalize. The $`22`$ state could be called the conjugate state to the $`11`$ state and vice versa. This generates two additional states which are maximally entangled. Thus we have generated the full set of all four maximally entangled states: $`1/\sqrt{2}(11+22)`$ $`1/\sqrt{2}(1122)`$ $`1/\sqrt{2}(12+21)`$ $`1/\sqrt{2}(1221)`$ Example 2. Tripartite two-level system - 3 qubits We now extend this procedure to the case of the tripartite system. The state space consists of $`SU(2)\times SU(2)\times SU(2)`$ and is of dimension $`2\times 2\times 2`$= 8. The decomposition proceeds on the lines of $`S_3\times SU(2)`$. The group S<sub>3</sub> has three irreducible representations which correspond to the three types of states in the state space. First, we have the symmetric representation which has the 4 symmetric wave functions as the basis. Then we have two states of mixed symmetry which correspond to two degenerate representations, each of dimension 2. There are no completely antisymmetric states for a system of three particles each with two levels. This is essentially a statement of the Pauli exclusion principle in terms of the group representations. These are the only possible states for this system ( $`2\times 2\times 2`$ = 4 + 2 + 2 + 0 ). In this case, the classification based on symmetry proceeds as follows. There are two sets of states which are fully symmetric and of mixed symmetry. $`j,m`$ Symmetric $`3/2,3/2=111`$ $`3/2,1/2=1/\sqrt{3}(112+121+211)`$ $`3/2,1/2=1/\sqrt{3}(221+212+122)`$ $`3/2,3/2=222`$ $`j,m;d`$ Mixed symmetry $`1/2,1/2,1=1/\sqrt{6}(2211112121)`$ $`1/2,1/2,1=1/\sqrt{6}(212+2212122)`$ and $`1/2,1/2,2=1/\sqrt{2}(112121)`$ $`1/2,1/2,2=1/\sqrt{2}(221212)`$ In the symmetric subspace, $`111`$ and $`222`$ are unentangled and are conjugate to each other. The states with $`j,m;3/2,1/2`$ and $`3/2,1/2`$ are also conjugate to each other but have partial entanglement (with non-zero one particle entropy). The conjugate state is generated by the action of the flip operators on the original state and are obtained by interchanging 1 and 2. To produce maximally entangled states by level symmetrization we add and subtract these states. The states are labeled by $`j,m;d`$ where j = N/2, N/2 -1,1/2 or 0, and m is the magnetic quantum number and d denotes the representation. For the three particle case of mixed symmetry there are two degenerate representations. In general the state $`j,m`$ is conjugate to $`j,m`$ and their addition and subtraction produces a maximally entangled state . It is maximally entangled in the sense that the concurrence is one. It is well known that concurrence in itself can be used as a measure of entanglement . For example, the Bell diagonal states in the case of a bipartite system have a concurrence equal to one and are maximally entangled. The unentangled pure states have a concurrence of zero. The degree of entanglement in a manifold of states of a certain specific symmetry increases from zero to higher entanglement in the middle of the ladder of states and then decreases. The states with lower m have higher degree of entanglement. It is interesting to note, however, that the concurrence for each conjugate by itself is zero. There is no fully antisymmetric state for the three spins and the procedure stops here generating eight possible maximally entangled states (with proper normalizations): $`3/2,3/2+3/2,3/2`$ , $`3/2,3/23/2,3/2`$ , $`3/2,1/2+3/2,1/2`$ , $`3/2,1/23/2,1/2`$ , $`1/2,1/2;1+1/2,1/2;1`$ $`1/2,1/2;11/2,1/2;1`$ $`1/2,1/2;2+1/2,1/2;2`$ $`1/2,1/2;21/2,1/2;2`$ We now discuss the nature of entanglement of these states. From the first pair of conjugate states of the three particle states we get the GHZ state: $`GHZ=3/2,3/2+(\pm )3/2,3/2=111+(\pm )222`$, (4.2.1) For the GHZ state the single particle entropy is maximum and thus it is truly three particle maximal entangled. From the next pair of states we get the following state that we call the Z-state $`Z=1/\sqrt{3}(112+121+211)`$, $`Z+\overline{Z}=3/2,1/2+3/2,1/2`$ $`=\{(1_A+2_A)\psi ^+_{BC}\}+\{2_A11_{BC}+1_A22_{BC}\}`$ (4.2.2) The above state consists of two parts. The first part shows two particle maximal entanglement in B and C (Bell diagonal state) while A is unentangled. The two-particle entanglement, could be between any two of the particles, shown by suitable algebraic manipulation. The second part is a state of the GHZ type. This shows that there is both two particle and three particle entanglement. We also define $`Y=1/\sqrt{6}(2211112121)`$ $`X=1/\sqrt{2}(112121)`$ Combining the conjugates $`Y+\overline{Y}=\psi ^{}_{AC}(1_B+2_B)+\psi ^{}_{AB}(1_C+2_C)`$ (4.2.3) The above state is a sum of two two-particle maximally entangled states with one particle separable from each. Again, any one particle can be separated out depending on the representation we start from. $`X+\overline{X}=(1_A+2_A)\psi ^{}_{BC}`$ (4.2.4) This state is a product state of a one particle unentangled and a two-particle maximal entangled state (Bell diagonal). So we see from Eqs. (4.2.1 - 4.2.4), that we obtain all types of maximal entanglement of their own kind. The number of maximally entangled states is exactly equal to the dimension of the space and they are mutually orthogonal. The maximally entangled states could form a useful basis. It is very interesting to note, that further linear combinations of these states give rise to fundamental entanglements. We first take the case of the states, from $`(Y+\overline{Y})+1/\sqrt{3}(X+\overline{X})`$, Eqs (4.2.3-4.2.4). We see that these states are a sum of three , two-particle maximally entangled terms. We have verified that this combination is three-particle maximally entangled i.e. its single particle entropy is maximum. Thus it is in a way comparable to the GHZ state which cannot be broken down into two-particle entanglements. Next, we take the states $`Z+\overline{Z}GHZ`$ , Eqs (4.2.1-4.2.2). We find that it is the product state of a two-particle maximal entangled (Bell diagonal) state and a one particle unentangled state. Here we see that by combining three-particle entangled states we can get two-particle entanglement . Example 3. N-partite two-level system It is straight forward to extend this procedure to the case of N spins. The state space corresponds to SU(2)X SU(2)X SU(2) (N copies), Each SU(2) group describes the states of the two-level atoms. As there are N copies in the direct product, it is possible to decompose the above direct product as S<sub>N</sub> X SU(2). There are 2<sup>N</sup> maximally entangled states which can be generated using our procedure. We combine the states $`j,m;d+j,m;d`$, where j = N/2, N/2 -1,…,1/2 or 0. Each possible value of j, corresponds to a set of states with a definite symmetry in the exchange of particles. The maximum value of j = N/2, corresponds to the fully symmetric states which are 2j + 1 = N+1 in number. The next value of j = N/2 - 1, gives states which are symmetric in the exchange of N-1 particles and antisymmetric in the exchange of two particles. This procedure is continued to obtain a complete decomposition of the states of the entire system. The states for the N qubit system, $`j,m`$ are shown in the table below. j=N/2; $`N/2,N/2,N/2,N/21,\mathrm{},N/2,N/2`$ j=N/2 - 1; $`N/21,N/21,N/21,N/22,\mathrm{},N/21,(N/21)`$ till … j=1/2 or 0. The angular momentum states $`j,m`$ are well known through Dicke’s work on superradiance . They form a complete and orthogonal basis. The way we have generated the entangled states can be related to Dicke’s interpretation of superradiance. The conjugate state corresponds to a value of m which is negative of the original m. Entanglement is maximum for low values of m which is well known to correspond to the superradiant state which occurs for the value of m around zero. The entangled state thus could be interpreted as the state where all the atoms of the radiating system are maximally interacting with each other. During superradiant emission the entanglement is transferred to the radiation called swapping of entanglement from the atom to the photon. Acknowledgments We would like to thank N. Mukunda, P. Zoller, I. Cirac, G. Vidal, N. Cerf and A. Zeilinger for discussions. This work was supported in part by the Erwin Schrodinger Institute, Vienna. References 1. J. Rai and S. Rai, e-print quant-ph/003055. 2. William K. Wooters, Phys. Rev. Lett. 80, 2245 (1998). 3. Jagdish Rai, C. L. Mehta and N. Mukunda, J. Math. Phys. 29, 510 (1988). 4. Jagdish Rai, C. L. Mehta and N. Mukunda, J. Math. Phys. 29, 2443 (1988). 5. N. Linden and S. Popescu, e-print quant-ph/9711016. 6. R. F. Werner, Physical Review A, 58, 1827 (1998). See also, D. Bruss, A. Ekurt and C. Macchiavello, Phys. Rev. Lett. 81, 2598 (1998). 7. J. Eisert et. al., Phys. Rev. Lett., 84, 1611 (2000). 8. J. I. Cirac, A. K. Ekert and C. Macchiavello, Phys. Rev. Lett., 82, 4344 (1999). 9. N. Gisin and H. Bechmann-Pasquinucci, e-print quant-ph/9804045. 10 . R. H. Dicke, Phys. Rev. A 93, 99 (1954).
warning/0006/hep-th0006211.html	ar5iv	text	# A new holographic entropy bound from quantum geometry ## Abstract A new entropy bound, tighter than the standard holographic bound due to Bekenstein, is derived for spacetimes with non-rotating isolated horizons, from the quantum geometry approach in which the horizon is described by the boundary degrees of freedom of a three dimensional Chern Simons theory. The Holographic Principle (HP) \- and the holographic Entropy Bound (EB) have been the subject of a good deal of attention lately. In its original form , , the HP asserts that the maximum possible number of degrees of freedom within a macroscopic bounded region of space is given by a quarter of the area (in units of Planck area) of the boundary. This takes into account that a black hole for which this boundary is (a spatial slice of) its horizon, has an entropy which obeys the Bekenstein-Hawking area law and also the generalized second law of black hole thermodynamics . Given the relation between the number of degrees of freedom and entropy, this translates into a holographic EB valid generally for spacetimes with boundaries. The basic idea underlying both these concepts is a network, at whose vertices are variables that take only two values (‘binary’, ‘Boolean’ or ‘pixel’), much like a lattice with spin one-half variables at its sites. Assuming that the spin value at each site is independent of that at any other site (i.e., the spins are randomly distributed on the sites), the dimensionality of the space of states of such a network is simply $`2^p`$ for a network with $`p`$ vertices. In the limit of arbitrarily large $`p`$, such a network can be taken to approximate the macroscopic surface alluded to above, a quarter of whose area bounds the entropy contained in it. Thus, any theory of quantum gravity in which spacetime might acquire a discrete character at length scales of the order of Planck scale, is expected to conform to this counting and hence to the HP. Let us consider now a slightly altered situation: one in which the binary variables at the vertices of the network considered are no longer distributed randomly, but according to some other distribution. Typically, for example, one could distribute them binomially, assuming, without loss of generality, a large lattice with an even number of vertices. Consider now the number of cases for which the binary variable acquires one of its two values, at exactly $`p/2`$ of the $`p`$ vertices. In case of a lattice of spin 1/2 variables which can either point ‘up’ or ‘down’, this corresponds to a situation of net spin zero, i.e., an equal number of spin-ups and spin-downs. Using standard formulae of binomial distributions, this number is $$N(\frac{p}{2}\|a)=2^p\left(\begin{array}{c}p\\ p/2\end{array}\right)[a(1a)]^{p/2},$$ (1) where, $`a`$ is the probability of occurrence of a spin-up at any given vertex. Clearly, this number is maximum when the probability of occurrence $`a=1/2`$; it is given by $`p!/(\frac{p}{2}!)^2`$. Thus, the number of degrees of freedom is now no longer $`2^p`$ but a smaller number. This obviously leads to a lowering of the entropy. For very large $`p`$ corresponding to a macroscopic boundary surface, this number is proportional to $`2^p/p^{\frac{1}{2}}`$. The new EB can therefore be expressed as $$S_{max}=ln\left(\frac{\mathrm{exp}S_{BH}}{S_{BH}^{1/2}}\right),$$ (2) where, $`S_{BH}=A_H/4l_P^2`$ is the Bekenstein-Hawking entropy. This is a tighter bound than that of ref. mentioned above. The ‘tightening’ of holographic EB is the subject of this paper. We shall show below that, in the quantum geometry framework, it is possible to have an even tighter bound than that depicted in eq. (2). There are of course examples of situations where the EB is violated , and must be generalized. However, generalizations proposed so far appear to be tied to fixed classical background spacetimes, and may not hold when gravitational fluctuations are taken into account . In this note, we restrict ourselves to the older version of the EB appropriate to stationary spacetimes, but with allowance for the existence of radiation in the vicinity of the boundary. In this sense, the appropriate conceptual framework is that of the Isolated Horizon . We consider generic 3+1 dimensional isolated horizons without rotation, on which one assumes an appropriate class of boundary conditions. These boundary conditions require that the gravitational action be augmented by the action of an $`SU(2)`$ Chern-Simons theory living on the isolated horizon . Boundary states of the Chern-Simons theory contribute to the entropy. These states correspond to conformal blocks of the two-dimensional Wess-Zumino model that lives on the spatial slice of the horizon, which is a 2-sphere of area $`A_H`$. The dimensionality of the boundary Hilbert space has been calculated thus - by counting the number of conformal blocks of two-dimensional $`SU(2)_k`$ Wess-Zumino model, for arbitrary level $`k`$ and number of punctures $`p`$ on the 2-sphere. We shall show, from the formula for the number of conformal blocks specialized to macroscopic black holes characterized by large $`k`$ and $`p`$ , that the restricted situation described above, ensues, thus realizing a more stringent EB. We may mention that similar ideas relating the quantum geometry approach to the HP and EB have been pursued by Smolin , although, as far as we understand, the issue of tightening the Bekenstein bound has not been addressed. We start with the formula for the number of conformal blocks of two-dimensional $`SU(2)_k`$ Wess-Zumino model that lives on the punctured 2-sphere. For a set of punctures $`𝒫`$ with spins $`\{j_1,j_2,\mathrm{}j_p\}`$ at punctures $`\{1,2,\mathrm{},p\}`$, this number is given by $$N^𝒫=\frac{2}{k+2}\underset{r=0}{\overset{k/2}{}}\frac{\underset{l=1}{\overset{p}{}}sin\left(\frac{(2j_l+1)(2r+1)\pi }{k+2}\right)}{\left[sin\left(\frac{(2r+1)\pi }{k+2}\right)\right]^{p2}}.$$ (3) Observe now that Eq. (3) can be rewritten as a multiple sum, $$N^𝒫=\left(\frac{2}{k+2}\right)\underset{l=1}{\overset{k+1}{}}sin^2\theta _l\underset{m_1=j_1}{\overset{j_1}{}}\mathrm{}\underset{m_p=j_p}{\overset{j_p}{}}\mathrm{exp}\{2i(\underset{n=1}{\overset{p}{}}m_n)\theta _l\},$$ (4) where, $`\theta _l\pi l/(k+2)`$. Expanding the $`\mathrm{sin}^2\theta _l`$ and interchanging the order of the summations, this becomes $$N^𝒫=\underset{m_1=j_1}{\overset{j_1}{}}\mathrm{}\underset{m_p=j_p}{\overset{j_p}{}}\left[\overline{\delta }_{(_{n=1}^pm_n),0}\frac{1}{2}\overline{\delta }_{(_{n=1}^pm_n),1}\frac{1}{2}\overline{\delta }_{(_{n=1}^pm_n),1}\right],$$ (5) where, we have used the standard resolution of the periodic Kronecker deltas in terms of exponentials with period $`k+2`$, $$\overline{\delta }_{(_{n=1}^pm_n),m}=\left(\frac{1}{k+2}\right)\underset{l=0}{\overset{k+1}{}}\mathrm{exp}\{2i[(\underset{n=1}{\overset{p}{}}m_n)m]\theta _l\}.$$ (6) Our interest focuses on the limit of large $`k`$ and $`p`$, appropriate to macroscopic black holes of large area. Observe, first of all, that as $`k\mathrm{}`$, the periodic Kronecker delta’s in (6) reduce to ordinary Kronecker deltas, $$\underset{k\mathrm{}}{lim}\overline{\delta }_{m_1+m_2+\mathrm{}+m_p,m}=\delta _{m_1+m_2+\mathrm{}+m_p,m}.$$ (7) In this limit, the quantity $`N^𝒫`$ counts the number of $`SU(2)`$ singlet states, rather than $`SU(2)_k`$ singlets states. For a given set of punctures with $`SU(2)`$ representations on them, this number is larger than the corresponding number for the affine extension. This is desirable for the purpose of deducing an upper bound on the number of degrees of freedom in any spacetime. Next, recall that the eigenvalues of the area operator for the horizon, lying within one Planck area of the classical horizon area $`A_H`$, are given by $$\widehat{A}_H\mathrm{\Psi }_S=8\pi \beta l_P^2\underset{l=1}{\overset{p}{}}[j_l(j_l+1)]^{\frac{1}{2}}\mathrm{\Psi }_S,$$ (8) where, $`l_P`$ is the Planck length, $`j_l`$ is the spin on the $`l`$th puncture on the 2-sphere and $`\beta `$ is the Barbero-Immirzi parameter . We consider a large fixed classical area of the horizon, and ask what the largest value of number of punctures $`p`$ should be, so as to be consistent with (8); this is clearly obtained when the spin at each puncture assumes its lowest nontrivial value of 1/2, so that, the relevant number of punctures $`p_0`$ is given by $$p_0=\frac{A_H}{4l_P^2}\frac{\beta _0}{\beta },$$ (9) where, $`\beta _0=1/\pi \sqrt{3}`$. We are of course interested in the case of very large $`p_0`$. Now, with the spins at all punctures set to 1/2, the number of states for this set of punctures $`𝒫_0`$ is given by $$N^{𝒫_0}=\underset{m_1=1/2}{\overset{1/2}{}}\mathrm{}\underset{m_{p_0}=1/2}{\overset{1/2}{}}\left[\delta _{(_{n=1}^{p_0}m_n),0}\frac{1}{2}\delta _{(_{n=1}^{p_0}m_n),1}\frac{1}{2}\delta _{(_{n=1}^{p_0}m_n),1}\right]$$ (10) The summations can now be easily performed, with the result: $$N^{𝒫_0}=\left(\begin{array}{c}p_0\\ p_0/2\end{array}\right)\left(\begin{array}{c}p_0\\ (p_0/21)\end{array}\right)$$ (11) There is a simple intuitive way to understand the result embodied in (11). This formula simply counts the number of ways of making $`SU(2)`$ singlets from $`p_0`$ spin $`1/2`$ representations. The first term corresponds to the number of states with net $`J_3`$ quantum number $`m=0`$ constructed by placing $`m=\pm 1/2`$ on the punctures. However, this term by itself overcounts the number of SU(2) singlet states, because even non-singlet states (with net integral spin, for $`p`$ is an even integer) have a net $`m=0`$ sector. Beside having a sector with total $`m=0`$, states with net integer spin have, of course, a sector with overall $`m=\pm 1`$ as well. The second term basically eliminates these non-singlet states with $`m=0`$, by counting the number of states with net $`m=\pm 1`$ constructed from $`m=\pm 1/2`$ on the $`p_0`$ punctures. The difference then is the net number of $`SU(2)`$ singlet states that one is interested in for that particular set of punctures. To get to the entropy from the counting of the number of conformal blocks, we need to calculate $`N_{bh}=_𝒫N^𝒫`$, where, the sum is over all sets of punctures. Then, $`S_{bh}=lnN_{bh}`$. It may be pointed out that the first term in (11) also has another interpretation. It represents the counting of boundary states for an effective $`U(1)`$ Chern-Simons theory. It counts the number of ways unit positive and negative $`U(1)`$ charges can be placed on the punctures to yield a vanishing total charge. This would then correspond to an entropy bound given by the same formula (2) above for binomial distribution of charges. On the other hand the combination of both terms in (11), which corresponds to counting of states in the $`SU(2)`$ Chern-Simons theory, yields an even tighter bound for entropy than that in eq. (2). One can show that , the contribution to $`N_{bh}`$ for this set of punctures $`𝒫_0`$ with all spins set to 1/2, is by far the dominant contribution; contributions from other sets of punctures are far smaller in comparison. Thus, the entropy of an isolated horizon is given by the formula derived in ref. . We may mention that very recently Carlip has presented compelling arguments that this formula may possibly be of a universal character. Here, the formula follows readily from eq. (11) and Stirling approximation for factorials of large integers. The number of punctures $`p_0`$ is rewritten in terms of area $`A_H`$ through eq. (9) with the identification $`\beta =\beta _0ln2`$. This allows us to write the entropy of an isolated horizon in terms of a power series in horizon area $`A_H`$: $$S_{bh}=lnN^{𝒫_0}=\frac{A_H}{4l_p^2}\frac{3}{2}ln\left(\frac{A_H}{4l_p^2}\right)\frac{1}{2}ln\left(\frac{\pi }{8(ln2)^3}\right)O(A_H^1).$$ (12) Notice that the constant term here is negative and so is the order $`A_H^1`$ term. This then implies that the entropy is bound from above by a tighter bound which can be written in terms of Bekenstein-Hawking entropy ($`S_{BH}=A_H/4l_p^2)`$ as: $$S_{max}=ln\left(\frac{\mathrm{exp}S_{BH}}{S_{BH}^{3/2}}\right)$$ (13) Inclusion of other than spin $`1/2`$ representations on the punctures does not affect this bound. For example, we may place spin 1 on one or more punctures and spin $`1/2`$ on the rest. The number of ways singlets can be made from this set of representations can be computed in a straight forward way. Adding these new states to the already counted ones above, just changes the constant and order $`A_H^1`$ terms in formula (12). However, these additional terms continue to be negative and hence the entropy bound (13) still holds.<sup>§</sup><sup>§</sup>§Using the Cardy formula with the prefactor (á la Carlip ) appears to lead to entropy corrections for certain black holes not in accord with eq. (13) (although the bound (2) is indeed respected). This could be an artifact of the application of the Cardy formula. We refrain from further comment on these works since the precise relation of the Cardy formula approach to the present framework is not clear. The steps leading to the EB now follows the standard route of deriving the Bekenstein bound (see, e.g., ): we assume, for simplicity that the spatial slice of the boundary of an asymptotically flat spacetime has the topology of a 2-sphere on which is induced a spherically symmetric 2-metric. Let this spacetime contain an object whose entropy exceeds the bound. Certainly, such a spacetime cannot have an isolated horizon as a boundary, since then, its entropy would have been subject to the bound. But, in that case, its energy should be less than that of a black hole which has the 2-sphere as its (isolated) horizon. Let us now add energy to the system, so that it does transform adiabatically into a black hole with the said horizon, but without affecting the entropy of the exterior. But we have already seen above that a black hole with such a horizon must respect the bound; it follows that the starting assumption that the object, to begin with, had an entropy violating the bound is not tenable. There is, however, an important caveat in the foregoing argument. Strictly speaking, there is as yet no derivation of the second law of black hole mechanics within the framework of the isolated horizon. But, that is perhaps not a conceptual roadblock as far as deriving the EB is concerned. One has to assume that if matter or radiation crosses the isolated horizon adiabatically in small enough amounts, the isolated character of the horizon will not be seriously affected. This is perhaps not too drastic an assumption. Thus, for a large class of spacetimes, one may propose Eq.(13) as the new holographic entropy bound. Finally, we should mention that we prefer to think of the above holographic principle and the consequent entropy bound as ‘weak’ rather than ‘strong’ in the sense of Smolin . The work of SD was supported by NSF grant NSF-PHY-9514240 and by Eberley Research Funds of Penn State University. PM thanks J. Ambjorn, A. Ashtekar, A. Ghosh, H. Nicolai, S. Kalyana Rama, R. Loll and L. Smolin for illuminating discussions and the Center for Gravitational Physics and Geometry at Penn State University, the Niels Bohr Institute and the Albert Einstein Institute for their very kind hospitality during which this work was completed.
warning/0006/nlin0006021.html	ar5iv	text	# 1 Direct computation of 𝐾_{𝑐⁢𝑟} in standard map: average transition time through the destroyed critical curve vs. supercriticality. Circles show numerical results for 𝑁_{𝑡⁢𝑟}=400; stars represent 3 single–trajectory runs, including one with the minimal Δ⁢𝐾 (); straight line is relation () with parameters () fitted from 15 left–most points (circles). Critical perturbation in standard map: A better approximation Boris Chirikov<sup>1</sup><sup>1</sup>1Email: chirikov@inp.nsk.su Budker Institute of Nuclear Physics 630090 Novosibirsk, Russia ## Abstract Direct computation of the transition time between neighbor resonances in the standard map, as a function of the perturbation parameter $`K`$, allows for improving the accuracy of the critical perturbation value up to $`K_{cr}K_g<2.5\times 10^4`$ that is by a factor of about 50 as compared to the previous result due to MacKay and Percival. As is well known by now a typical structure of the phase space of a few–freedom nonlinear dynamical system is characterized by a very complicated admixture of both chaotic as well as regular (integrable) components of motion (the so–called divided phase space, see, e.g., ). Statistical properties of such a motion are very intricate and unusual. One of the most interesting (and important for many applications) problem is the conditions for transition from a local (restricted to relatively small regions in phase space) to the global chaos covering the whole available phase space. The most studied model of such a transition is described by the so–called (canonical) standard map (for history of this model see ): $$\overline{y}=y\frac{K}{2\pi }\mathrm{sin}(2\pi x),\overline{x}=x+\overline{y}$$ (1) where $`K`$ is the perturbation parameter. In this simple model the transition to global chaos corresponds to some exact critical value $`K=K_{cr}`$. For $`K>K_{cr}`$ the motion becomes infinite (in momentum $`y`$) for some initial conditions while for $`KK_{cr}`$ all the trajectories are confined within a period of map (1): $`\mathrm{\Delta }y=1`$. The first idea how to solve this difficult problem was due to Greene . First, he was able to solve a much simpler problem of the critical perturbation $`K(r)`$ at which a particular invariant Kolmogorov - Arnold - Moser (KAM) curve with the rotation number $`r`$ is destroyed. Critical function $`K(r)`$ is extremely singular with big dips at everywhere dense set of rational $`r`$ values (see, e.g., ). The physical mechanism of this behavior (known since Poincaré) is explained by resonances in the system (1) as the rotation number is the ratio of oscillation/perturbation frequencies. Whence, the main Greene’s idea: to find the ’most irrational’ $`r=r_g`$ which would correspond to the motion ’most far–off’ all the resonances. The former is well known in the number theory: $`r_g=[111\mathrm{}]=(\sqrt{5}1)/2`$ where the first representation is a continued fraction. This ’golden’ curve was found to be critical at the parameter $`K=K_g=0.97163540631\mathrm{}`$ . It was conjectured that for $`K>K_g`$ all invariant curves are destroyed , that is $`K_{cr}=K_g`$. The ’most–irrational’ assumption - as plausible as it is - remains a hypothesis. The main difficulty is here in that the resonance interaction and overlap, destroying invariant curves, depend not only on the resonance spacings, which are indeed maximal for $`r=r_g`$, but also on the amplitudes of those which are not simply an arithmetical property. Another argument, based on the analysis of the critical function $`K(r)`$ , also does not prove this principal hypothesis. A different approach to the problem - the so–called converse KAM theory - was developed in . It relies upon a rigorous criterion for the absence of any invariant curve in a certain region. Unfortunately, this criterion can only be checked numerically, and besides it provides the upper bound $`K_{cr}^+`$ only (the lower bound $`K_{cr}^{}=K_g`$). The remaining gap, or the accuracy of $`K_{cr}`$: $$(\mathrm{\Delta }K)_{cr}=K_{cr}^+K_{cr}$$ (2) can be made arbitrarily small at the expense of computation time $`t_C`$ which scales as $$t_C(K_{cr}^+K_{cr})^p$$ (3) Facing this difficulty, it is natural to recall the first method for calculating the critical perturbation used in . The method was based on the direct computation of trajectories for different $`KK_g`$. The criterion of supercriticality of a particular $`K`$ value was very simple: the transition if only a single trajectory in one of two neighbor integer resonances ($`y_r=0mod1`$) through the destroyed critical curve. With the computers available at that time the minimal $`K=1`$ has been reached only which corresponds to the uncertainty $`(\mathrm{\Delta }K)_{min}=K_{min}K_g=0.0284`$. This may be compared to the later result $`(\mathrm{\Delta }K)_{min}=0.0127`$ . Remarkably, the dependence of the average transition time on parameter $`K`$ was found to be similar to scaling (3): $$<t>=\frac{A}{(KK_{cr})^p}$$ (4) Fitting three unknown parameters gave: $`A=103,p=2.55`$, and $`K_{cr}=0.989`$. The latter result was rather different from the present value $`K_{cr}K_g`$, again because of the computation restrictions mentioned above: $`K1,t10^7`$ iterations. Nevertheless, the fitting Eq.(4) provided a less uncertainty $`(\mathrm{\Delta }K)_f=K_fK_g=0.0174`$ as compared to the result from the minimal $`K`$. The same is true for data from where $`(\mathrm{\Delta }K)_f=K_fK_g=0.00236`$. The latter value was apparently obtained by the direct fitting the relation (3). Fitting in log–log scale provides a much better result: $`(\mathrm{\Delta }K)_f=K_fK_g=0.000128\pm 0.000288`$ that is the remainig uncertainty reduces down to $`0.000288`$. In both cases the fitted value for the critical perturbation $`K_{cr}`$ is only true up to a certain confidence probability while the minimal $`K`$ is an exact result: $`K_{cr}^+=K_{min}`$. In the present paper the studies are continued with much better computers. The main result is farther considerable increasing of the accuracy $`(\mathrm{\Delta }K)_{cr}`$. To reduce the computation expenses, the transition time was calculated for a number of trajectories $`N_{tr}`$ started near the unstable fixed point of a half–integer resonance ($`y_r=1/2mod1`$), and then run until each of them crossed over to a neighbor integer resonance. The minimal $`K`$ value is determined already by the first trajectory escaped from the half–integer resonance. In this way the minimal uncertainty $$(\mathrm{\Delta }K)_{min}=K_{min}K_g=0.00025$$ (5) has been achieved with the escape time $`t6.77\times 10^{11}`$ itterations which took about 72 hours of CPU time on ALPHA–4100 computer (see Fig.1). The average transition time was computed from $`N_{tr}=400`$ trajectories for each of 100 values of $`K`$ in the interval: $`0.0035KK_g0.35`$. This costed 36 hours of computation. The results are shown in Fig.1. In the whole interval of $`\mathrm{\Delta }K`$ the dependence $`<t(K)>`$ is not exactly a power–law. It becomes so asymptotically for $`KK_{cr}`$ as expected from the theory . For this reason, only few smallest $`K`$ values of the function $`<t(K)>`$ were taken for the final fitting which is also shown in Fig.1 by the solid line. It is obtained from the fitting 15 left–most points (just up to the first big fluctuation) in log–log scale, and corresponds to the following parameters in Eq.(4): $$(\mathrm{\Delta }K)_f=\mathrm{\hspace{0.17em}0.000125}\pm 0.000267,p=\mathrm{\hspace{0.17em}2.959}\pm 0.0771,A=\mathrm{\hspace{0.17em}33}\pm 8$$ (6) The fitting relative accuracy $`rms=0.071`$ is close to, but somewhat larger than, the standard $`rms=1/\sqrt{N_{tr}}=0.05`$. This is seen from the data of 3 single trajectories in Fig.1, too. Notice also 2 very big deviations for the average over 400 trajectories which nature remains unclear. Interestingly, the relative fitting accuracy of the data is considerably higher: $`rms=0.02`$. This would require as many as about 5000 trajectories in the present method. However, it does not mean that the computation of the procedure in would be shorter. The most important parameter in (6) is $`(\mathrm{\Delta }K)_f`$ which is zero within statistical errors. This farther confirms the Greene hypothesis $`K_{cr}=K_g`$. The exponent $`p`$ is also equal to the theoretical value $`p_{th}=3.011722`$ to the fitting accuracy. The present value of parameter $`A`$ is much less than in because of a different (shorter) transition between resonances used. The summary of all results is presented in the Table below. \| Table. Accuracy of $`K_{cr}`$ in standard map \| \| \| \| --- \| --- \| --- \| \| $`\left(\mathrm{\Delta }K\right)_{min}`$ \| $`\left(\mathrm{\Delta }K\right)_{fit}`$ \| Reference \| \| exact \| probable \| \| \| $`2.84\times 10^2`$ \| $`1.74\times 10^2`$ \| \| \| $`1.27\times 10^2`$ \| $`3.36\times 10^3`$ \| \| \| \| $`\pm 1.\times 10^3`$ \| \| \| \| $`1.28\times 10^4`$ \| \| \| \| $`\pm 2.88\times 10^4`$ \| our fit \| \| $`2.5\times 10^4`$ \| $`1.25\times 10^4`$ \| present \| \| \| $`\pm 2.67\times 10^4`$ \| paper \| A serious difficulty in such a numerical approach to the problem is the computation accuracy. This was mentioned also in but no estimate for the computation errors was given, apparently because of a very complicated numerical procedure. Even in a much simpler method , accepted in the present study, the effect of noise turned out to be rather complicated. Special numerical experiments were done to clarify the question. To this end, a random perturbation of amplitude $`\nu `$ was introduced in both equations (1). The results are shown in Fig.2. Typically, the transition time becomes less than that without noise, and saturates below some critical noise–dependent value of $`K`$: $`\mathrm{\Delta }K\stackrel{<}{}\text{ }B(\nu )`$. However, in some cases the average transition time considerably grows, as an example in Fig.2 demonstrates, apparently due to a sharp increase of the fluctuations near the crossover from normal (noisefree) dependence of $`<t(\mathrm{\Delta }K)>`$ to the saturation. In turn, these fluctuations are apparently explained by the noise–induced diffusion into some of many small domains of regular motion within the critical structure. A rough estimate for unknown function $`B(\nu )`$ can be obtained as follows. The transition time is primarily determined by the width $`\delta y(\mathrm{\Delta }K)^2`$ of the chaotic layer around destroyed critical curve while the diffusion time through this layer $`t_01/\mathrm{\Delta }K`$ . Noise decreases this time down to $`t_\nu (\delta y)^2/\nu ^2`$. Hence, the crossover corresponds to $`t_\nu t_0`$, whence: $$B(\nu )A\nu ^b$$ (7) with $`b=2/5`$. Fitting the empirical data in Fig.2 in log–log scale gives: $`b0.39\pm 0.012`$, which is surprisingly close to the theoretical estimate, and $`A0.9716\pm 0.054`$ (Fig.3). The fitting accuracy is also fairly good: the relative $`rms=0.019`$. Moreover, below crossover ($`\mathrm{\Delta }K<B(\nu )`$) the width $`\delta y`$ as well as the diffusion time depend on $`\nu `$ only, and hence the transition time remains approximately constant for a given $`\nu `$ (Fig.2). In any event, the minimal $`(\mathrm{\Delta }K)_{min}`$ (5), which is the main result of the present study, is well above the expected limitation for the double–precision computation (see Fig.3). In conclusion, the direct approach a la to the problem of the critical perturbation in the standard map does further confirm Greene’s hypothesis $`K_{cr}=K_g`$ with a much better exact upper bound (5): $`K_{cr}K_g<\mathrm{\hspace{0.17em}2.5}\times 10^4`$. Still another recent confirmation of this conjecture (curiously, with roughly the same statistical accuracy (6)) has been inferred from a detailed study of the critical structure at the chaos–chaos border in standard map for $`K=K_g`$ . Acknowledgements. I am grateful to D.L. Shepelyansky for interesting discussions. This work was partially supported by the Russia Foundation for Fundamental Research, grant 97–01–00865.
warning/0006/astro-ph0006022.html	ar5iv	text	# The Space Density of low redshift AGN ## 1 Introduction The QSO optical luminosity function (OLF) and its evolution with redshift has been studied extensively for over three decades (see e.g. Schmidt 1968, Marshall et al. 1983, Boyle et al. 1988, Hewett, Foltz & Chaffee 1993, La Franca & Cristiani 1997). This has led to a detailed picture of the QSO OLF over a wide range in redshift from $`z0.3`$ to $`z>4`$. In contrast, the local ($`z<0.15`$) QSO OLF is actually much more poorly determined, frustrating attempts to link QSO evolution at moderate to high redshifts with nuclear activity in galaxies at the present epoch. This is due to a number of factors associated with the compilation of a suitable sample of local active galactic nuclei (AGN) with which to derive the local OLF. First, local AGN are relatively rare. Their space density is approximately 100 times less than that of normal galaxies, and large area surveys are required to yield a statistically useful sample. Secondly, many selection techniques for local AGN suffer from morphological biases. While surveys for stellar-like objects are clearly biased against resolved AGN, galaxy-based surveys are equally biased against objects with a dominant nuclear component. Finally, accurate knowledge of the nuclear OLF requires accurate subtraction of the light from the host galaxy. In the low luminosity AGN ($`M_B>23`$) that constitute the vast majority of the low redshift population, the light from the host galaxy may dominate the nuclear luminosity. Even at relatively low redshifts ($`z0.1`$), seeing limitations imposed by ground based observations limit accurate modelling of the luminosity profiles of the central regions of AGN host galaxies to scales typically larger than 1h$`{}_{}{}^{1}{}_{50}{}^{}`$kpc. A recent attempt to estimate the local AGN OLF has been carried out by Köhler et al. (1997), hereinafter K97. Using a sample of 27 candidates selected from the Hamburg/ESO objective prism survey, K97 derived a local AGN LF that exhibited a featureless power law form over a wide range in absolute magnitudes $`24<M_B<18`$. The form of the low redshift AGN LF is thus very different from the two-power-law luminosity function at higher redshifts ($`z0.5`$). This is a significant challenge for any theoretical model which seeks to connect the evolution of QSOs at high redshift with the local AGN population. The Hamburg/ESO survey covers an extensive area (611deg<sup>2</sup>; now extended to 3700deg<sup>2</sup>, see Wisotzki 2000) and is free of morphological bias. Unfortunately the spatial resolution (1–2 arcsec) of the survey is not sufficiently good to permit an accurate deconvolution of the galaxy and nuclear light even for the lowest redshift AGN ($`z<0.1`$) in the sample. For the 0.07 $`<z<`$ 0.3 sample K97 used small-aperture, zero-point corrected $`B`$ band CCD magnitudes which were subsequently corrected to reflect nuclear luminosities by subtracting a template host galaxy value of $`M_B=21`$ ; for the AGN with $`z<0.07`$ corrections were calculated individually and ranged from 0.21 to 1.61 mag. Until recently, AGN data sets studied with HST were either too small or the samples on which they were based were too heterogeneous to construct a reliable estimate of the local OLF. We report here on the estimate of the local nuclear OLF based on HST observations of 76 AGN selected from a unbiased sample of X-ray selected AGN. The sample is part of the extensive Einstein Medium Sensitivity Survey (EMSS, Stocke et al. 1991) which covers over 400deg<sup>2</sup> and has near-complete ($`>96`$ per cent) optical spectroscopic identification with no morphological bias. An earlier attempt to derive the low redshift OLF using the EMSS was made by della Ceca et al. (1996). They used 226 broad line AGN with $`z<0.3`$ to obtain a total (nuclear + host) OLF. This OLF was then convolved with the observed distribution of nuclear-to-total flux ratios for Seyfert 1 and 1.5 galaxies to yield a nuclear OLF. In this paper we propose to improve significantly on this work by using the HST observtions to correct explicity for the host galaxy light in each AGN. In section 2 we report on the measurement of the OLF from this sample, and in sections 3 and 4 we present and discuss the results obtained, comparing them to the K97 results. We present our conclusions in section 5. ## 2 Analysis ### 2.1 The data Details of the comprehensive HST imaging survey from which the sample of AGN used in this analysis were drawn is presented in a paper by Schade et al. (2000, hereinafter SBL). A full discussion of the methods used to select and observe these AGN is presented by SBL, thus only a brief description will be given here. HST observations of 76 $`z<0.15`$ AGN selected from the EMSS survey were carried out in the F814W ($`I`$) band, chosen to assist in the detection of the redder host galaxy components over the bluer nucleus. These data were complemented by deeper ground-based observations in the $`B`$ and $`R`$ bands for 69 AGN in the survey. A simultaneous three-component parametric model fit to the $`B`$, $`R`$ and $`I`$ images was performed for each AGN in the sample to derive magnitudes for the nuclear point source, bulge and disk components in each object. Despite the improved spatial resolution afforded with the HST, the fitting procedure is complex and error estimation required significant modelling of the fitting process. For host-dominated objects uncertainties in $`M_B(\mathrm{host})`$<sup>1</sup><sup>1</sup>1In the SBL study all magnitudes were based on the AB system. Nuclear $`B(\mathrm{AB})`$ magnitudes were derived by applying a mean $`(BI)_{(\mathrm{AB})}=0.2`$ colour correction to the nuclear $`I(\mathrm{AB})`$-band magnitudes obtained from the fit to the HST data. For objects of this colour, there is a negligible colour term between $`B`$ and $`B(\mathrm{AB})`$ passbands. For the purposes of this analysis we have therefore assumed $`M_B(\mathrm{AB})=M_B`$ for the nuclear regions. were typically $`\pm 0.25`$ magnitudes in the region of the $`M_B,z`$ plane where the AGN are found, but increased to $`\pm 0.5`$mag where the nucleus was dominant. Similarly, errors in nuclear magnitudes were $`\pm `$0.25 mag for bright nuclei but as much as $`\pm 0.5`$ mag for host-dominated objects. In total, nuclear $`M_B`$ magnitudes were obtained for 66 AGN in the sample (10 had no detectable nuclear component), and these data form the basis for the calculation of the OLF below. Nuclear absolute magnitudes were found to lie in the range $`14.6>M_B>24.1`$. The region of the AGN $`M_B(\mathrm{nuc}),z`$ plane sampled in this study is shown in Fig. 1. Since we are attempting to contruct an OLF from an X-ray-selected sample we need to ensure that there is a good correlation between nuclear optical and X-ray luminosity for objects in the SBL sample. We used the monochromatic $`2`$keV X-ray luminosity, $`L_{2\mathrm{k}\mathrm{e}\mathrm{V}}`$(nuc) and 2500Å UV fluxes, $`L_{2500\mathrm{A}}`$(nuc), listed in the SBL paper. $`L_{2500}`$Å(nuc) is based on the fitted nuclear $`M_B(\mathrm{nuc})`$, assuming a power-law optical/UV spectrum of the form $`f_\nu \nu ^{0.5}`$. Errors on the X-ray flux range from 5 per cent for the brightest X-ray sources to 25 per cent for the faintest X-ray sources (Gioia 1990). Although the X-ray luminosity for each source is expected to be dominated by the AGN, it is impossible to rule out some contribution from the host galaxy, particularly for the lowest luminosity sources. The least-squares fit to the observed relationship between $`L_{2\mathrm{k}\mathrm{e}\mathrm{V}}`$(nuc) and $`L_{2500\mathrm{A}}`$(nuc) plotted in Fig. 2 gives a relation of the form $`L_{2\mathrm{k}\mathrm{e}\mathrm{V}}(\mathrm{nuc})L_{2500\mathrm{A}}(\mathrm{nuc})^{0.82\pm 0.08}`$, consistent with other studies (e.g. Green et al. 1995). ### 2.2 The $`1/V_a`$ OLF estimate Space densities were derived using the $`1/V_a`$ estimator (Avni & Bahcall 1980). The OLF was constructed from the summed contributions of $`n`$ AGN using: $$\mathrm{\Phi }(M_B),z)=_{i=1}^n\frac{1}{V_a^i}\delta (M_B^iM_B)$$ $`V_a^i`$ being the accessible co-moving volume of the $`i^{\mathrm{th}}`$ AGN. The estimate of $`V_a`$ was based on $`z_{\mathrm{max}}`$ derived from the AGN’s X-ray flux and EMSS flux limits, assuming an X-ray spectral index $`\alpha _\mathrm{X}=1`$; $`f_\nu \nu ^{\alpha _\mathrm{X}}`$. To construct an optical OLF, we binned the $`1/V_a`$ estimates according to their optical nuclear $`M_B`$ magnitudes. We computed Poisson errors on the binned estimates of the OLF using $`\sigma =\left(\frac{1}{(V_a^i)^2}\right)^{0.5}`$. Not all 127 AGN with $`z<0.15`$ in the EMSS were observed by SBL and so a straightforward normalising factor of 0.6 (76/127) was applied to the area coverage function when computing the accessible volume. KS tests confirmed that the redshift and flux distribution for the SBL sample is consistent with the sample being drawn at random from the $`z<0.15`$ EMSS parent sample. The OLF was calculated at 1-mag intervals for the full redshift range $`z0.15`$. We made no correction for evolution across the redshift bin. We also constructed separate OLFs for AGN in elliptical and spiral hosts to investigate any host-related trends. ## 3 Results The differential OLF calculated for an Einstein-de Sitter universe in which $`H_0=50h_{50}`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_\mathrm{M}=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, is presented in Table 1 and plotted in Fig. 3, over plotted with data from K97 (their Table 5). As a comparison values using total galaxy luminosity (host + nucleus) are shown in Table 2 and plotted in Figure 5. The K97 data comprises 27 objects extending out to a redshift of 0.3. Of these, eight were at redshifts greater than the $`z=0.15`$ cut-off adopted in the SBL sample, all of which have $`M_B<24`$, i.e. brighter than the most luminous AGN in the SBL sample. In Fig. 3 we have also plotted two predictions of the $`z<0.15`$ OLF based on the luminosity evolution models of Boyle et al. (2000). These authors fit a variety of evolutionary models to a data set comprising over 6000 QSOs with $`M_B<23`$ and $`0.35<z<2.3`$ selected from the 2dF QSO redshift survey (Boyle et al. 1999) and the Large Bright QSO survey (LBQS, Hewett et al. 1995). Boyle et al. (2000) found that luminosity evolution models provided acceptable fits to the data, with exponential evolution ($`L^{}e^{k\tau }`$) as a function of look back time ($`\tau `$) favoured for a $`q_0=0.05`$ universe and as a general second order polynomial with redshift ($`L^{}10^{k_1z+k_2z^2}`$) for $`q_0=0.5`$. The extrapolated $`z<0.15`$ OLFs for the best-fitting ‘exponential’ and ‘polynomial’ models are shown as the short- and the long-dashed lines respectively in Fig. 3. The model OLFs have been plotted over the magnitude range consistent with the corresponding range (with respect to $`M_B^{}`$) over which they were derived at $`z>0.35`$. We obtained a reduced $`\chi ^2=1.0`$ for the exponential model fit to the SBL data at $`M_B<19`$, but were able to reject the extrapolation of the polynomial model at the 99 per cent confidence level. ## 4 DISCUSSION There is good agreement both in slope and normalisation between our estimate of the OLF and the K97 OLF at $`M_B<20`$. However, at fainter magnitudes the two estimates diverge. Our estimate of the OLF turns over to a much flatter slope whereas the K97 OLF continues to rise steeply. However, there are only three AGN in the K97 sample with $`M_B>20`$, whereas the SBL sample contains 37 AGN at these fainter magnitudes. It is therefore most probable that the difference between the two data-sets ($`2\sigma `$) at these magnitudes is simply due to small number statistics in the K97 sample. It is possible that the X-ray selection used to generate the SBL dataset is systematically biased against AGN with low optical luminosity. However, the correlation between $`L_{2\mathrm{k}\mathrm{e}\mathrm{V}}(\mathrm{nuc})`$ and $`L_{2500\mathrm{A}}(\mathrm{nuc})`$ plotted in Fig. 2 demonstrates that there is no systematic trend for objects with lower optical luminosities to exhibit relatively weaker X-ray-to-optical flux ratios. Indeed, the derived relation $`L_{2\mathrm{k}\mathrm{e}\mathrm{V}}(\mathrm{nuc})L_{2500\mathrm{A}}(\mathrm{nuc})^{0.82\pm 0.08}`$ implies the reverse, i.e. that AGN with lower optical luminosities have stronger X-ray-to-optical flux ratios. In common with other groups (La Franca & Cristiani 1997, Goldschmidt & Miller 1998), K97 have used their determination of the low redshift OLF to claim that the slope of the bright end of the OLF ($`\mathrm{\Phi }(L)`$) flattens significantly from $`\mathrm{\Phi }(L)L^{3.6}`$ at $`z>0.6`$ to $`\mathrm{\Phi }(L)L^{2.5}`$ at $`z<0.3`$. Such an observation would rule out pure luminosity evolution models, in which the shape of the OLF remains invariant with redshift. This is in marked constrast with our result that an extrapolation of the exponential form of a pure luminosity evolution model derived at $`z>0.35`$ still provides an adequate fit to the $`z<0.15`$ OLF. To investigate the discrepancy between these results, we fitted our binned estimate of the $`z<0.15`$ OLF with a two-power-law model of the form: $$\mathrm{\Phi }(L)L^\alpha L>L^{}$$ $$\mathrm{\Phi }(L)L^\beta L<L^{}$$ Fixing a ‘break’ luminosity at $`L^{}M_B^{}=20.5`$, we derived slopes of $`\alpha =2.1\pm 0.3`$ and $`\beta =1.1\pm 0.1`$ using a weighted least squares technique. This fit is over-plotted as the solid line on Fig. 3. This slope for $`\alpha `$ is indeed flatter than that derived at high redshift and is consistent with other estimates of the slope of the low redshift OLF, including the most recent determination of the bright end slope of the X-ray QSO LF ($`\alpha =2.6`$) by Miyaji et al. (1998). However, the value of $`\alpha `$ derived in this crude fashion is strongly dependent the choice of $`M_B^{}`$, and the inclusion of the two brightest bins that each contain a single object. By choosing a ‘break’ magnitude of $`M_B=21.5`$, and restricting consideration of the data points in the OLF to those bins which contain more than one object, we can increase this value to $`\alpha =2.6\pm 0.3`$. This is, admittedly, a very crude analysis and more sophisticated fitting of a smooth two-power-law function similar to that used to fit the higher redshift OLF would yield a more accurate estimate of the statistical errors associated with fitting the OLF. However, it is also likely that systematic errors play an equally important role in the determination of the local OLF. We attempted to estimate the sizes of such errors by first exploring the factor used to correct for the host galaxy luminosity. In the SBL sample the host galaxy luminosity was explicitly removed using fits to the individual HST images. For the bulk of their sample (i.e. 0.07 $`<z<0.3`$, or $`M_B<22`$) K97 relied on a two-step procedure using corrected CCD magnitudes measured in an aperture of diameter approximately equal to that of the seeing disk, hence replacing total magnitudes with small aperture magnitudes. Subsequently, a further correction factor for host galaxy luminosity was applied by adopting a template host galaxy of $`M_B=21`$. In Fig. 4 we have plotted nuclear luminosity against the total galaxy luminosity for the SBL sample. As also found by della Ceca et al. (1996), although the more powerful AGN reside in the more luminous hosts, the distribution of the ratio between nuclear and total luminosity is not constant; indeed the spread becomes very large at total absolute magnitudes fainter than $`M_B(\mathrm{nuc})=22.5`$. We re-computed our estimate of the $`z<0.15`$ OLF based on the SBL dataset using total, instead of nuclear, absolute magnitudes. The resulting OLF is shown in Fig. 5, with the both the original fit derived for the OLF and the extrapolated model fits plotted as a comparison. We find that although the bright end of the OLF has steepened appreciably to $`\alpha 3`$, both model fits are now clearly incompatible with the OLF computed in this fashion. The treatment of galaxy luminosities can thus result in significant differences to the estimate of the OLF at zero redshift. Note also that the evolutionary models have been derived from high redshift OLFs uncorrected for host galaxy light. Although the assumption that the host galaxy light makes an increasingly small contribution to the total luminosity of QSOs at high redshift may well be correct, some spectacular counter-examples have already been discovered (Aretxaga et al. 1995, Brotherton et al. 1999). We conclude that the large statisitical and systematic errors associated with the determination of the low redshift OLF and extrapolation of the OLF at higher redshifts make it difficult to rule out luminosity evolution models on the basis of shape of the low redshift OLF. We also attempted to derive nuclear OLFs for AGN with bulge-dominant (E/S0) and disk-dominant (spiral) hosts. Following SBL, the distinction between the two broad classes and elliptical was made on the basis of the bulge-to-total light ratio, $`B/T`$, for the host galaxy. Galaxies with $`B/T>0.5`$ were classified as E/S0 (44 in the sample), those with $`B/T0.5`$ as spiral (22 in the sample). The nuclear OLFs for different types of host galaxy is plotted in Fig. 6 (after re-normalising the E/SO LF to the same space density as the Sa/Sb LF). There are no significant differences between the two OLFs; confirmed by a KS test on the cumulative luminosity distributions for spiral and E/S0 hosts. This result is consistent with observation by SBL that $`B/T`$ for the host galaxy was independent of nuclear luminosity. In contrast, we note that McLure et al. (1999) found that the fraction of elliptical hosts increased significantly amongst the highest luminosity AGN. However, these observations are still based on relatively small datasets and the McLure et al. (1999) study predominantly samples a different luminosity regime ($`M_B<23`$) to that under investigation in this analysis. ## 5 Conclusions Results from a comprehensive, unbiased X-ray selected sample of AGN using the 0.1 arcsec resolving power of the HST, have enabled the first direct estimate of the nuclear OLF for AGN to be constructed. The OLF derived illustrates a two power law form similar to that derived for QSOs at higher redshifts and as such is different to the largely featureless power-law OLF claimed for the low redshift AGN identified in the Hamburg/ESO QSO survey. However, any discrepancy only occurs at $`M_B>20`$, where previous estimates of the OLF from the Hamburg/ESO survey are dominated by statistical errors arising from small number statistics. The OLF is consistent with an extrapolation of the exponential pure luminosity evolution derived at $`z>0.35`$ by Boyle et al. (2000), although the ’best-fit’ slope for the bright-end slope of the OLF is flatter than predicted by such pure luminosity evolution models. Given the large uncertainties associated with current estimates of the low redshift OLF (not least in the present analysis) and the extrapolation of evolutionary models to low redshift, it is almost certainly premature to rule out luminosity evolution on the basis of the current determinations of the low redshift OLF. Further detailed imaging work on optically-selected samples of low-moderate redshift AGN/QSOs will clearly help resolve this issue. With the superior imaging capability of the new generation of ground-based telescopes (Keck, Gemini) we may look forward to such data being obtained in the near future. ## 6 Acknowledgements We thank the referee, Lutz Wisotzki, for a number of useful suggestions which significantly improved the presentation and discussion of the results in this paper.
warning/0006/hep-th0006060.html	ar5iv	text	# Warped AdS near-horizon geometry of completely localized intersections of M5-branes ## 1 Introduction Intersecting branes are ubiquitous in string theory. Although much effort has been put into finding supergravity solutions of these systems, only solutions with branes localized in the overall transverse dimensions and those with at most one set of completely localized branes are known -. The cases where all branes are localized have so far eluded solution. In the past year some progress has been made in finding solutions of partially localized intersecting branes. Some of these results can be found in -. An interesting development has been to interpret brane delocalization physically . In this approach the delocalization seen in the supergravity solution is interpreted via the AdS/CFT correspondence of Maldacena as a Coleman-Mermin-Wagner theorem in the field theory. These results do not directly apply to the case we study and we will see that our near-horizon geometry describes completely localized branes. In this paper we report an exact solution of 11-dimensional supergravity for a system of intersecting M5-branes in the near-horizon limit. The particular system we study is the supergravity dual of $`𝒩=2`$ superconformal field theory with gauge group $`\mathrm{SU}(N)`$ and $`N_f=2N`$ fundamental flavors. This paper is a continuation of our work in which we solved the supersymmetry preservation conditions for the system. The full solution requires solving for a Kähler metric satisfying a non-linear partial differential equation in 7 variables! In our previous paper we solved this equation in an approximation where one set of branes were localized while the second set were smeared out over the worldvolume directions of the first set (these partially localized solutions were also found independently in .) This equation was studied in to yield an iterative expansion around the asymptotically flat region. This is the opposite limit to the one we pursue here. In the present paper we solve this differential equation exactly in the near horizon limit which is relevant to the AdS/CFT duality . The paper is structured as follows. We start with a brief description of the system under study. We then take a scaling limit where the Planck scale is taken to infinity while keeping field theory quantities fixed. Finally, we solve for the metric in this “near-horizon” limit, finding a warped AdS geometry. Warped AdS metrics have recently been discussed in the context of the AdS/CFT correspondence and semi-localized intersecting branes for brane configurations similar to ours. We conclude with some comments. ## 2 The system In this section we set up the problem and summarize some results from which we will need. One way of studying $`𝒩=2`$ gauge theories is to generalize the Hanany-Witten set-up to a system relevant to four dimensional gauge theories . The idea is to suspend D4-branes between a pair of NS5-branes which are separated by a finite coordinate distance $`L`$. The gauge theory living on the D4-branes will be, in the infrared, a four dimensional Yang-Mills theory with $`𝒩=2`$ supersymmetry and gauge group $`\mathrm{SU}(N)`$. There are many ways of introducing fundamental matter, but the easiest method is to introduce semi-infinite D4-branes on either side of the NS5-branes. The gauge D4-branes detect the semi-infinite D4-branes through strings which have ends on both types of D4-branes. These strings carry Chan-Paton factors with respect to the gauge groups of both types of D4-branes. From the gauge theory point of view these represent fundamental matter transforming in the fundamental representation of (a subgroup of) the flavor group. Witten pointed out that this system can be lifted to M-theory where this web of D4-branes and NS5-branes can be viewed as a single M5-brane wrapping a non-compact Riemann surface which coincides with the Seiberg-Witten Riemann surface. The same picture was derived in a different way in . In the remainder of this paper we will study a configuration of branes consisting of a set of coincident infinite D4-branes intersecting a pair of separated NS5-branes. This configuration can be viewed as one particular realization of $`\mathrm{SU}(N)`$ gauge theory with $`N_f=2N`$ according to the recipe described above. Our set-up can be arrived at from any generic Hanany-Witten configuration describing this field theory by moving the D6-branes (on which the semi-infinite D4-branes end ‘at infinity’) through the NS5-branes so that we have an equal number of semi-infinite D4-branes on both sides. We have also tuned the moduli so that all the D4-branes are coincident and collinear. In the gauge theory this corresponds to both tuning the bare masses of the fundamental matter to zero and sitting at the origin of the Coulomb branch where the gauge group is enhanced to the full $`\mathrm{SU}(N)`$. When viewed from the point of view of M-theory this looks simply like a system of intersecting M5-branes. It is convenient to pick a coordinate system such that the $`N`$ D4-branes have world-volume directions along $`x^0,x^1,x^2,x^3,x^6`$ while the NS5-branes have world-volume directions along $`x^0,x^1,x^2,x^3,x^4,x^5`$. The two sets of branes then intersect along $`x^0,x^1,x^2,x^3`$. The positions of the NS5-branes are $`x^6=\pm L/2`$. This configuration of branes can be lifted to M-theory with two sets of M5-branes intersecting along $`x^0,x^1,x^2,x^3`$. Let us denote one set as M5(1) branes, they have world-volume directions along $`x^0,x^1,x^2,x^3,x^6,x^7`$ and the other two M5-branes as M5(2), they have world volume directions along $`x^0,x^1,x^2,x^3,x^4,x^5`$. The M5(1) branes descend to D4-branes when $`x^7`$ is compactified to give type IIA string theory while the M5(2) branes are localized in the compactified $`x^7`$ direction and become NS5-branes. It is convenient to define a complex structure in the subspace $`x^4,x^5,x^6,x^7`$ as follows: $`vz^1`$ $`=`$ $`x^4+ix^5`$ (1) $`sz^2`$ $`=`$ $`x^6+ix^7.`$ (2) We also take $`x^7`$ to be a compact direction with radius $`R`$. In we solved the supersymmetry variation equations for M5-brane configurations which preserve at least 8 real supersymmetries. The general solution is given by the metric: $$ds^2=g^{\frac{1}{3}}dx_{3+1}^2+g^{\frac{1}{3}}g_{m\overline{n}}dz^mdz^{\overline{n}}+g^{\frac{2}{3}}\delta _{\alpha \beta }dx^\alpha dx^\beta ,$$ (3) and the 4-form field strength: $`F_{m\overline{n}\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{i}{4}}ϵ_{\alpha \beta \gamma }_\gamma g_{m\overline{n}}`$ (4) $`F_{m89(10)}`$ $`=`$ $`{\displaystyle \frac{i}{2}}_mg`$ (5) $`F_{\overline{m}89(10)}`$ $`=`$ $`{\displaystyle \frac{i}{2}}_{\overline{m}}g.`$ (6) The Greek indices run over the overall transverse coordinates $`x^8,x^9,x^{10}`$. Both the metric and 4-form are expressed in terms of the Kähler metric $`g_{m\overline{n}}`$. The source equations for the 4-form $`F`$ force $`g_{m\overline{n}}`$ to satisfy the non-linear partial differential equations: $$_\gamma _\gamma g_{m\overline{n}}+4_m_{\overline{n}}g=J_{m\overline{n}}$$ (7) where $`J`$ is the source specifying the positions of the M5-branes. The quantity $`g`$ appearing in the above equations is the square root of the determinant of the Kähler metric: $`g=g_{v\overline{v}}g_{s\overline{s}}g_{v\overline{s}}g_{s\overline{v}}`$. For the particular configuration that we will be studying the source equations are: $`^2g_{s\overline{s}}+4_s_{\overline{s}}g`$ $`=`$ $`8\pi ^3l_p^3\delta ^{(3)}(r)(\delta ^{(2)}(sL/2)+\delta ^{(2)}(s+L/2))`$ $`^2g_{v\overline{v}}+4_v_{\overline{v}}g`$ $`=`$ $`N8\pi ^3l_p^3\delta ^{(3)}(r)\delta ^{(2)}(v)`$ $`^2g_{v\overline{s}}+4_v_{\overline{s}}g`$ $`=`$ $`0`$ (8) where $`^2`$ is the flat Laplacian in the overall transverse space. To summarize, a given M5-brane configuration determines a source $`J`$ in (7). Solving this source equation for the metric $`g_{m\overline{n}}`$ then determines all other quantities in the supergravity solution. ## 3 Field theory (“near-horizon”) limit Maldacena proposed a certain scaling limit of string theory quantities to isolate the world-volume gauge theory from bulk interactions. The idea is simply to take a limit in which the Planck length goes to zero while keeping field theory quantities fixed. In we pointed out the relevant scalings of supergravity variables in type IIA theory. We can express these scalings in M-theory units by defining $`w,t`$ and $`y`$ as follows: $`w`$ $`=`$ $`{\displaystyle \frac{v}{\alpha ^{}}}={\displaystyle \frac{vR}{l_p^3}}`$ $`t^2`$ $`=`$ $`{\displaystyle \frac{r}{g_s\alpha ^{\frac{3}{2}}}}={\displaystyle \frac{r}{l_p^3}}`$ (9) $`y`$ $`=`$ $`{\displaystyle \frac{s}{R}}.`$ Note that $`w`$ and $`y`$ are complex variables while $`t`$ is a real variable. The field theory limit is one in which we keep $`w,y,t`$ fixed while taking $`l_p`$ to zero. We take the metric to be: $$\frac{1}{l_p^2}ds^2=g^{\frac{1}{3}}\eta _{\mu \nu }dx^\mu dx^\nu +g^{\frac{1}{3}}g_{m\overline{n}}dz^mdz^{\overline{n}}+g^{\frac{2}{3}}(4t^2dt^2+t^4d\mathrm{\Omega }_2^2)$$ (10) Where now $`m,n`$ run over $`y,w`$, $`g=g_{w\overline{w}}g_{y\overline{y}}g_{w\overline{y}}g_{y\overline{w}}`$ and $`d\mathrm{\Omega }_2^2`$ is the metric on the round unit 2-sphere. The source equations become: $`{\displaystyle \frac{1}{4t^5}}_t(t^3_t)g_{y\overline{y}}+4_y_{\overline{y}}g`$ $`=`$ $`\pi ^2{\displaystyle \frac{\delta (t)}{t^5}}\left(\delta ^{(2)}(y{\displaystyle \frac{1}{2g_{YM}^2}})+\delta ^{(2)}(y+{\displaystyle \frac{1}{2g_{YM}^2}})\right)`$ $`{\displaystyle \frac{1}{4t^5}}_t(t^3_t)g_{w\overline{w}}+4_w_{\overline{w}}g`$ $`=`$ $`N\pi ^2{\displaystyle \frac{\delta (t)}{t^5}}\delta ^{(2)}(w)`$ $`{\displaystyle \frac{1}{4t^5}}_t(t^3_t)g_{w\overline{y}}+4_w_{\overline{y}}g`$ $`=`$ $`0`$ (11) where $`g_{YM}^2=R/L`$ is the Yang-Mills coupling constant in the field theory. We have also assumed that the metric depends only on the “radial” coordinate $`t`$ and not the “angular” variables in the overall transverse directions. This is consistent with the requirement of having an $`\mathrm{SU}(2)`$ isometry corresponding to the field theory R-symmetry. We end this section with some comments. The source equations now have no powers of the Planck length $`l_p`$, they are expressed in terms of quantities which have a field theory interpretation. This will allow us to express the metric in terms of only these rescaled variables with no further dependence on $`l_p`$ aside from the overall multiplicative factor. Secondly, we would like to point out that while the initial set up treated the M5(1) and M5(2) branes on an equal footing, the scaling limit we take breaks that symmetry. This is clear from the type IIA picture since there the D4-branes play a distinguished role in that the field theory of interest lives on their world-volume. ## 4 Near-horizon geometry of intersecting M5-branes Since we are looking for a supergravity dual of a four-dimensional conformal field theory, we expect a solution where the metric contains an $`AdS_5`$ factor. The most general metric of this type is: $$\frac{1}{l_P^2}ds^2=\mathrm{\Omega }^2\left(u^2dx_{3+1}^2+\frac{1}{u^2}du^2\right)+ds_6^2$$ (12) where the metric for the six-dimensional transverse space and $`\mathrm{\Omega }`$ are independent of the $`AdS_5`$ coordinates. It is convenient to choose variables where there is only one dimensionful variable. We choose to do this by defining $`\rho `$ with dimensions of mass and dimensionless angular variables $`\theta `$ and $`\varphi `$ by: $`t`$ $`=`$ $`\rho \mathrm{cos}\theta `$ (13) $`w`$ $`=`$ $`\rho \mathrm{sin}\theta e^{i\varphi }`$ (14) Now we see on dimensional grounds that $`u`$ must be related to $`\rho `$ by $`u=\rho \alpha `$ where $`\alpha `$ is some function of the dimensionless variables $`\theta `$, $`\varphi `$, $`y`$ and $`\overline{y}`$. By substituting this expression for $`u`$ into the above metric, we can compare the metric components with those in the known form of the solution, eq. (10). In particular, by examining the factor multiplying $`dx_{3+1}^2`$ and the metric components $`g_{\rho \rho }`$, $`g_{\rho y}`$, $`g_{\rho \theta }`$ and $`g_{\rho \varphi }`$ we find: $`g`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Omega }^6\alpha ^6\rho ^6}}`$ $`g_{w\overline{w}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }^6\alpha ^44\mathrm{cos}^4\theta }{\rho ^4\mathrm{\Omega }^6\alpha ^6\mathrm{sin}^2\theta }}`$ $`g_{y\overline{w}}`$ $`=`$ $`{\displaystyle \frac{2e^{i\varphi }_y\alpha }{\rho ^3\alpha ^3\mathrm{sin}\theta }}`$ (15) $`_\theta \alpha `$ $`=`$ $`{\displaystyle \frac{(\mathrm{\Omega }^6\alpha ^44\mathrm{cos}^2\theta )\mathrm{cos}\theta }{\mathrm{\Omega }^6\alpha ^3\mathrm{sin}\theta }}`$ $`_\varphi \alpha `$ $`=`$ $`0`$ Since we are looking at $`𝒩=2`$ superconformal field theories we would like to preserve a $`\mathrm{SU}(2)\times \mathrm{U}(1)`$ isometry. This we have incorporated in the above ansatz by requiring that the metric preserve a $`\mathrm{U}(1)`$ which rotates $`w`$ by a phase and the $`\mathrm{SU}(2)`$ symmetry of the transverse 2-sphere. In fact, these symmetries are consequences of our required form of the metric. For example we see that $`\alpha `$ is independent of $`\varphi `$ and the equation for $`_\theta \alpha `$ shows that $`\mathrm{\Omega }`$ must also be independent of $`\varphi `$. If we assume, for the moment, that $`\mathrm{\Omega }`$ is constant<sup>3</sup><sup>3</sup>3As we will see below the metric has a warped AdS structure and so our assumption that $`\mathrm{\Omega }`$ is constant is incorrect. Nevertheless, for the purposes of solving the equations we find it easier to begin with an incorrect assumption which can be easily modified to yield the correct metric than to solve the equations directly., say $`\mathrm{\Omega }_0`$, we can solve for the $`\theta `$-dependence of $`\alpha `$ (which we denote by $`\alpha _0`$) in terms of an arbitrary function $`A(y,\overline{y})`$: $$\mathrm{\Omega }_0^6\alpha _0^4=4\mathrm{cos}^4\theta +4A(y,\overline{y})\mathrm{sin}^4\theta $$ (16) We can then write the above equations as: $`g`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_0^3}{8\rho ^6\beta ^3}}`$ (17) $`g_{w\overline{w}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_0^3A\mathrm{sin}^2\theta }{2\rho ^4\beta ^3}}`$ (18) $`g_{y\overline{w}}`$ $`=`$ $`{\displaystyle \frac{e^{i\varphi }\mathrm{\Omega }_0^3\mathrm{sin}^3\theta _yA}{4\rho ^3\beta ^3}}`$ (19) where we have defined: $$\beta =\left(\mathrm{cos}^4\theta +A\mathrm{sin}^4\theta \right)^{\frac{1}{2}}$$ (20) The metric component $`g_{y\overline{y}}`$ can be determined from the determinant $`g`$ and the other components of the metric given in the above equations. However, the metric determined in this way fails to be Kähler <sup>4</sup><sup>4</sup>4Hence the assumption of constant $`\mathrm{\Omega }`$ is incorrect.. It is easier instead to determine $`g_{y\overline{y}}`$ by requiring the metric to be Kähler. The Kähler condition is satisfied if: $$g_{y\overline{y}}=\frac{\mathrm{\Omega }_0^3\mathrm{sin}^4\theta \|_yA\|^2}{8\rho ^2A\beta ^3}$$ (21) provided that $`A=\|F(y)\|^2`$, where $`F`$ is a holomorphic function of $`y`$. One can easily check that the source equations are satisfied everywhere away from the support of the delta functions<sup>5</sup><sup>5</sup>5The normalizations and precise form of $`F`$ relevant to the delta function sources are determined below.. The metric as it stands now has a vanishing determinant so it is not a valid solution. However, it is easy to see that the metric can be modified in such a way as to get the correct determinant while continuing to satisfy the source equations. The idea is simply to add to all the metric components additional terms which are themselves Kähler (so as not to destroy the Kähler properties of our initial ansatz), such that the determinant is correctly reproduced. The source equations will continue to be satisfied provided these additional terms do not depend on $`t`$. From these simple requirements one determines the solution: $`g_{w\overline{w}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_0^3A\mathrm{sin}^2\theta }{2\rho ^4\beta ^3}}+{\displaystyle \frac{\|f\|^2}{\rho ^4\mathrm{sin}^4\theta }}`$ $`g_{y\overline{y}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_0^3\mathrm{sin}^4\theta \|_yA\|^2}{8\rho ^2A\beta ^3}}+{\displaystyle \frac{\|_yf\|^2}{\rho ^2\mathrm{sin}^2\theta }}`$ (22) $`g_{y\overline{w}}`$ $`=`$ $`{\displaystyle \frac{e^{i\varphi }\mathrm{\Omega }_0^3\mathrm{sin}^3\theta _yA}{4\rho ^3\beta ^3}}{\displaystyle \frac{e^{i\varphi }\overline{f}_yf}{\rho ^3\mathrm{sin}^3\theta }}`$ Where $`f(y)`$ is a holomorphic function determined from the requirement that the determinant has the form (17). This requirement can be stated succinctly as a differential equation: $$_y(f^2F)=f$$ (23) The general solution of this equation is: $$f(y)=\frac{_a^yF^{\frac{1}{2}}(z)𝑑z}{2F^{\frac{1}{2}}(y)}$$ (24) where $`a`$ is an arbitrary constant of integration. Note that the additional terms in the metric (involving $`f`$) are independent of $`t`$ ($`r`$ in the original coordinates) while the other terms in $`g_{m\overline{n}}`$ can be expressed in the form $`g_m_{\overline{n}}(\|Fw^2\|^2)`$. So the source equations can be conveniently written as: $$\frac{1}{4t^5}_t(t^3_tg)_m_{\overline{n}}(\|Fw^2\|^2)+4_m_{\overline{n}}g=J_{m\overline{n}}$$ (25) It is now a straightforward calculation to check that these equations are satisfied with the source terms: $`J_{w\overline{w}}`$ $`=`$ $`\pi ^2N{\displaystyle \frac{\delta (t)}{t^5}}\delta ^{(2)}(w)`$ (26) $`J_{y\overline{y}}`$ $`=`$ $`\pi ^2{\displaystyle \frac{\delta (t)}{t^5}}\left(\delta ^{(2)}(y{\displaystyle \frac{1}{2g_{YM}^2}})+\delta ^{(2)}(y+{\displaystyle \frac{1}{2g_{YM}^2}})\right)`$ (27) $`J_{y\overline{w}}`$ $`=`$ $`0.`$ (28) The source equations determine $`F`$ (and consequently also $`f`$) since the M5(2) branes are localized at the zeroes of $`F`$, as well as fixing the constant: $$\mathrm{\Omega }_0^3=4\pi N.$$ (29) We will now give the explicit form of the metric, exhibiting the warped product structure before solving for $`F`$. We will then consider the form of $`F`$ in the large radius limit before solving it in the general case. ### 4.1 Warped anti-de Sitter structure of the metric According to Maldacena’s conjecture conformal field theories have anti-de Sitter supergravity duals. Our metric is not a product manifold of anti-de Sitter space with a transverse manifold, but as mentioned earlier in this section, the metric can be written as a warped product consistent with Maldacena’s conjecture. To see this warped product structure, one simply returns to equations (15) and solves for $`\alpha `$ and $`\mathrm{\Omega }`$ using the explicit metric appearing in (22). This yields: $`\alpha ^2`$ $`=`$ $`{\displaystyle \frac{2\pi N}{(\mathrm{cos}^4\theta +\|F\|^2\mathrm{sin}^4\theta )^{\frac{1}{2}}}}+{\displaystyle \frac{\|f\|^2}{\mathrm{sin}^2\theta }}`$ $`(\mathrm{\Omega }\alpha )^6`$ $`=`$ $`{\displaystyle \frac{\pi N}{2}}{\displaystyle \frac{1}{(\mathrm{cos}^4\theta +\|F\|^2\mathrm{sin}^4\theta )^{\frac{3}{2}}}}.`$ (30) These expressions are consistent with all the metric components derived and it can easily be checked that $`_\theta \alpha `$ has the correct form required by eq. (15). The metric, therefore, can be written as a warped product of AdS space with a transverse manifold. The metric, while messy, can be written relatively concisely if one expresses it in terms of $`\alpha `$ and $`\mathrm{\Omega }`$: $`{\displaystyle \frac{1}{l_P^2}}ds^2`$ $`=`$ $`\mathrm{\Omega }^2(u^2\eta _{\mu \nu }dx^\mu dx^\nu +{\displaystyle \frac{du^2}{u^2}})+{\displaystyle \frac{4\mathrm{cos}^2\theta }{\mathrm{\Omega }^4\alpha ^4\mathrm{sin}^2\theta }}(1{\displaystyle \frac{4\mathrm{cos}^4\theta }{\mathrm{\Omega }^6\alpha ^4}})d\theta ^2`$ (31) $`+`$ $`{\displaystyle \frac{8\mathrm{cos}^3\theta }{\mathrm{sin}\theta \mathrm{\Omega }^4\alpha ^5}}_y\alpha d\theta dy+{\displaystyle \frac{8\mathrm{cos}^3\theta }{\mathrm{sin}\theta \mathrm{\Omega }^4\alpha ^5}}_{\overline{y}}\alpha d\theta d\overline{y}`$ $`+`$ $`\mathrm{\Omega }^2(1{\displaystyle \frac{4\mathrm{cos}^4\theta }{\mathrm{\Omega }^6\alpha ^4}})d\varphi ^22i\mathrm{\Omega }^2{\displaystyle \frac{_y\alpha }{\alpha }}d\varphi dy+2i\mathrm{\Omega }^2{\displaystyle \frac{_{\overline{y}}\alpha }{\alpha }}d\varphi d\overline{y}`$ $`+`$ $`{\displaystyle \frac{\mathrm{\Omega }^2\alpha ^2}{\mathrm{\Omega }^6\alpha ^44\mathrm{cos}^4\theta }}(\mathrm{sin}^2\theta +(2\mathrm{\Omega }^6\alpha ^4+8\mathrm{cos}^4\theta )\|{\displaystyle \frac{_{\overline{y}}\alpha }{\alpha ^2}}\|^2)dy\|^2`$ $``$ $`{\displaystyle \frac{\mathrm{\Omega }^2}{\alpha ^2}}(_y\alpha )^2dy^2{\displaystyle \frac{\mathrm{\Omega }^2}{\alpha ^2}}(_{\overline{y}}\alpha )^2d\overline{y}^2+{\displaystyle \frac{\mathrm{cos}^4\theta }{\mathrm{\Omega }^4\alpha ^4}}d\mathrm{\Omega }_2^2`$ Everything is now determined explicitly in terms of $`F(y)`$. We will now consider the form of $`F`$ in various cases, including the simple generalisation to conformal theories with gauge group $`\mathrm{SU}(N)^n`$. ### 4.2 Large $`R`$ or M-theory limit In the limit that $`R`$, the radius of $`x^7`$, becomes infinite we can ignore the periodicity of $`y`$: $`yy+i2\pi `$. Notice that the field theory is not sensitive to the value of $`R`$ but only to the ratio $`R/L`$ which determines the gauge coupling constant. We are thus simultaneously taking $`L`$ to infinity while keeping $`R/L`$ fixed. In this limit we can solve for $`F`$ taking into account the normalization of the sources: $$F=\left((y\frac{1}{2g_{YM}^2})(y+\frac{1}{2g_{YM}^2})\right)^{\frac{2}{N}}.$$ (32) Using this explicit expression we can calculate $`f`$: $$f=\frac{1}{2}(4g_{YM}^4)^{\frac{1}{N}}y(y^2\frac{1}{4g_{YM}^4})^{\frac{1}{N}}(\frac{1}{N},\frac{1}{2};\frac{3}{2};4g_{YM}^4y^2),$$ (33) where $``$ denotes the hypergeometric function. It is easy to see how one can generalize this to an arbitrary number of M5(2) branes. If there are $`n`$ M5(2) branes located at $`y=y_i`$ then: $$F=\underset{i=1}{\overset{n}{}}(yy_i)^{2/N}.$$ (34) We can then determine, at least in principle, $`f`$ from this expression. The dual conformal field theory will have a product gauge group $`\mathrm{SU}(N)^{n1}`$ with the gauge coupling of the $`i`$’th factor being given by: $$\frac{1}{g_{YM,i}^2}=y_{i+1}y_i$$ (35) ### 4.3 Solution for arbitrary $`R`$ As noted above the zeroes of $`F`$ determine the locations of the M5(2) branes. From the previous section it is easy to see how to generalize to an arbitrary radius of $`x^7`$ (i.e. when we take into account the periodicity of $`y`$). For our sources with periodic $`y`$ the correct $`F`$ is: $$F=\left(\mathrm{sinh}(y\frac{1}{2g_{YM}^2})\mathrm{sinh}(y+\frac{1}{2g_{YM}^2})\right)^{2/N}.$$ (36) In this case we have not been able to express $`f`$ in terms of a known function but it is still given by the integral in (24). One can similarly generalize this for a collection of $`n`$ M5(2) branes: $$F=\underset{i=1}{\overset{n}{}}\mathrm{sinh}(yy_i)^{2/N}.$$ (37) This again determines $`f`$ in principle through eq. (24). ## 5 Conclusions and discussion In this paper we presented an exact solution of 11-dimensional supergravity describing localized intersections of M5-branes. The solution has some surprising features worth pointing out. The geometry of our intersecting brane configuration is a warped anti-de Sitter product geometry, consistent with the fact that the dual quantum field theory is a conformal field theory. Another feature concerns the ‘t Hooft coupling. Taking the large $`N`$ limit to remain in the domain of validity of supergravity does not imply anything about the value of $`g_{YM}^2=R/L`$. This ratio is an arbitrary constant. In the solutions known thus far for 4-dimensional field theories the relevant combination appearing in the supergravity solution is always $`g_{YM}^2N`$, forcing the ‘t Hooft coupling to be large in the small curvature limit relevant to supergravity. In our case there appears to be no such restriction on the ‘t Hooft coupling. It is thus surprising that in principle one can tune the ‘t Hooft coupling to be small or large while remaining in the domain of validity of supergravity. However, the large $`N`$ limit may be rather subtle in this case and this issue is currently under investigation. Our solution does not have any simple $`N`$ dependence: there are terms of different orders in $`N`$ in a $`1/N`$ expansion despite the fact that we have taken the decoupling limit. Unlike the $`AdS_5\times S^5`$ case, the $`1/N`$ suppressed terms do not come with powers of the Planck scale. Certainly further terms relevant to the asymptotically flat solution will contain the Planck scale, however, it is surprising that the $`1/N`$ corrections do not appear to be directly connected to an expansion in the Planck scale. There are a number of directions which open up from this analysis. One is to consider other intersecting branes which are connected to this configuration through compactification and T-duality. The present system of M5-branes can be viewed as a special case of a more general problem of an M5-brane wrapped on a Riemann surface. The supergravity description of the general problem will be of interest for finding supergravity duals for more interesting field theories, including non-conformal field theories. The solution to the latter problem will be presented in . ## 6 Acknowledgements AF would like to thank Subir Mukhopadhyay for discussions, and the Durham University Department of Mathematical Sciences, where part of this work was done, for their hospitality. AF is supported by a grant from the Swedish Research Council, he would also like to acknowledge support from the NSF under grant PHY99-73935 during the academic year ‘98-’99 when this project was initiated.
warning/0006/quant-ph0006004.html	ar5iv	text	# Fast parallel circuits for the quantum Fourier transform ## 1 Introduction and summary of results In this paper we consider the quantum circuit complexity of the quantum Fourier transform (QFT). The quantum Fourier transform is the key quantum operation at the heart of Shor’s quantum algorithms for factoring and computing discrete logarithms and the known extensions and variants of these algorithms (see, e.g., Kitaev , Boneh and Lipton , Grigoriev , and Cleve, Ekert, Macchiavello, and Mosca ). The quantum Fourier transform also plays a key role in extensions of Grover’s quantum searching technique , due to Brassard, Høyer, and Tapp and Mosca . In order to discuss the quantum Fourier transform in greater detail we recall the discrete Fourier transform (DFT); for a given dimension $`m`$ the discrete Fourier transform is a linear operator on $`^m`$ mapping $`(a_0,a_1,\mathrm{},a_{m1})`$ to $`(b_0,b_1,\mathrm{},b_{m1})`$, where $$b_x=\underset{y=0}{\overset{m1}{}}(e^{2\pi i/m})^{xy}a_y.$$ (1) The discrete Fourier transform has many important applications in classical computing, essentially due to the efficiency of the fast Fourier transform (FFT), which is an algorithm that computes the DFT with $`O(m\mathrm{log}m)`$ arithmetic operations, as opposed to the obvious $`O(m^2)`$ method. The FFT algorithm was proposed by Cooley and Tukey in 1965 , though its origins can be traced back to Gauss in 1866 . The FFT plays an important role in digital signal processing, and it has been suggested as a contender for the second most important nontrivial algorithm in practice, after fast sorting. The DFT (and the FFT algorithm) generalize to certain algebraic structures, such as rings containing primitive $`m^{\text{th}}`$ roots of unity (which can play the role of $`e^{2\pi i/m}`$ in Eq. 1). This more abstract type of FFT is a principal component in Schönhage and Strassen’s fast multiplication algorithm , which can be expressed as circuits of size $`O(n\mathrm{log}n\mathrm{log}\mathrm{log}n)`$ for multiplying $`n`$-bit integers. For more applications—of which there are many—and historical information, see . The quantum Fourier transform (QFT) is a unitary operation that essentially performs the DFT on the amplitude vector of a quantum state—the QFT maps the quantum state $`_{x=0}^{m1}\alpha _x\|x`$ to the state $`_{x=0}^{m1}\beta _x\|x`$, where $$\beta _x=\frac{1}{\sqrt{m}}\underset{y=0}{\overset{m1}{}}(e^{2\pi i/m})^{xy}\alpha _y.$$ (2) For certain values of $`m`$ there are very efficient quantum algorithms for the QFT. The fact that the quantum circuit size can be polynomial in $`\mathrm{log}m`$ for some values of $`m`$ was first observed by Shor and is of critical importance in his polynomial-time algorithms for prime factorization and discrete logarithms. Shor’s original method may be described as a “mixed-radix” method, and is discussed further in Section 7.2. In the particular case where $`m=2^n`$, there exist quantum circuits performing the quantum Fourier transform with $`O(n^2)`$ gates, which was proved by Coppersmith (see also ). These circuits are based on a recursive description of the QFT that is analogous to the description of the DFT exploited by the FFT. While in some sense these quantum circuits are exponentially faster than the classical FFT, the task that they perform is quite different. The QFT does not explicitly produce any of the values $`\beta _0,\beta _1,\mathrm{},\beta _{m1}`$ as output (nor does it explicitly obtain any of the values $`\alpha _0,\alpha _1,\mathrm{},\alpha _{m1}`$ as input). Intuitively, the difference between performing a DFT and a QFT can be thought of as being analogous to the difference between computing all the probabilities that comprise a probability distribution and sampling a probability distribution—the latter task being frequently much easier. Coppersmith also proposed quantum circuits that approximate the QFT with error bounded by $`\epsilon `$, and showed that such approximations can be computed by circuits of size $`O(n\mathrm{log}(n/\epsilon ))`$ for modulus $`2^n`$. Such approximations can be thought of as unitary operations whose distance from the QFT (in the operator norm induced by Euclidean distance) is bounded by $`\epsilon `$. Kitaev showed how the QFT for an arbitrary modulus $`m`$ can be approximated by circuits with size polynomial in $`\mathrm{log}(m/\epsilon )`$. For most information processing purposes, it suffices to use such approximations of quantum operations (for $`\epsilon `$ ranging from constant down to $`1/n^{O(1)}`$). Indeed, since it seems rather implausible to physically implement quantum gates with perfect accuracy, the need to ultimately consider approximations is likely inevitable. Thus, we believe that the most relevant consideration is to approximately compute the QFT, though exact computations of the QFT are still of interest as part of the mathematical theory of quantum computation. Moore and Nilsson showed how to obtain logarithmic-depth circuits that perform encoding and decoding for standard quantum error-correcting codes. For the QFT, in both the exact and approximate case, the gates in Coppersmith’s circuits can be arranged so as to have depth $`2n1`$, as noted in , but not less depth than this. Similarly, the techniques of Shor and of Kitaev have polynomial depth. Our first result shows that it is possible to compute good approximations of the QFT with logarithmic-depth quantum circuits. ###### Theorem 1 For any $`n`$ and $`\epsilon `$ there is a quantum circuit approximating the QFT modulo $`2^n`$ with precision $`\epsilon `$ that has size $`O(n\mathrm{log}(n/\epsilon ))`$ and depth $`O(\mathrm{log}n+\mathrm{log}\mathrm{log}(1/\epsilon ))`$. By an approximation of a unitary operation $`U`$ with precision $`\epsilon `$, we mean a unitary operation $`V`$ (possibly acting on additional ancilla qubits) with the following property. For any input (pure) quantum state, the Euclidean distance between applying $`U`$ to the state and $`V`$ to the state is at most $`\epsilon `$ (in the Hilbert space that includes the input/output qubits and the ancilla qubits). Also, whenever we refer to circuits, unless otherwise stated, there is an implicit assumption that the circuits belong to a logarithmic-space uniformly generated family in the usual way (via a classical Turing machine). In Section 7.2, we consider a different approach for parallelizing Shor’s QFT method, which gives somewhat worse bounds. The proof of Theorem 1 follows the general approach introduced by Kitaev , with several efficiency improvements as well as parallelizations. In particular, we introduce a new parallel method for performing multiprecision phase estimation. We also show that, if size rather than depth is the primary consideration, it is possible to compute the QFT exactly with a near-linear number of gates. ###### Theorem 2 For any $`n`$ there is a quantum circuit that exactly computes the QFT modulo $`2^n`$ that has size $`O(n(\mathrm{log}n)^2\mathrm{log}\mathrm{log}n)`$ and depth $`O(n)`$. Theorem 2 is based on a nonstandard recursive description of the QFT combined with an asymptotically fast multiplication algorithm . There are several reasons why we believe results regarding quantum circuit complexity, such as in the above theorems, are important. First, circuit depth is likely to be particularly relevant in the quantum setting for physical reasons. Perhaps most notably, fault tolerant quantum computation necessarily requires parallelization anyway —under various noise models, error correction must continually be applied in parallel to the qubits of a quantum computer, even when the qubits are doing nothing. In such models, parallelization saves not only the total amount of time, but also the total amount of work. Furthermore, informally speaking, the depth of a quantum circuit corresponds to the amount of time coherence must be preserved, so in addition to saving work, parallelization may allow for larger quantum circuits to be implemented within systems having shorter decoherence times or using less extensive error correction. A final reason is that such results suggest alternate methods for performing various operations, which may in turn suggest or shed light on quantum algorithms for other problems or more general methods for improving efficiency of quantum algorithms. It has long been known that the main bottleneck of the quantum portion of Shor’s factoring algorithm is not the QFT, but rather is the modular exponentiation step. If it were possible to perform modular exponentiation by classical circuits with poly-logarithmic depth and polynomial size then it would be possible to implement Shor’s factoring algorithm in poly-logarithmic depth with a polynomial number of qubits. Although no such algorithm is known for modular exponentiation, we can prove the following weaker result, which nevertheless implies that quantum computers need only run for poly-logarithmic time for factoring to be feasible. ###### Theorem 3 There is an algorithm for factoring $`n`$-bit integers that consists of: a classical pre-processing stage, computed by a polynomial-size classical circuit; followed by a quantum information processing stage, computed by an $`O(\mathrm{log}n)`$-depth $`O(n^5(\mathrm{log}n)^2)`$-size quantum circuit<sup>1</sup><sup>1</sup>1In this case, the underlying circuit family is polynomial-time uniform rather than logarithmic-space uniform.; followed by a classical post-processing stage, computed by a polynomial-size classical circuit. Furthermore, the size of the quantum circuit can be reduced if a larger depth is allowed. In particular, the size can be reduced to $`O(n^3)`$ if the depth is increased to $`O((\mathrm{log}n)^2)`$. If we define the complexity class BQNC as all computational problems that can be solved by quantum circuits with poly-logarithmic depth and polynomial size—a reasonably natural extension of previous notation (see, e.g., )—then Theorem 3 implies that the factoring problem is in $`\text{ZPP}^{\text{BQNC}}`$. Finally, we consider the minimum depth required for approximating the QFT. It is fairly easy to show that computing the QFT exactly requires depth at least $`\mathrm{log}n`$. However, this is less clear in the case of approximations—and we exhibit a problem related to the QFT whose depth complexity decreases from $`\mathrm{log}n`$ in the exact case to $`O(\mathrm{log}\mathrm{log}n)`$ for approximations with precision $`\epsilon `$, whenever $`\epsilon 1/n^{O(1)}`$. Nevertheless, we show the following. ###### Theorem 4 Any quantum circuit consisting of one- and two-qubit gates that approximates the QFT with precision $`\frac{1}{10}`$ or smaller must have depth at least $`\mathrm{log}n`$. This implies that the depth upper bound in Theorem 1 is asymptotically optimal for a reasonable range of values of $`\epsilon `$. The remainder of this paper is organized as follows. In Section 2, we review some definitions and introduce notation that is used in subsequent sections. In Section 3 we prove the depth and size bounds for quantum circuits approximating the quantum Fourier transform for any power-of-2 modulus as claimed in Theorem 1, and in Section 4 we prove the size bound claimed in Theorem 2 for exactly computing the quantum Fourier transform. In Section 5 we prove Theorem 3 by demonstrating how Shor’s factoring algorithm can be arranged so as to require only logarithmic-depth quantum circuits. In Section 6 we prove the lower bound for the QFT in Theorem 4. In Section 7 we discuss the situation when the modulus for the quantum Fourier transform is not necessarily a power of 2, including arbitrary moduli and the special case of “smooth” moduli considered in Shor’s original method for performing quantum Fourier transform. We conclude with Section 8, which mentions some directions for future work relating to this paper. ## 2 Definitions and notation Notation for special quantum states: For an $`n`$-bit modulus $`m`$, we will identify each $`x_m`$ with its binary representation $`x_{n1}\mathrm{}x_1x_0\{0,1\}^n`$. For $`x_m`$, the state $`\|x=\|x_{n1}\mathrm{}x_1x_0`$ is called a computational basis state. For $`x_m`$, the state $$\|\psi _x=\frac{1}{\sqrt{m}}\underset{y=0}{\overset{m1}{}}(e^{2\pi i/m})^{xy}\|y,$$ (3) is a Fourier basis state with phase parameter $`x`$. As noted in , when $`m=2^n`$, $`\|\psi _x`$ can be factored as follows $$\|\psi _{x_{n1}\mathrm{}x_1x_0}=\frac{1}{\sqrt{2^n}}(\|0+e^{2\pi i(0.x_0)}\|1)(\|0+e^{2\pi i(0.x_1x_0)}\|1)\mathrm{}(\|0+e^{2\pi i(0.x_{n1}\mathrm{}x_1x_0)}\|1).$$ (4) For convenience, we define the state $$\|\mu _\theta =\frac{1}{\sqrt{2}}(\|0+e^{2\pi i\theta }\|1),$$ (5) where $`\theta `$ is a real parameter. Using this notation, we can rewrite Eq. 4 as $$\|\psi _{x_{n1}\mathrm{}x_1x_0}=\|\mu _{0.x_0}\|\mu _{0.x_1x_0}\mathrm{}\|\mu _{0.x_{n1}\mathrm{}x_1x_0}.$$ (6) Definition of the QFT: The quantum Fourier transform (QFT) is the unitary operation that maps $`\|x`$ to $`\|\psi _x`$ (for all $`x_m`$). Mappings related to the QFT: A quantum Fourier state computation (QFS) is any unitary operation that maps $`\|x`$$`\|0`$ to $`\|x`$$`\|\psi _x`$ (for all $`x_m`$). When the input is a computational basis state, this computes the corresponding Fourier state, but without erasing the input. We refer to approximations of a QFS as Fourier state estimation. A quantum Fourier phase computation (QFP) is any unitary operation that maps $`\|\psi _x`$$`\|0`$ to $`\|\psi _x`$$`\|x`$ (for all $`x_m`$). When the input is a Fourier basis state, this computes the corresponding phase parameter, but without erasing the input. We refer to approximations of a QFP as Fourier phase estimation. As pointed out by Kitaev , the QFT can be computed by composing a QFS and the inverse of a QFP: $`\|x\|0\|x\|\psi _x\|0\|\psi _x`$. Quantum gates: All of the quantum circuits that we construct will be composed of three types of unitary gates. One is the one-qubit Hadamard gate, $`H`$, which maps $`\|x`$ to $`\frac{1}{\sqrt{2}}(\|0+(1)^x\|1)`$ (for $`x\{0,1\}`$). Another is the one-qubit phase shift gate, $`P(\theta )`$, where $`\theta `$ is a parameter of the form $`x/2^n`$ (for $`x_{2^n}`$). $`P(\theta )`$ maps $`\|x`$ to $`e^{2\pi i\theta x}\|x`$ (for $`x\{0,1\}`$). Finally, we use two-qubit controlled-phase shift gates, $`\text{controlled-}\text{P}(\theta )`$ ($`\text{c-}\text{P}(\theta )`$ for short), which map $`\|x`$$`\|y`$ to $`e^{2\pi i\theta xy}\|x\|y`$ (for $`x,y\{0,1\}`$). Note that this set is universal, and in particular that any (classical) reversible circuit can be composed of these gates. ## 3 New depth bounds for the QFT The main purpose of this section is to prove Theorem 1. First, we review the approach of Kitaev for performing the QFT for an arbitrary modulus $`m`$. By linearity, it is sufficient to give a circuit that operates correctly on computational basis states. Given a computational basis state $`\|x`$, first create the Fourier basis state with phase parameter $`x`$ (which can be done easily if $`\|x`$ is not erased in the process). The system is now in the state $`\|x`$$`\|\psi _x`$. Now, by performing Fourier phase estimation, the state $`\|x`$$`\|\psi _x`$ can be approximated from the state $`\|0`$$`\|\psi _x`$. Therefore, by performing the inverse of Fourier phase estimation on the state $`\|x`$$`\|\psi _x`$, a good estimate of the state $`\|0`$$`\|\psi _x`$ is obtained. The particular phase estimation procedure used by Kitaev does not readily parallelize, but, in the case where the modulus is a power of 2, we give a new phase estimation procedure that does parallelize. This procedure requires several copies of the Fourier basis state rather than just one. To insure that the entire process parallelizes, we must parallelize the creation of the Fourier basis state as well as the process of copying and uncopying this state. The basic steps of our technique are as follows: 1. Creation of the Fourier basis state, which is the mapping $$\|x\|0\|x\|\psi _x.$$ 2. Copying the Fourier basis state, which is the mapping $$\|\psi _x\|0\mathrm{}\|0\|\psi _x\|\psi _x\mathrm{}\|\psi _x.$$ 3. Erasing the computational basis state by means of estimating the phase of the Fourier basis state, which is the mapping $$\|x\|\psi _x\|\psi _x\mathrm{}\|\psi _x\|0\|\psi _x\|\psi _x\mathrm{}\|\psi _x.$$ 4. Reverse step 2, which is the mapping $$\|\psi _x\|\psi _x\mathrm{}\|\psi _x\|\psi _x\|0\mathrm{}\|0.$$ Each of these components is discussed in detail in the subsections that follow. Throughout we assume the modulus is $`m=2^n`$. ### 3.1 Parallel Fourier state computation and estimation The first step is the creation of the Fourier basis state corresponding to a given computational basis state $`\|x`$. This corresponds to the mapping $$\|x\|0\|x\|\psi _x.$$ (7) First let us consider a circuit that performs this transformation exactly. By Eq. 6 (equivalently, Eq. 4), it suffices to compute the states $`\|\mu _{0.x_0},\|\mu _{0.x_1x_0},\mathrm{},\|\mu _{0.x_{n1}\mathrm{}x_1x_0}`$ individually. The circuit suggested by Figure 1 performs the required transformation for $`\|\mu _{0.x_j\mathrm{}x_0}`$. In this figure we have not labelled the controlled phase shift gates, $`\text{c-}\text{P}(\theta )`$ (such gates are defined in Section 2), which are the gates in the center drawn as two solid circles connected by a line. In the above case, the phase $`\theta `$ depends on $`j`$ and on the particular qubit of $`\|x_{n1}\mathrm{}x_1x_0`$ on which the gate acts. The value of $`\theta `$ for the controlled phase shift acting on $`\|x_i`$ is $`2^{ij1}`$ (for $`i\{0,1,\mathrm{},j\}`$). From this, it may be verified that the circuit acts as indicated. The depth of this circuit is $`O(\mathrm{log}n)`$ and the size is $`O(n)`$. If such a circuit is to be applied for each value of $`j\{0,1,\mathrm{},n1\}`$, in order to perform the mapping (7), then the qubits $`\|x_{n1},\mathrm{}\|x_1,\|x_0`$ must first be copied several times ($`ni`$ times for $`\|x_i`$) to allow the controlled phase shift gates to operate in parallel. This may be performed (and inverted appropriately) in size $`O(n^2)`$ and depth $`O(\mathrm{log}n)`$ in the most obvious way. We conclude that the transformation (7) can be performed by circuits of size $`O(n^2)`$ and depth $`O(\mathrm{log}n)`$ in the exact case. In order to reduce the size of the circuit in the approximate case, we use a similar procedure, except we only perform the controlled phase shifts when the phase $`\theta `$ is significant. An illustration of such a circuit is given in Figure 2. Here $`k`$ denotes the number of significant phase shift gates that are used. The condition $`\|\mu _{0.x_j\mathrm{}x_0}\|\mu _{0.x_j\mathrm{}x_{jk+1}}(\epsilon /n)^{O(1)}`$ occurs when $`kO(\mathrm{log}(n/\epsilon ))`$. With such a setting of $`k`$, the precision of the approximation of $`\|\mu _{0.x_{n1}\mathrm{}x_0}\mathrm{}\|\mu _{0.x_0}`$ can be $`O(\epsilon )`$. Note that the size of the resulting circuit is $`O(n\mathrm{log}(n/\epsilon ))`$ and the depth is $`O(\mathrm{log}\mathrm{log}(n/\epsilon ))`$. ### 3.2 Copying a Fourier state in parallel In this section, we show how to efficiently produce $`k`$ copies of an $`n`$-qubit Fourier state from one copy. This is a unitary operation that acts on $`k`$ $`n`$-qubit registers (thus $`kn`$ qubits in all) and maps $`\|\psi _x\|0^n\mathrm{}\|0^n`$ to $`\|\psi _x\|\psi _x\mathrm{}\|\psi _x`$ for all $`x\{0,1\}^n`$. The copying circuit will be exact and have size $`O(kn)`$ and depth $`O(\mathrm{log}(kn))`$. The setting of $`k`$ will be $`O(\mathrm{log}(n/\epsilon ))`$. Let us begin by considering the problem of producing two copies of a Fourier state from one. First, define the (reversible) addition and (reversible) subtraction operations as the mappings $`\|x`$$`\|y`$ $``$ $`\|x`$$`\|y+x`$ $`\|x`$$`\|y`$ $``$ $`\|x`$$`\|yx`$ (respectively), where $`x,y\{0,1\}^n`$ and additions and subtractions are performed as integers modulo $`2^n`$. By appealing to classical results about the complexity of arithmetic , one can construct quantum circuits of size $`O(n)`$ and depth $`O(\mathrm{log}n)`$ for these operations (using an ancilla of size $`O(n)`$). It is straightforward to show that applying a subtraction to the state $`\|\psi _x`$$`\|\psi _y`$ results in the state $`\|\psi _{x+y}`$$`\|\psi _y`$. Also, the state $`\|\psi _0`$ can be obtained from $`\|0^n`$ by applying a Hadamard transform independently to each qubit. Therefore, the copying operation can begin with a state of the form $`\|0^n`$$`\|\psi _x`$ and consist of these two steps: 1. Apply $`H`$ to each of the first $`n`$ qubits. 2. Apply the subtraction operation to the $`2n`$ qubits. The resulting state will be $`\|\psi _x`$$`\|\psi _x`$. An obvious method for computing $`k`$ copies of a Fourier state is to repeatedly apply the above doubling operation. This will result in a quantum circuit of size $`O(kn)`$; however, its depth will be $`O((\mathrm{log}k)(\mathrm{log}n))`$, which is too large for our purposes. The depth bound can be improved to $`O(\mathrm{log}(kn))`$ by applying other classical circuit constructions to efficiently implement the (reversible) prefix addition and (reversible) telescoping subtraction operations, which are the mappings $`\|x_1\|x_2\mathrm{}\|x_k`$ $``$ $`\|x_1\|x_1+x_2\mathrm{}\|x_1+x_2+\mathrm{}+x_k`$ $`\|x_1\|x_2\mathrm{}\|x_k`$ $``$ $`\|x_1\|x_2x_1\mathrm{}\|x_kx_{k1}`$ (respectively), where $`x_1,x_2,\mathrm{},x_k\{0,1\}^n`$. Before addressing the issue of efficiently implementing these operations, let us note that the copying operation can be performed by starting with the state $`\|0^n\mathrm{}\|0^n\|\psi _x`$ and performing these two steps: 1. Apply $`H`$ to all of the first $`(k1)n`$ qubits. 2. Apply the telescoping subtraction operation to the $`kn`$ qubits. The resulting state will be $`\|\psi _x\mathrm{}\|\psi _x`$. Now, to implement the prefix addition and telescoping subtraction, note that they are inverses of each other. This means that it is sufficient to implement each one efficiently by a classical (nonreversible) circuit, and then combine these to produce a reversible circuit by standard techniques in reversible computing . The telescoping subtraction clearly consists of $`k1`$ subtractions that can be performed in parallel, so the nonreversible size and depth bounds are $`O(kn)`$ and $`O(\mathrm{log}n)`$ respectively. The prefix addition is a little more complicated. It relies on a combination of well-known tools in classical circuit design. One of them is the following general result of Ladner and Fischer about parallel prefix computations. ###### Theorem 5 (Ladner and Fischer) For any associative binary operation $``$, the mapping $$(x_1,x_2,\mathrm{},x_k)(x_1,x_1x_2,\mathrm{},x_1x_2\mathrm{}x_k)$$ (8) can be computed by a circuit consisting of $`(x,y)(x,xy)`$ gates that has size $`O(k)`$ and depth $`O(\mathrm{log}k)`$. Another tool is the so-called three-two adder, which is a circuit that takes three $`n`$-bit integers $`x,y,z`$ as input and produces two $`n`$-bit integers $`s,c`$ as output, such that $`x+y+z=s+c`$ (recall that addition is in modulo $`2^n`$ arithmetic). It is remarkable that a three-two adder can be implemented with constant depth and size $`O(n)`$. By combining two three-two adders, one can implement a size $`O(n)`$ and depth $`O(1)`$ four-two adder, that performs the mapping $`(x,y,z,w)(x,y,s,c)`$, where $`x+y+z+w=s+c`$. Now, consider the pairwise representation of each $`n`$-bit integer $`z`$ as a pair of two $`n`$-bit integers $`(z^{},z^{\prime \prime })`$ such that $`z=z^{}+z^{\prime \prime }`$. This representation is not unique, but it is easy to convert to and from the pairwise representation: the respective mappings are $`z(z,0^n)`$ and $`(z^{},z^{\prime \prime })z^{}+z^{\prime \prime }`$. The useful observation is that the four-two adder performs integer addition in the pairwise representation scheme, and it does so in constant depth and size $`O(n)`$. Now, the following procedure computes prefix addition in size $`O(kn)`$ and depth $`O(\mathrm{log}k+\mathrm{log}n)=O(\mathrm{log}(kn))`$. The input is $`(x_1,x_2,\mathrm{},x_k)`$. 1. Convert the $`k`$ integers into their pairwise representation. 2. Apply the parallel prefix circuit of Theorem 5 to perform the prefix additions in the pairwise representation scheme. 3. Convert the $`k`$ integers from their pairwise representation to their standard form. The output will be $`(x_1,x_1+x_2,\mathrm{},x_1+x_2+\mathrm{}+x_k)`$, as required. Note that step 4 of the main algorithm has a circuit of identical size and depth to the one just described, as it is simply its inverse. ### 3.3 Estimating the phase of a Fourier state in parallel Finally, we will discuss the third step of the main algorithm, which corresponds to the mapping $$\|\psi _x\|\psi _x\mathrm{}\|\psi _x\|x\|\psi _x\|\psi _x\mathrm{}\|\psi _x\|0$$ (9) for $`x\{0,1\}^n`$. The number of copies of $`\|\psi _x`$ required for this step depends on the error bound $`\epsilon `$; we will require $`kO(\mathrm{log}(n/\epsilon ))`$ copies. As discussed in subsection 3.1, any Fourier basis state $`\|\psi _x`$ may be decomposed as $`\|\psi _x=\|\mu _{x2^1}\|\mu _{x2^2}\mathrm{}\|\mu _{x2^n}`$. Thus, we may assume that we have $`k`$ copies of each of the states $`\|\mu _{x2^j}`$. First, for each $`j=1,\mathrm{},n`$, the circuit will simulate measurements of the $`k`$ copies of $`\|\mu _{x2^j}`$ (in the bases discussed below) in order to obtain an approximation $`l_j/4`$ to the fractional part of $`2^jx`$. The approximation is with respect to the function $`\|\|_1`$ defined as $$\|y\|_1=\mathrm{min}\left\{z[0,1):\text{either }yz\text{ or }y+z\text{ }\right\}$$ for $`y`$ (i.e., “modulo 1” distance). With high probability the approximations will result in $`l_1,\mathrm{},l_n`$ satisfying $`\|l_j/42^jx\|_1<\frac{1}{4}`$ for each $`j`$. The (simulated) measurements can be performed in parallel, and each $`l_j`$ will be determined by considering the mode of the outcomes of the measurements and thus can be computed in parallel as well. Next the circuit will reconstruct an approximation $`\stackrel{~}{x}`$ to $`x`$ (in parallel) from $`l_1,\mathrm{},l_n`$. The circuit then XORs this value of $`\stackrel{~}{x}`$ to the register containing $`x`$, thereby “erasing” it with high probability. Finally, the circuit inverts the computation of this $`\stackrel{~}{x}`$ to clean up any garbage from the computation. As in subsection 3.2, standard techniques may be used to implement these computations as reversible circuits. We now describe each of the above steps in more detail. Let us first recall the following fact from probability theory (see, e.g., Goldreich ). If $`X_1,\mathrm{},X_t`$ are independent Bernoulli trials with probability $`p_X`$ of success and $`Y_1,\mathrm{},Y_t`$ are independent Bernoulli trials with $`p_Y`$ of success, where $`p_X<p_Y`$, then $$\mathrm{Pr}\left[\underset{i=1}{\overset{t}{}}X_i\underset{i=1}{\overset{t}{}}Y_i\right]<2e^{(p_Yp_X)^2t/2}.$$ Now, define $$\begin{array}{cc}\|b_0=\frac{1}{\sqrt{2}}\|0+\frac{1}{\sqrt{2}}\|1=\|\mu _0,\hfill & \|b_1=\frac{1}{\sqrt{2}}\|0+\frac{i}{\sqrt{2}}\|1=\|\mu _{\frac{1}{4}},\hfill \\ \|b_2=\frac{1}{\sqrt{2}}\|0\frac{1}{\sqrt{2}}\|1=\|\mu _{\frac{1}{2}},\hfill & \|b_3=\frac{1}{\sqrt{2}}\|0\frac{i}{\sqrt{2}}\|1=\|\mu _{\frac{3}{4}},\hfill \end{array}$$ and consider measurements of the states $`\|\mu _{x2^j}`$ in the bases $`\{\|b_0,\|b_2\}`$ and $`\{\|b_1,\|b_3\}`$ (these measurements correspond to measurements of the Pauli operators $`\sigma _x`$ and $`\sigma _y`$, respectively). In particular, given that we have $`k`$ copies of each $`\|\mu _{x2^j}`$, we suppose that each of the above two measurements is performed independently on $`k/2`$ of the copies. Let $`l_j\{0,1,2,3\}`$ represent the basis state that occurs with the highest frequency in these measurements for each $`j`$, breaking ties arbitrarily. We claim that the inequality $`\|l_j/42^jx\|_1<\frac{1}{4}`$ is satisfied with high probability: $$\mathrm{Pr}\left[\left\|l_j/42^jx\right\|_1\frac{1}{4}\right]<4e^{k/8}.$$ (10) To prove that the inequality (10) holds, let us suppose that $`x`$ and $`j`$ are fixed, and let us define $$p_0=\|b_0\|\mu _{x2^j}\|^2,p_1=\|b_1\|\mu _{x2^j}\|^2,p_2=\|b_2\|\mu _{x2^j}\|^2,p_3=\|b_3\|\mu _{x2^j}\|^2.$$ These are the probabilities associated with the above measurements, meaning that the probability that a measurement of $`\|\mu _{x2^j}`$ in the $`\{\|b_0,\|b_2\}`$ basis yields 0 is $`p_0`$, the probability that the measurement yields 2 is $`p_2`$, and similar for $`p_1`$ and $`p_3`$ when the measurement in the $`\{\|b_1,\|b_3\}`$ basis is performed. Now, note the following two facts: (i) it must be that $`\mathrm{max}\{p_0,p_1,p_2,p_3\}1/2+\sqrt{2}/4`$ (for any choice of $`x`$ and $`j`$), and (ii) if $`\|l/42^jx\|_1\frac{1}{4}`$, then we must have $`p_l1/2`$. Therefore, if the inequality is not satisfied for some $`j`$ (i.e., if $`\|l_j/42^jx\|_1\frac{1}{4}`$), then it must be the case that $`p_l^{}p_{l_j}\sqrt{2}/4`$ for some different value of $`l^{}l_j`$. Based on the inequalities above, we conclude that a very improbable event has taken place: the probability of the result $`l_j`$ appearing more frequently than $`l^{}`$ is at most $`2e^{k/8}`$. Unless $`\|2^jx\|_1\{0,\frac{1}{4},\frac{1}{2},\frac{3}{4}\}`$ there are at most two values of $`l_j`$ that give $`\|l_j/42^jx\|_1\frac{1}{4}`$, and so in this case we conclude that (10) holds. (In the special case $`\|2^jx\|_1\{0,\frac{1}{4},\frac{1}{2},\frac{3}{4}\}`$, the inequality (10) follows trivially.) From (10) we determine that $`\|l_j/42^jx\|_1<\frac{1}{4}`$ holds for all values of $`j`$ with probability at least $`14ne^{k/8}`$. Now consider the following problem: \| Input: \| $`l_1,\mathrm{},l_n\{0,1,2,3\}`$. \| \| --- \| --- \| \| Promise: \| There exists $`x\{0,1\}^n`$ such that $`\left\|l_j/42^jx\right\|_1<\frac{1}{4}`$ for $`j=1,\mathrm{},n`$. \| \| Output: \| $`x`$ satisfying the promise. \| The following algorithm solves this problem: * Define $$A_0=(\begin{array}{cc}1& 0\\ 0& 1\end{array}),\text{ }A_1=(\begin{array}{cc}1& 1\\ 0& 0\end{array}),\text{ }A_2=(\begin{array}{cc}0& 1\\ 1& 0\end{array}),\text{ }A_3=(\begin{array}{cc}0& 0\\ 1& 1\end{array}).$$ * Let $`x_j=A_{l_j}A_{l_{j1}}\mathrm{}A_{l_1}[2,1]`$ for each $`j`$, and output $`x=x_n\mathrm{}x_1`$. Let us now demonstrate that the algorithm is correct. We note that it is straightforward to show that for a given input $`l_1,\mathrm{},l_n`$ there is at most one $`x`$ satisfying the promise, and thus the solution is uniquely determined if the promise holds. To show that the algorithm computes this $`x`$ correctly, we prove by induction on $`j`$ that $`x_j`$ is output correctly. The set $`\{A_0,A_1,A_2,A_3\}`$ is closed under matrix multiplication, so we must have that the first column of $`A_{l_i}\mathrm{}A_{l_1}`$ is either $$e_1:=(\begin{array}{c}1\\ 0\end{array})\text{ }\text{or}\text{ }e_2:=(\begin{array}{c}0\\ 1\end{array})$$ for each $`i`$. Thus it suffices to prove that the first column of $`A_{l_j}\mathrm{}A_{l_1}`$ is $`e_1`$ if $`x_j=0`$ and is $`e_2`$ if $`x_j=1`$. The base case is $`j=1`$. Either $`x_1=0`$, in which case the fractional part of $`2^1x`$ is $`0`$, or $`x_1=1`$, in which case the fractional part of $`2^1x`$ is $`1/2`$. By the promise, we must therefore have $`l_1\{0,1\}`$ in case $`x_1=0`$ and $`l_1\{2,3\}`$ in case $`x_1=1`$. Thus the first column of $`A_{l_1}`$ is $`e_1`$ if $`x_1=0`$ and is $`e_2`$ if $`x_1=1`$ as required. Now suppose $`x_j,\mathrm{},x_1`$ are output correctly. We want to show that the first column of $`A_{l_{j+1}}\mathrm{}A_{l_1}`$ is $`e_1`$ if $`x_{j+1}=0`$ and is $`e_2`$ if $`x_{j+1}=1`$. There are four possibilities for the pair $`(x_{j+1},x_j)`$ that, along with the promise, give rise to the following implications: $$\begin{array}{c}x_{j+1}=0,x_j=0l_{j+1}\{0,1\}\hfill \\ x_{j+1}=0,x_j=1l_{j+1}\{1,2\}\hfill \\ x_{j+1}=1,x_j=0l_{j+1}\{2,3\}\hfill \\ x_{j+1}=1,x_j=1l_{j+1}\{3,0\}\hfill \end{array}$$ Suppose $`x_j=0`$, implying that the first column of $`A_{l_j}\mathrm{}A_{l_1}`$ is $`e_1`$. If $`x_{j+1}=0`$ then either $`l_{j+1}=0`$ or $`l_{j+1}=1`$, in either case implying that the first column of $`A_{l_{j+1}}\mathrm{}A_{l_1}`$ is $`e_1`$, as required. Similarly, if $`x_{j+1}=1`$ then either $`l_{j+1}=2`$ or $`l_{j+1}=3`$, in either case implying that the first column of $`A_{l_{j+1}}\mathrm{}A_{l_1}`$ is $`e_2`$, as required. The case $`x_j=1`$ is similar. Thus we have shown that the algorithm operates correctly. The above algorithm lends itself well to parallelization, following from the parallel prefix method discussed in subsection 3.2; by Theorem 5 all values of $`x_j=A_{l_j}A_{l_{j1}}\mathrm{}A_{l_1}[2,1]`$, $`j=1,\mathrm{},n`$ can be computed by a single circuit of size $`O(n)`$ and depth $`O(\mathrm{log}n)`$ (following from the fact that multiplication of the $`2\times 2`$ matrices, in modulo 2 arithmetic, can of course be done by constant-size circuits). It follows that the entire circuit for approximating the mapping (9) given $`kO(\mathrm{log}(n/\epsilon ))`$ copies of $`\|\psi _x`$ has size $`O(n\mathrm{log}(n/\epsilon ))`$ and depth $`O(\mathrm{log}n+\mathrm{log}\mathrm{log}(n/\epsilon ))=O(\mathrm{log}n+\mathrm{log}\mathrm{log}(1/\epsilon ))`$. It remains to argue that the circuit operates with error $`O(\epsilon )`$. This follows from standard results based on ideas in about converting quantum circuits that perform measurements and produce classical information with small error probability into unitary operations (without measurements) that can operate on data in superposition. It should be noted that a state $`\|\psi _x`$ can be conserved throughout the computation to ensure that errors corresponding to different values of $`x`$ are orthogonal. ## 4 New size bounds for the QFT In this section, we prove Theorem 2. Let $`F_{2^n}`$ denote the Fourier transform modulo $`2^n`$, which acts on $`n`$ qubits. The Hadamard transform is $`H=F_2`$. The standard circuit construction for $`F_{2^n}`$ can be described recursively as follows (where the two-qubit controlled-phase shift gates of the form $`\text{c-}\text{P}(\theta )`$ are defined in Section 2). Standard recursive circuit description for $`F_{2^n}`$: 1. Apply $`F_{2^{n1}}`$ to the first $`n1`$ qubits. 2. For each $`j\{1,2,\mathrm{},n1\}`$, apply $`\text{c-}\text{P}(1/2^{nj+1})`$ to the $`j^{\text{th}}`$ and $`n^{\text{th}}`$ qubit. 3. Apply $`H`$ to the $`n^{\text{th}}`$ qubit. The resulting circuit consists of $`n(n1)/2`$ two-qubit gates and $`n`$ one-qubit gates. Below is a more general recursive circuit description for $`F_{2^n}`$, parameterized by $`m\{1,\mathrm{},n1\}`$. This coincides with the above circuit when $`m=1`$. When $`m>1`$, it can be verified that the circuit does not change very much. It has exactly the same gates, though the relative order of the two-qubit gates (which all commute with each other) changes. Generalized recursive circuit description for $`F_{2^n}`$: 1. Apply $`F_{2^{nm}}`$ to the first $`nm`$ qubits. 2. For each $`j\{1,2,\mathrm{},nm\}`$ and $`k\{1,2,\mathrm{},m\}`$, apply $`\text{c-}\text{P}(1/2^{kj+1})`$ to the $`j^{\text{th}}`$ and $`(nm+k)^{\text{th}}`$ qubit. 3. Apply $`F_{2^m}`$ to the last $`m`$ qubits. Our new quantum circuits are based on this generalized recursive construction with $`m=n/2`$, except that they use a more efficient method for performing the transformation in Step 2. As is, Step 2 consists of $`(nm)m`$ two-qubit gates, which is approximately $`n^2/4`$. The key observation is that Step 2 computes the mapping which, for $`x\{0,1\}^{nm}`$ and $`y\{0,1\}^m`$, takes the state $`\|x`$$`\|y`$ to the state $`(e^{2\pi i/2^n})^{xy}\|x\|y`$, where $`xy`$ denotes the product of $`x`$ and $`y`$ interpreted as binary integers. From this, it can be shown that Step 2 can be computed using any classical method for integer multiplication in conjunction with some one-qubit phase shift gates (of the form $`P(\theta )`$, defined in Section 2). The best asymptotic circuit size for integer multiplication, due to Schönhage and Strassen , is $`O(n\mathrm{log}n\mathrm{log}\mathrm{log}n)`$, which can be translated into a reversible computation of the same size that we will denote as $`S`$. For $`x\{0,1\}^{nm}`$ and $`y\{0,1\}^m`$, $`S`$ maps the state $`\|x`$$`\|y`$$`\|0^n`$ to $`\|x`$$`\|y`$$`\|xy`$. (There are $`O(n)`$ additional ancilla qubits that are not explicitly indicated. Each of these begins and ends in state $`\|0`$.) Improved Step 2 in general circuit description for $`F_{2^n}`$: 1. Apply $`S`$ to the $`2n`$ qubits. 2. For each $`k\{1,2,\mathrm{},n\}`$ apply $`P(1/2^k)`$ to the $`(n+k)^{\text{th}}`$ qubit. 3. Apply $`S^1`$ to the $`2n`$ qubits. Using this improved Step 2 in the generalized recursive circuit description for $`F_{2^n}`$ results in a total number of gates that satisfies the recurrence $$T_n=T_{n/2}+T_{n/2}+O(n\mathrm{log}n\mathrm{log}\mathrm{log}n),$$ (11) which implies that $`T_nO(n(\mathrm{log}n)^2\mathrm{log}\mathrm{log}n)`$. It is straightforward to also show that the circuit has depth $`O(n)`$ and width $`O(n)`$ (where ancilla qubits are counted as part of the width). ## 5 Factoring via logarithmic-depth quantum circuits In this section we discuss a simple modification of Shor’s factoring algorithm that factors integers in polynomial time using logarithmic-depth quantum circuits. It is important to note that we are not claiming the existence of logarithmic-depth quantum circuits that take as input some integer $`n`$ and output a non-trivial factor of $`n`$ with high probability—the method will require (polynomial time) classical pre-processing and post-processing that is not known to be parallelizable. The motivation for this approach is that, under the assumption that quantum computers can be build, one may reasonably expect that quantum computation will be expensive while classical computation will be inexpensive. The main bottleneck of the quantum portion of Shor’s factoring algorithm is the modular exponentiation. Whether or not modular exponentiation can be parallelized is a long-standing open question, and we do not address this question here. Instead, we show that sufficient classical pre-processing allows parallelization of the part of the quantum circuit associated with the modular exponentiation. Combined with logarithmic-depth circuits for quantum Fourier transform, we obtain the result claimed in Theorem 3. In order to describe our method, let us briefly review Shor’s factoring algorithm, including the reduction from factoring to order-finding. It is assumed the input is a $`n`$-bit integer $`N`$ that is odd and composite. * (Classical) Randomly select $`a\{2,\mathrm{},N1\}`$. If $`\mathrm{gcd}(a,N)>1`$ then output $`\mathrm{gcd}(a,N)`$, otherwise continue to step 2. * (Quantum) Attempt to find information about the order of $`a`$ in $`_N`$: + Initialize a $`2n`$-qubit register and an $`n`$-qubit register to state $`\|0`$$`\|0`$. + Perform a Hadamard transform on each qubit of the first register. + (Modular exponentiation step.) Perform the unitary mapping: $$\|x\|0\|x\|a^xmodN.$$ + Perform the quantum Fourier transform on the first register and measure (in the computational basis). Let $`y`$ denote the result. * (Classical) Use the continued fraction algorithm to find relatively prime integers $`k`$ and $`r`$ such that $`0k<r<N`$ and $`\|y/2^mk/r\|2^{2n}`$. If $`a^r1(\mathrm{mod}N)`$ then continue to step 4, otherwise repeat step 2. * (Classical) If $`r`$ is even, compute $`d=\mathrm{gcd}(a^{r/2}1,N)`$ and output $`d`$ if it is a nontrivial factor of $`N`$. Otherwise go to step 1. The key observation is that much of the work required for the modular exponentiation step can be shifted to the classical computation in step 1 of the procedure. In step 1, the powers $`b_0=a`$, $`b_1=a^2modN`$, $`b_2=a^{2^2}modN,\mathrm{}`$, $`b_{2n1}=a^{2^{2n1}}modN`$ can be computed in polynomial-time. With this information available in step 2, the modular exponentiation step reduces to applying a unitary operation that maps $`\|b_0\|b_1\mathrm{}\|b_{2n1}\|x\|0`$ to $`\|b_0\|b_1\mathrm{}\|b_{2n1}\|x\|b_0^{x_0}b_1^{x_1}\mathrm{}b_{2n1}^{x_{2n1}}modN`$. This is essentially an iterated multiplication problem, where one is given $`2n`$ $`n`$-bit integers $`b_0^{x_0},b_1^{x_1},\mathrm{},b_{2n1}^{x_{2n1}}`$ as input and the goal is to compute their product. The most straightforward way to do this is to perform pairwise multiplications following the structure of a binary tree with $`2n`$ leaves. Each multiplication can be performed with depth $`O(\mathrm{log}n)`$ and size $`O(n^2)`$. The underlying binary tree has depth $`\mathrm{log}(2n)`$ and $`2n1`$ internal nodes. Thus, the entire process can be performed with depth $`O((\mathrm{log}n)^2)`$ and size $`O(n^3)`$. There are alternative methods for performing iterated multiplication achieving various combinations of depth and size. In particular, it was proved by Beame, Cook and Hoover that a product such as we have above can be computed by $`O(\mathrm{log}n)`$ depth boolean circuits of size $`O(n^5(\mathrm{log}n)^2)`$. While $`O(n^5\mathrm{log}n)`$ qubits may seem a high price to pay in order to save a factor of $`O(\mathrm{log}n)`$ in the circuit depth, the result has an interesting consequence regarding simulations of logarithmic-depth quantum circuits: if logarithmic-depth quantum circuits can be simulated in polynomial time, then factoring can be done in polynomial time as well. It should be noted that the circuits of Beame, Cook and Hoover are not logspace-uniform but rather are polynomial-time uniform; the best known bound on circuit depth for iterated products in the case of logspace uniform circuits is $`O(\mathrm{log}n\mathrm{log}\mathrm{log}n)`$ due to Reif . ## 6 Lower bounds Logarithmic-depth lower bounds for exact computations with two-qubit gates are fairly easy to obtain, based on the fact that the state of some output qubit (usually) critically depends on every input qubit. Since, by Eq. 4, the last qubit of $`\|\psi _{x_{n1}\mathrm{}x_1x_0}`$ is in state $`\frac{1}{\sqrt{2}}(\|0+e^{2\pi i(0.x_{n1}\mathrm{}x_1x_0)}\|1)`$, its value depends on all $`n`$ input qubits to the QFT when its input state is $`\|x_{n1}\mathrm{}x_1x_0`$. The depth of the circuit must be at least $`\mathrm{log}n`$ for this to be possible. This lower bound proof applies not only to the QFT, but also to QFS computations (which are defined in Section 2). This is because the output of a QFS on input $`\|x`$$`\|0`$ includes the state $`\|\psi _x`$. On the other hand, approximate computations can sometimes be performed with much lower depth than their exact counterparts. For example, in Section 3.1, it is shown that a QFS can be computed with precision $`\epsilon `$ by a quantum circuit with depth $`O(\mathrm{log}\mathrm{log}(n/\epsilon ))`$. Note that this is $`O(\mathrm{log}\mathrm{log}n)`$ whenever $`\epsilon 1/n^{O(1)}`$. Although this suggests that it is conceivable for a sub-logarithmic-depth circuit to approximate the QFT with precision $`1/n^{O(1)}`$, Theorem 4 implies that this is not possible. We now prove this theorem. Let $`C`$ be a quantum circuit that approximates the inverse QFT with precision $`\frac{1}{10}`$. In this section, since we will need to consider distances between mixed states, we adopt the trace distance as a measure of distance (see, e.g., ). The trace distance between two states with respective density operators $`\rho `$ and $`\sigma `$ is given as $$D(\rho ,\sigma )=\frac{1}{2}\text{Tr}\|\rho \sigma \|,$$ (12) where, for an operator $`A`$, $`\|A\|=\sqrt{A^{}A}`$. For a pair of pure states $`\|\varphi `$ and $`\|\varphi ^{}`$, their trace distance is $`\sqrt{1\|\varphi \|\varphi ^{}\|^2}`$, which is upper bounded by their Euclidean distance. On input $`\|\psi _{x_{n1}\mathrm{}x_1x_0}`$, the output state of $`C`$ contains an approximation of $`\|x_{n1}\mathrm{}x_1x_0`$. In particular, one of the output qubits of $`C`$ should be in a state that is an approximation of $`\|x_{n1}`$ within $`\frac{1}{10}`$. Let us refer to this as the high-order output qubit of $`C`$. If the depth of $`C`$ is less than $`\mathrm{log}n`$ then the high-order output qubit of $`C`$ cannot depend on all $`n`$ of its input qubits. Let $`k\{0,1,\mathrm{},n1\}`$ be such that the high-order output qubit does not depend on the $`k^{\text{th}}`$ input qubit (where we index the input qubits right to left starting from 0). Let $`r=nk1`$. Set $`z=2^n1`$, which is $`11\mathrm{}1=1^n`$ in binary. Following Eq. 6, $`\|\psi _z`$ can be written as $$\|\psi _z=\|\mu _{0.1}\|\mu _{0.11}\mathrm{}\|\mu _{0.1^n}.$$ (13) Consider the state $`\|\psi _{z+2^r}`$. Since $`z+2^r=0^{nr}1^r(\mathrm{mod}\mathrm{\hspace{0.17em}2}^n)`$, $$\|\psi _{z+2^r}=\|\mu _{0.1}\|\mu _{0.11}\mathrm{}\|\mu _{0.1^r}\|\mu _{0.01^r}\|\mu _{0.001^r}\mathrm{}\|\mu _{0.0^{nr}1^r}.$$ (14) Note that, on input $`\|\psi _z`$, the high-order output qubit of $`C`$ approximates $`\|1`$ with precision $`\frac{1}{10}`$; whereas, on input $`\|\psi _{z+2^r}`$, the high-order output qubit of $`C`$ approximates $`\|0`$ with precision $`\frac{1}{10}`$. Now, we consider a state $`\|\psi _z^{}`$, which has an interesting relationship with both $`\|\psi _z`$ and $`\|\psi _{z+2^r}`$. Define $$\|\psi _z^{}=\|\mu _{0.1}\|\mu _{0.11}\mathrm{}\|\mu _{0.1^r}\|\mu _{0.01^r}\|\mu _{0.1^{r+2}}\|\mu _{0.1^{r+3}}\mathrm{}\|\mu _{0.1^n}.$$ (15) The states $`\|\psi _z^{}`$ and $`\|\psi _z`$ are identical, except in their $`k^{\text{th}}`$ qubit positions (which are orthogonal: $`\|\mu _{0.01^r}`$ vs. $`\|\mu _{0.1^{r+1}}`$). Since the high-order output qubit of $`C`$ does not depend on its $`k^{\text{th}}`$ input qubit, it is the same for input $`\|\psi _z^{}`$ as for input $`\|\psi _z`$. Therefore, the state of the high-order output qubit of $`C`$ on input $`\|\psi _z^{}`$ is within $`\frac{1}{10}`$ of $`\|1`$. On the other hand, the trace distance between $`\|\psi _z^{}`$ and $`\|\psi _{z+2^r}`$ can be calculated to be below 0.7712, as follows. The two states are identical in qubit positions $`n1,n2,\mathrm{},k`$. In qubit position $`k1`$, the two states differ by an angle of $`\frac{\pi }{4}`$, in qubit position $`k2`$ the two states differ by an angle of $`\frac{\pi }{8}`$, and so on. Therefore, $`\psi _z^{}\|\psi _{z+2^r}`$ $`=`$ $`\mu _{0.1^{r+2}}\|\mu _{0.001^r}\mu _{0.1^{r+3}}\|\mu _{0.0001^r}\mathrm{}\mu _{0.1^n}\|\mu _{0.0^{nr}1^r}`$ $`=`$ $`\mathrm{cos}(\frac{\pi }{2^2})\mathrm{cos}(\frac{\pi }{2^3})\mathrm{}\mathrm{cos}(\frac{\pi }{2^{nk1}})`$ $`>`$ $`\mathrm{cos}(\frac{\pi }{2^2})\mathrm{cos}(\frac{\pi }{2^3})\mathrm{cos}(\frac{\pi }{2^4})\mathrm{}`$ $`>`$ $`0.6366,`$ where the numerical value for the last inequality is proved in Lemma 6 (below). This implies that the trace distance between $`\|\psi _z^{}`$ and $`\|\psi _{z+2^r}`$ is less than $`\sqrt{1(0.6366)^2}=0.7712`$. Since the trace distance is contractive, it follows that the state of the high-order output of $`C`$ on input $`\|\psi _z^{}`$ has trace distance less than 0.7712 from the state of high-order output of $`C`$ on input $`\|\psi _{z+2^r}`$. But, by the triangle inequality, this implies that the trace distance between $`\|0`$ and $`\|1`$ is less than $`\frac{1}{10}+0.7712+\frac{1}{10}<1`$, which is a contradiction, since $`\|0`$ and $`\|1`$ are orthogonal. This completes the proof of Theorem 4. ###### Lemma 6 $`\mathrm{cos}(\frac{\pi }{2^2})\mathrm{cos}(\frac{\pi }{2^3})\mathrm{cos}(\frac{\pi }{2^4})\mathrm{}>0.6366`$. Proof: We first lower bound the tails of the above infinite product by showing that, for any $`i1`$, $`\mathrm{cos}(\frac{\pi }{2^{i+1}})\mathrm{cos}(\frac{\pi }{2^{i+2}})\mathrm{cos}(\frac{\pi }{2^{i+3}})\mathrm{}>1\frac{\pi ^2}{64^i}`$. Since, for $`t>0`$, $`\mathrm{cos}(t)>1\frac{t^2}{2}`$, $`\mathrm{cos}(\frac{\pi }{2^{i+1}})\mathrm{cos}(\frac{\pi }{2^{i+2}})\mathrm{cos}(\frac{\pi }{2^{i+3}})\mathrm{}`$ $`>`$ $`\left(1\frac{\pi ^2}{24^{i+1}}\right)\left(1\frac{\pi ^2}{24^{i+2}}\right)\left(1\frac{\pi ^2}{24^{i+3}}\right)\mathrm{}`$ $``$ $`1\frac{1}{2}\left(\frac{\pi ^2}{4^{i+1}}+\frac{\pi ^2}{4^{i+2}}+\frac{\pi ^2}{4^{i+3}}+\mathrm{}\right)`$ $`=`$ $`1\frac{\pi ^2}{64^i}.`$ Now it follows that, for any $`i1`$, $`\mathrm{cos}(\frac{\pi }{2^2})\mathrm{cos}(\frac{\pi }{2^3})\mathrm{cos}(\frac{\pi }{2^4})\mathrm{}>\mathrm{cos}(\frac{\pi }{2^2})\mathrm{}\mathrm{cos}(\frac{\pi }{2^i})(1\frac{\pi ^2}{64^i})`$. Setting $`i=8`$ in this inequality gives the numerical lower bound. ## 7 Other moduli In this section we discuss the quantum Fourier transform with respect to moduli that are not powers of 2. First we briefly sketch a method for performing (in parallel) the QFT for an arbitrary modulus that uses the QFT with a power of 2 modulus as a black box. We then discuss Shor’s original method for performing the QFT with respect to a “smooth” modulus, and mention how this method may be parallelized as well. ### 7.1 Arbitrary moduli Consider the QFT with respect to an arbitrary modulus $`m`$. In this subsection we note that it is possible to approximate such a QFT with high accuracy in parallel using circuits for the quantum Fourier transform modulo $`2^k`$ for $`k=\mathrm{log}m+O(1)`$. Using the circuits for the quantum Fourier transform modulo $`2^k`$ described previously, we have that for any $`\epsilon `$ and modulus $`m`$ there exists a depth $`O(\mathrm{log}n\mathrm{log}\mathrm{log}(1/\epsilon ))`$ quantum circuit that approximates the QFT modulo $`m`$ to within accuracy $`\epsilon `$ for $`n=\mathrm{log}m`$. The size of the circuit is polynomial in $`n+\mathrm{log}(1/\epsilon )`$. The method exploits a relation between QFTs with different moduli that was used by Hales and Hallgren in regard to the Fourier Sampling problem (see also Høyer for an extension and simplified proof). The basic components of the technique are as follows: 1. Create a Fourier state with modulus $`m`$, which is the mapping $$\|x\|0\|x\|\psi _x.$$ 2. Copy the Fourier state, which is the mapping $$\|x\|\psi _x\|0\mathrm{}\|0\|x\|\psi _x\|\psi _x\mathrm{}\|\psi _x.$$ 3. Apply the inverse Fourier transform modulo $`2^k`$ on each state $`\|\psi _x`$, which is the mapping $$\|x\|\psi _x\mathrm{}\|\psi _x\|x\left(F_{2^k}^{}\|\psi _x\right)\mathrm{}\left(F_{2^k}^{}\|\psi _x\right).$$ 4. For each (computational basis state) $`y`$ occurring among the collections of qubits on which $`F_k^{}`$ was performed, compute $`\mathrm{round}(ym\mathrm{\hspace{0.17em}2}^k)`$ $`modm`$, and compute the mode of these results. With high probability the result will be $`x`$. (A reasonably straightforward calculation shows that observation of $`F_{2^k}^{}\|\psi _x`$ in the computational basis yields some $`y`$ with $`\mathrm{round}(ym2^k)=x`$ with probability greater than $`1/2+\delta `$ for some constant $`\delta `$.) XOR this result to the qubits in state $`\|x`$, and reverse the computation of each $`\mathrm{round}(ym\mathrm{\hspace{0.17em}2}^k)`$ and $`y`$. With high probability the mapping $$\|x\left(F_{2^k}^{}\|\psi _x\right)\mathrm{}\left(F_{2^k}^{}\|\psi _x\right)\|0\left(F_{2^k}^{}\|\psi _x\right)\mathrm{}\left(F_{2^k}^{}\|\psi _x\right)$$ has been performed. 5. Reverse steps 3 and 2, giving the mapping $`\|0\left(F_{2^k}^{}\|\psi _x\right)\mathrm{}\left(F_{2^k}^{}\|\psi _x\right)\|0\|\psi _x\|0\mathrm{}\|0`$. Unfortunately some of the methods used in the power of 2 case (such as using three-two adders and approximating the individual qubits of the Fourier basis states) do not seem to work in this case, which results in the slightly worse depth bound. The overall size bound increases as well, but is still polynomial. It is interesting to note that this method does not require the larger modulus to be a power of 2—effectively the method shows that the QFT modulo $`m`$ for any modulus $`m`$ can be efficiently approximated given a black box that approximates the QFT modulo $`m^{}`$ for any sufficiently large $`m^{}`$. The technical details regarding this method will appear in the final version of this paper. ### 7.2 Shor’s “mixed-radix” QFT We conclude with a brief discussion of Shor’s original “mixed radix” method for computing the quantum Fourier transform, as it too can be parallelized (although to our knowledge not as efficiently as the power-of-2 case discussed previously in this paper). Shor’s original method for computing the QFT is based on the Chinese Remainder Theorem and its consequences regarding $`_m`$ for given modulus $`m`$. Here the modulus is $`m=m_1m_2\mathrm{}m_k`$ for $`m_1,\mathrm{},m_k`$ pairwise relatively prime and $`m_jO(\mathrm{log}m)`$. Thus $`kO(\mathrm{log}m/\mathrm{log}\mathrm{log}m)`$ is somewhat less than the number of bits of $`m`$, and each $`m_j`$ has length logarithmic in the length of $`m`$. Taking $`m_j`$ to be the $`j^{\text{th}}`$ prime results in a sufficiently dense collection of moduli $`m`$ for factoring (see Rosser and Schoenfeld for explicit bounds and a detailed analysis of such bounds). Although stated somewhat differently by Shor, the mixed radix QFT method may be described as follows: 1. For $`j=1,\mathrm{},k`$ define $`f_j=\frac{m}{m_j}`$ and set $`g_j\{0,\mathrm{},m_j1\}`$ such that $`g_jf_j^1(modm_j)`$. 2. Define $`C`$ to be the (reversible) operator acting as follows for each $`x\{0,\mathrm{},m1\}`$: $$C:\|x\|(xmodm_1),\mathrm{},(xmodm_k)$$ 3. Define $`A`$ to be a (reversible) operator such that $$A:\|x_1,\mathrm{},x_k\|g_1x_1,\mathrm{},g_kx_k$$ for each $`(x_1,\mathrm{},x_k)\{0,\mathrm{},m_11\}\times \mathrm{}\times \{0,\mathrm{},m_k1\}`$. 4. Let $`F_m`$ and $`F_{m_j}`$ denote the QFT for moduli $`m`$ and $`m_j`$, $`j=1,\mathrm{},k`$, respectively. Then the following relation holds: $$F_m=C^{}(F_{m_1}\mathrm{}F_{m_k})AC.$$ (16) Thus, to perform the QFT modulo $`m`$ on $`\|x`$, first convert $`x`$ to its modular representation $`(x_1,\mathrm{},x_k)`$ using the operator $`C`$, multiply each $`x_j`$ by $`g_j`$ (modulo $`m_j`$), perform the QFT modulo $`m_j`$ independently on coefficient $`j`$ (for each $`j`$), then apply the inverse of $`C`$ to convert back to the ordinary representation of elements in $`\{0,\mathrm{},m1\}`$. The numbers computed in step 1 are used in the standard proof of the Chinese Remainder Theorem: given $`x_1,\mathrm{},x_k`$, we may compute $`x\{0,\mathrm{},m1\}`$ satisfying $`xx_j(modm_j)`$ for each $`j`$ by taking $`x=_{j=1}^kf_jg_jx_jmodm`$. Thus the operator $`C`$ can be implemented efficiently, since the mappings $`x((xmodm_1),\mathrm{},(xmodm_k))`$ and $`((xmodm_1),\mathrm{},(xmodm_k))x`$ are efficiently computable (e.g., with size $`O(\mathrm{log}^2m)`$ circuits ). In the present case $`C`$ can be parallelized to logarithmic depth, since each of the moduli are small. Similarly, the operator $`A`$ can be parallelized to logarithmic depth. To see that the relation (16) holds, we may simply examine the action of the operator on the right hand side on computational basis states: $`C^{}(F_{m_1}\mathrm{}F_{m_k})AC\|x`$ $`=`$ $`C^{}(F_{m_1}\mathrm{}F_{m_k})\|g_1x_1,\mathrm{},g_kx_k`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{m}}}C^{}{\displaystyle \underset{y_1,\mathrm{},y_k}{}}\mathrm{exp}(2\pi ig_1x_1y_1/m_1)\mathrm{}\mathrm{exp}(2\pi ig_kx_ky_k/m_k)\|y_1,\mathrm{},y_k`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{m}}}C^{}{\displaystyle \underset{y_1,\mathrm{},y_k}{}}\mathrm{exp}(2\pi i(f_1g_1x_1y_1+\mathrm{}+f_kg_kx_ky_k)/m)\|y_1,\mathrm{},y_k`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{m}}}{\displaystyle \underset{y_1,\mathrm{},y_k}{}}\mathrm{exp}(2\pi i(f_1g_1x_1y_1+\mathrm{}+f_kg_kx_ky_k)/m)\|f_1g_1y_1+\mathrm{}+f_kg_ky_k(modm)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{m}}}{\displaystyle \underset{y}{}}\mathrm{exp}(2\pi ixy/m)\|y`$ $`=`$ $`F_m\|x`$ Finally, each of the independent QFTs modulo $`m_1,\mathrm{},m_k`$ can of course be done in parallel. Here, however, a problem arises if our goal is to parallelize the entire process. Originally Shor suggests implementing each of these operations by circuits of size $`m_j`$ (not $`\mathrm{log}m_j`$), since any quantum operation can be computed by circuits with exponential-size quantum circuits . This results in a linear-depth circuit overall, although the circuit will be exact. However, we may try to compute each $`F_{m_j}`$ more efficiently. There are a few possibilities for how to do this, all (apparently) requiring approximations of each $`F_{m_j}`$. First, we may apply the method of Kitaev to approximate these QFTs. Alternately, we may use the arbitrary modulus method we have proposed in section 7.1. Finally, we have noted that this method works for any two moduli (not just for the larger modulus a power of 2) so that we may in fact recurse using the mixed-radix method to approximate each $`F_{m_j}`$. In all cases, our analysis has revealed that the mixed radix method results in worse size and/or depth bounds than the power of 2 method presented in Section 3. ## 8 Conclusion We have proved several new bounds on the circuit complexity of approximating the quantum Fourier transform, and have applied these bounds to the problem of factoring using quantum circuits. There are several related open questions, a few of which we will now discuss. First, is it possible to perform the quantum Fourier transform exactly using logarithmic- or poly-logarithmic-depth quantum circuits? The best currently known upper bound on the depth of the exact QFT is linear in the number of input qubits. Next, can the efficiency of our techniques be improved significantly? We have concentrated on asymptotic analyses of our circuits, and we believe it is certain that our circuits can be optimized significantly for “interesting” input sizes (perhaps several hundred to a few thousand qubits). Finally, the fact that the quantum Fourier transform can be performed in logarithmic depth suggests the following question: are there interesting natural problems in BQNC (bounded-error quantum NC) not known to be in NC or RNC? For instance, computing the gcd of two $`n`$-bit integers and computing $`a^b\mathrm{mod}c`$ and $`a^1\mathrm{mod}c`$ for $`n`$-bit integers $`a`$, $`b`$, and $`c`$ is not known to be possible using polynomial-size circuits with depth poly-logarithmic in $`n`$ in the classical setting. Are there logarithmic- or poly-logarithmic-depth quantum circuits for these problems? Greenlaw, Hoover and Ruzzo list several other problems not known to be classically parallelizable, all of which are interesting problems to consider in the quantum setting. ### Acknowledgments We thank Wayne Eberly for helpful discussions on classical circuit complexity, and Chris Fuchs and Patrick Hayden for an informative discussion regarding quantum state distance measures. (The Los Alamos Preprint Archive may be found at http://xxx.lanl.gov/ on the World Wide Web.)
warning/0006/math0006070.html	ar5iv	text	# Determinants of Hankel Matrices ## 1 Introduction The main purpose of this paper is to compute asymptotically the determinants of the finite matrices $`H_n(u)`$ defined by $$det(a_{i+j})_{i,j=0}^{n1}$$ where $$a_{i+j}=_0^{\mathrm{}}x^{i+j}u(x)𝑑x$$ with suitable conditions on the weight function $`u(x).`$ These determinant entries depend only the sum $`i+j`$ and are hence classified as Hankel determinants. Hankel determinants such as these were considered by Szegö in and also by Hirshmann in , but in both cases for finite intervals. These determinants are important in random matrix theory and its applications. Our main result is as follows. Suppose we replace $`u(x)`$ by a function given in the form $`w(x)U(x)`$ where $`w(x)`$ is the weight $`e^xx^\nu `$ with $`\nu 1/2`$. Then if $`U`$ is nowhere zero and $`U1`$ is a Schwartz function (a condition which can be considerably relaxed) the determinants are given asymptotically as $`n\mathrm{}`$ by $`det(H_n(u))=\mathrm{exp}\{c_1n^2\mathrm{log}n+c_2n^2+c_3n\mathrm{log}n+c_4n+c_5n^{1/2}+c_6\mathrm{log}n+c_7+o(1)\}`$ (1) where $$c_1=1,c_2=3/2,c_3=\nu ,$$ $$c_4=\nu +\mathrm{log}2\pi ,c_5=\frac{2}{\pi }_0^{\mathrm{}}\mathrm{log}(U(x^2))𝑑x$$ $$c_6=\nu ^2/21/6,c_7=4/3\mathrm{log}G(1/2)+(1/3+\nu /2)\mathrm{log}\pi +(\nu /21/18)\mathrm{log}2$$ $$\mathrm{log}G(1+\nu )\nu /2\mathrm{log}U(0)+\frac{1}{2\pi ^2}_0^{\mathrm{}}xS(x)^2𝑑x,$$ $`G`$ is the Barnes G-function, and $`S(x)=_0^{\mathrm{}}\mathrm{cos}(xy)\mathrm{log}U(y^2)𝑑y.`$ The idea behind the proof is to replace the matrix $`H_n(u)`$ with one whose $`i,j^{th}`$ entry is given by $$_0^{\mathrm{}}P_i(x)P_j(x)u(x)𝑑x$$ where the $`P_i`$’s are orthogonal (Laguerre) polynomials with respect to the weight $`e^xx^\nu .`$ These new determinants can then be evaluated using the ideas of the “linear statistic” method in random matrix theory. More precisely, the above determinant can be replaced by $$det(I+K_n)$$ where $`K_n`$ is an integral operator whose kernel is given by $$\underset{i=0}{\overset{n1}{}}L_i^\nu (x)L_i^\nu (y)(U1)(y)$$ with $`L_i^\nu `$ defined as the $`i^{th}`$ Laguerre function of order $`\nu .`$ The main computation in the paper is to approximate the above kernel with a different kernel which involves Bessel functions. This new kernel was fortunately already considered in . There, asymptotics for certain integral operators were computed and these results are then applied to give the result of formula (1). The kernels considered in arises in random matrix theory in the “hard-edge” scaling for ensembles of positive Hermitian matrices. Details about the random matrix connections can be found in . The formula in (1) was stated earlier in where a heuristic argument using the Coulomb gas approach was used to derive the same formula. The Coulomb gas approach was also used in to extend to the case where the interval of integration is the entire real line ## 2 Preliminaries We first show how to replace the powers in the Hankel determinants with the orthogonal polynomials. Let $$_n(u)=\left(_0^{\mathrm{}}P_i(x)P_j(x)u(x)dx\right)\|_{i,j=0}^{n1},$$ and write $$P_i(x)=\underset{ki}{}a_{ik}x^k.$$ ###### Lemma 2.1 Let $`P_i,_n(u),H_n(u)`$ be defined as above and let us assume that all moment integrals exist. Then $$det(_n(u))=A_ndet(H_n(u)),$$ where $$A_n=(\underset{i=0}{\overset{n1}{}}a_{ii})^2.$$ Proof. We have $$det(_n(u))=det(_0^{\mathrm{}}(\underset{mj}{}a_{jm}x^m)(\underset{lk}{}a_{kl}x^l)u(x)𝑑x).$$ The $`j,k`$ entry of the matrix is given by $$\underset{mj}{}\underset{lk}{}a_{jm}a_{kl}_0^{\mathrm{}}x^{m+l}u(x)𝑑x$$ $$=\underset{mj}{}\underset{lk}{}a_{jm}_0^{\mathrm{}}x^{m+l}u(x)𝑑xa_{kl}.$$ Form this it follows that $`_n(u)`$ is a product of three matrices $`TH_n(u)T^t.`$ The matrix $`T`$ is lower triangular with entries $`a_{jm},mj.`$ It is easy to see from this that the lemma follows. The next step is to evaluate the term $$A_n=(\underset{i=0}{\overset{n1}{}}a_{ii})^2,$$ which is of course straight forward. We normalize the polynomials so that they are orthonormal. Then it is well known that $$a_{ii}^2=\frac{1}{\mathrm{\Gamma }(1+i+\nu )\mathrm{\Gamma }(1+i)}.$$ The following lemma computes the product asymptotically using the Barnes G-function. This function is defined by $`G(1+z)`$ $`=`$ $`(2\pi )^{z/2}e^{(z+1)z/2\gamma _Ez^2/2}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left((1+z/k)^ke^{z+z^2/2k}\right)`$ (2) with $`\gamma _E`$ being Euler’s constant. ###### Lemma 2.2 $$A_n^1=\underset{i=0}{\overset{n1}{}}a_{ii}^2$$ is given asymptotically by $$\mathrm{exp}\left\{d_1n^2\mathrm{log}n+d_2n^2+d_3n\mathrm{log}n+d_4n+d_5\mathrm{log}n+d_6+o(1)\right\}$$ where $$d_1=1,d_2=3/2,d_3=\nu ,d_4=\nu +\mathrm{log}2\pi ,d_5=\nu ^2/21/6$$ $$d_6=4/3\mathrm{log}G(1/2)+(1/3+\nu /2)\mathrm{log}\pi +(\nu /21/18)\mathrm{log}2\mathrm{log}G(1+\nu )$$ Proof. It is well known that the Barnes G function satisfies the property $$G(1+z)=\mathrm{\Gamma }(z)G(z).$$ From this it is quite easy to see that $`A_n^1`$ can be written as $$\frac{G(1+n)G(1+n+\nu )}{G(1+\nu )}.$$ The asymptotics of the Barnes function are computed in and since $`G(1+a+n)`$ is asymptotic to $$n^{(n+a)^2/21/12}e^{3/4n^2an}(2\pi )^{(n+a)/2}G^{2/3}(1/2)\pi ^{1/6}2^{1/36}$$ we can directly apply this formula with $`a=0`$ and $`a=\nu `$ to obtain the desired result. This last result shows then that the asymptotics of the Hankel matrices can be reduced to those of the matrices $`_n(u).`$ We compute these asymptotics by replacing the determinant with a Fredholm determinant. To do this we must first consider some estimates for the various kernels. ## 3 Hilbert-Schmidt and trace norm estimates We proceed to write the determinant of $`_n(u)`$ as a Fredholm determinant. It is clear that if $`U(x)`$ is a bounded function then the matrix $`_n(u)`$ can be realized as $`P_nM_UP_n`$ where $`M_U`$ is multiplication by $`U`$ and $`P_n`$ is the projection onto the space spanned by the first $`n`$ Laguerre functions. This is clearly a bounded, finite rank operator defined on $`L_2(0,\mathrm{}).`$ Thus $$det_n(u)=det(I+K_n)$$ is defined where $`K_n`$ is the integral operator whose kernel is given by $$\underset{i=0}{\overset{n1}{}}L_i^\nu (x)L_i^\nu (y)(U(y)1).$$ Let us write a square root of the function $`U1`$ as $`V.`$ Then $$det(I+K_n)=det(I+P_nM_VM_VP_n)=det(I+M_VP_nM_V).$$ This last equality uses the general fact that $$det(I+AB)=det(I+BA)$$ for general operators defined on a Hilbert space. At this point we have replaced the kernel $`K_n`$ with $$\underset{i=0}{\overset{n1}{}}V(x)L_i^\nu (x)L_i^\nu (y)V(y).$$ Our primary goal in the paper is to replace this last kernel, which we will commonly refer to as the Laguerre kernel, with a more familiar one namely the compressed Bessel kernel, also defined on $`L_2(0,\mathrm{}),`$ $`V(x){\displaystyle \frac{J_\nu (2(nx)^{1/2})\sqrt{ny}J_\nu ^{^{}}(2(ny)^{1/2})J_\nu (2(ny)^{1/2})\sqrt{nx}J_\nu ^{^{}}(2(nx)^{1/2})}{xy}}V(y).`$ (3) We use the term “compressed” since the above without the factors $`V(x)`$ and $`V(y)`$ is generally called the Bessel kernel. We will use the general fact that if we have families of trace class operators $`A_n`$ and $`B_n`$ (thinking of the Laguerre kernel as $`A_n`$ and the compressed Bessel kernel as $`B_n`$) that depend on a parameter $`n`$ such that the Hilbert Schmidt norm $`A_nB_n_2=o(1)`$ and $`\text{tr }(A_nB_n)(I+B_n)^1=o(1)`$ then $$det(I+A_n)/det(I+B_n)1$$ as $`n\mathrm{}.`$ The proof of this fact uses the idea of generalized determinants found in , and also requires uniformity for the norms of the inverses of the operators $`I+B_n.`$ This will be made more precise later. Using the standard identity for the sum of the products of Laguerre functions allows us to write the Laguerre kernel $$\underset{i=0}{\overset{n1}{}}V(x)L_i^\nu (x)L_i^\nu (y)V(y)$$ as $$\sqrt{n(n+\nu )}V(x)\frac{L_{n1}^\nu (x)L_n^\nu (y)L_{n1}^\nu (y)L_n^\nu (x)}{xy}V(y).$$ Notice that the above kernel has somewhat the same form as the compressed Bessel kernel. It turns out that both the kernel above and the Bessel kernel have an integral representation that make the Hilbert-Schmidt computations simpler. The next lemma shows how this is done for in the Laguerre case. ###### Lemma 3.1 We have $$\frac{L_{n1}^\nu (x)L_n^\nu (y)L_{n1}^\nu (y)L_n^\nu (x)}{xy}=\frac{1}{2}_0^1[L_{n1}^\nu (tx)L_n^\nu (ty)+L_{n1}^\nu (ty)L_n^\nu (tx)]𝑑t.$$ Proof. Call the left-hand side $`\mathrm{\Phi }(x,y)`$ and the right-hand side $`\mathrm{\Psi }(x,y)`$. It follows from the differentiation formulas for Laguerre polynomials that there is a differentiation formula $$x\frac{d}{dx}\left(\begin{array}{c}L_{n1}^\nu (x)\\ \\ L_n^\nu (x)\end{array}\right)=\left(\begin{array}{cc}A& B\\ & \\ C& A\end{array}\right)\left(\begin{array}{c}L_{n1}^\nu (x)\\ \\ L_n^\nu (x)\end{array}\right),$$ where $$A(x)=\frac{1}{2}x\frac{\nu }{2}N,B(x)=C(x)=\sqrt{N(N+\nu )}.$$ An easy computation using this shows that $`\mathrm{\Phi }(x,y)`$ satisfies $$\left(x\frac{}{x}+y\frac{}{y}+1\right)\mathrm{\Phi }(x,y)=\frac{1}{2}[L_{n1}^\nu (x)L_n^\nu (y)+L_{n1}^\nu (y)L_n^\nu (x)].$$ On the other hand, if we temporarily assume that $`\nu >0`$ and differentiate under the integral sign and then integrate by parts we find that $$\left(x\frac{}{x}+y\frac{}{y}+1\right)\mathrm{\Psi }(x,y)=\frac{1}{2}[L_{n1}^\nu (x)L_n^\nu (y)+L_{n1}^\nu (y)L_n^\nu (x)].$$ Hence $$\left(x\frac{}{x}+y\frac{}{y}+1\right)(\mathrm{\Phi }(x,y)\mathrm{\Psi }(x,y))=0.$$ From this it follows that $`\mathrm{\Phi }(x,y)\mathrm{\Psi }(x,y)`$ is of the form $`\phi (x/y)e^{(x+y)/2}`$ for some function $`\phi `$ of one variable. Still assuming $`\nu >0`$, both $`\mathrm{\Phi }(x,y)`$ and $`\mathrm{\Psi }(x,y)`$ tend to 0 as $`x`$ and $`y`$ tend to 0 independently. This shows that $`\phi =0`$ and so the identity is established when $`\nu >0`$. Since both sides of the identity are analytic functions of $`\nu `$ for $`\text{Re }\nu >1`$ the identity holds generally. Thus our kernel may be rewritten $$\frac{1}{2}\sqrt{n(n+\nu )}V(x)V(y)_0^1[L_{n1}^\nu (tx)L_n^\nu (ty)+L_{n1}^\nu (ty)L_n^\nu (tx)]𝑑t.$$ (4) We assume from now on that $`\|\nu \|1/2`$. The compressed Bessel kernel appearing in (3) also has the following well-known integral representation $`nV(x)V(y){\displaystyle _0^1}J_\nu (2\sqrt{nxt})J_\nu (2\sqrt{nyt})𝑑t`$ (5) and the rest of the section is devoted to replacing the Laguerre functions by the appropriate Bessel functions in (4) and showing that this leads to a small Hilbert-Schmidt error and then showing that certain traces tend to zero. Notice that in the above notation kernel (4) is the same as $`A_n`$ and kernel (5) is the same as $`B_n.`$ We need only the following fact about the Laguerre functions. ###### Lemma 3.2 Let $`a`$ be any real positive constant and suppose that $`\nu 1/2.`$ Then as $`n\mathrm{},`$ the normalized Laguerre functions satisfy $$\underset{0<xa}{\mathrm{max}}x^{1/4}L_n^\nu (x)=O(n^{1/4}),$$ $$\underset{xa}{\mathrm{max}}L_n^\nu (x)=O(n^{1/4}).$$ We remark that the implied constants in the estimates may depend on $`a`$ but not on $`n`$. Proof. These estimates follow easily from formula (7.6.9) and Theorem 8.91.2 in . ###### Lemma 3.3 We have $$L_n^\nu (x)c_{n,\nu }J_\nu (2\sqrt{Nx})=\frac{c_{n,\nu }}{6\sqrt{\pi }}N^{3/4}x^{5/4}\mathrm{sin}(2\sqrt{Nx}\alpha )+O(N^{5/4}(x^{1/4}+x^3)),$$ (6) where $`N=n+(\nu +1)/2,c_{n,\nu }=N^{\nu /2}(\frac{\mathrm{\Gamma }(n+\nu +1)}{\mathrm{\Gamma }(n+1)})^{1/2},\alpha =(\nu /2+1/4)\pi `$ and the constant implied in the O depends only on $`\nu `$. Proof. Formula (8.64.3) of Szegö reads, in current notation, $$L_n^\nu (x)c_{n,\nu }J_\nu (2\sqrt{Nx})$$ $$=\frac{\pi }{4\mathrm{sin}\nu \pi }_0^x[J_\nu (2\sqrt{Nx})J_\nu (2\sqrt{Nt})J_\nu (2\sqrt{Nx})J_\nu (2\sqrt{Nt})]tL_n^\nu (t)𝑑t.$$ Using the Lemma 3.2 estimate for $`L_n^\nu (t)`$ we find that if $`Nx<1`$ the right side is at most a constant times $$(Nx)^{\|\nu \|/2}_0^x(Nt)^{\|\nu \|/2}tL_n^\nu (t)𝑑t=O(N^{\|\nu \|\frac{1}{4}}x^{\frac{7}{4}\|\nu \|})=O(N^2).$$ (If $`\nu `$ is an integer this must be multiplied by $`\mathrm{log}N`$.) So we assume $`Nx>1`$. The integral over $`Nt<1`$ is at most a constant times $$(Nx)^{1/4}_0^{1/N}(Nt)^{\|\nu \|/2}N^{1/4}t^{\frac{3}{4}}𝑑t=O(N^{\frac{9}{4}}x^{1/4})$$ with, possibly, an extra $`\mathrm{log}N`$ factor. So we confine our integral now to $`Nt>1`$. If we replace the Bessel functions by their first-order asymptotics the error is at most a constant times $$N^1_0^x(x^{1/2}+t^{1/2})(xt)^{1/4}N^{1/4}(t+t^{3/4})𝑑t=O(N^{5/4}(1+x)).$$ Therefore with this error we can replace the Bessel functions by their first-order asymptotics, obtaining, after using some some trigonometric identities, $$\frac{1}{4}N^{1/2}x^{1/4}_{1/N}^x\mathrm{sin}(2\sqrt{Nx}2\sqrt{Nt})t^{3/4}L_n^\nu (t)𝑑t.$$ (7) Notice that this has the uniform estimate $`O(N^{3/4}(1+x^{3/2}))`$, and the earlier errors combined are $`O(N^{5/4}(x^{1/4}+x)`$. It follows that if in this integral we replace $`L_n^\nu (t)`$ by $`c_{n,\nu }J_\nu (2\sqrt{Nt})`$ the error is $`O(N^{5/4}(1+x^3))`$. Then if we replace $`J_\nu (2\sqrt{Nt})`$ by its first-order asymptotics the error is at most a constant times $$N^{1/2}x^{1/4}_{1/N}^xt^{3/4}(Nt)^{3/4}𝑑t=O(N^{5/4}x^{3/4}).$$ Hence with the sum of the last-mentioned errors we may replace the factor $`L_n^\nu (t)`$ in (7) by the first-order asymptotics of $`c_{n,\nu }J_\nu (2\sqrt{Nt})`$. Using a trigonometric identity shows that this results in $$\frac{c_{n,\nu }}{8\sqrt{\pi }}N^{3/4}x^{1/4}_{1/N}^x[\mathrm{sin}(2\sqrt{Nx}\alpha )\mathrm{sin}(2\sqrt{Nx}4\sqrt{Nt}+\alpha )]t^{1/2}𝑑t$$ $$=\frac{c_{n,\nu }}{6\sqrt{\pi }}N^{3/4}x^{5/4}\mathrm{sin}(2\sqrt{Nx}\alpha )+O(N^{9/4}x^{1/4})+O(N^{5/4}x^{3/4}).$$ Putting these things together gives the statement of the lemma. In what follows we use the notation $`o_2()`$ or $`O_2()`$ for a family of operators whose Hilbert-Schmidt norm satisfies the corresponding $`o`$ or $`O`$ estimate, or for a kernel whose associated operator satisfies the estimate. Similarly, $`o_1()`$ and $`O_1()`$ refer to the trace norm. We also set $`N^{}=n+(\nu 1)/2`$. ###### Lemma 3.4 . Suppose $`V`$ is in $`L^{\mathrm{}}`$ and satisfies $`_0^{\mathrm{}}\|V(x)\|^2(x^{1/2}+x^6)𝑑x<\mathrm{}`$. Then the difference between the kernel (4) and $$\frac{1}{2}\sqrt{n(n+\nu )}V(x)V(y)c_{n1,\nu }c_{n,\nu }_0^1[J_\nu (2\sqrt{N^{}tx})J_\nu (2\sqrt{Nty})+J_\nu (2\sqrt{N^{}ty})J_\nu (2\sqrt{Ntx})]𝑑t$$ is equal to (i) $`O_1(N^{1/2})`$ plus a constant which is $`O(1)`$ times $$N^{1/4}_0^1\left[(tx)^{5/4}\mathrm{sin}(2\sqrt{N^{}tx}\alpha )J_\nu (2\sqrt{Nty})+(ty)^{5/4}\mathrm{sin}(2\sqrt{Nty}\alpha )J_\nu (2\sqrt{N^{}tx})\right]V(x)V(y)𝑑t$$ plus a similar term with $`x`$ and $`y`$ interchanged; (ii) $`O_2(N^{1/4})`$. Proof. The estimate of Lemma 3.3 holds with $`n`$ replaced by $`n1`$ if $`N`$ is replaced by $`N^{}`$. Write $$L_{n1}^\nu (x)L_n^\nu (y)c_{n1,\nu }c_{n,\nu }J_\nu (2\sqrt{N^{}x})J_\nu (2\sqrt{Ny})$$ $$=[L_{n1}^\nu (x)c_{n1,\nu }J_\nu (2\sqrt{N^{}x})]L_n^\nu (y)+L_{n1}^\nu (x)[L_n^\nu (y)c_{n,\nu }J_\nu (2\sqrt{Nx})]$$ $$+[c_{n1,\nu }J_\nu (2\sqrt{N^{}x})L_{n1}^\nu (x)][L_n^\nu (y)c_{n,\nu }J_\nu (2\sqrt{Nx})].$$ The contribution of the right side to the difference of the two kernels is a constant which is $`O(1)`$ times $$N_0^1([L_{n1}^\nu (tx)c_{n1,\nu }J_\nu (2\sqrt{N^{}tx})]L_n^\nu (ty)$$ $$+L_{n1}^\nu (tx)[L_n^\nu (ty)c_{n,\nu }J_\nu (2\sqrt{Ntx})]$$ $$+[c_{n1,\nu }J_\nu (2\sqrt{N^{}tx})L_{n1}^\nu (tx)][L_n^\nu (ty)c_{n,\nu }J_\nu (2\sqrt{Ntx})])V(x)V(y)dt.$$ The integrand is a sum of three terms, each of which is the kernel (in the $`x,y`$ variables) of a rank one operator. The trace norm of such an integral is at most the integral of the Hilbert-Schmidt norms.) Using this fact and Lemma 3.2 we see that the contribution of the error term in (6) to any of these terms is $`O_1(N^{1/2})`$, as is the contribution of the last term above. Then we see that replacing the two Laguerre functions by the corresponding Bessel functions leads to an even smaller error. The error term in Lemma 3.3 is seen also to contribute $`O_1(N^{1/2})`$. Applying this lemma, and then doing everything with $`x`$ and $`y`$ interchanged we arrive at the statement of part (i). For part (ii), we have to show that the integrals are $`O_2(N^{1/2})`$. It is easy to see that with error $`O_1(N^{1/2})`$ we may replace the Bessel functions by their first-order asymptotics, resulting in a constant which is $`O(N^{1/4})`$ times $$V(x)V(y)x^{5/4}y^{1/4}_0^1\mathrm{sin}(2\sqrt{N^{}tx}\alpha )\mathrm{cos}(2\sqrt{Nty}\alpha )t𝑑t$$ plus similar expressions. A trigonometric identity and integration by parts shows that the integral is at most a constant times $$(\mathrm{max}(1,\|2\sqrt{N^{}x}2\sqrt{Ny}\|))^1+(\mathrm{max}(1,\|2\sqrt{N^{}x}+2\sqrt{Ny}\|2\alpha ))^1,$$ and an easy exercise shows that, together with the outer factors, this gives $`O_2(N^{1/4})`$. Analogous argument applies to the other integrals and this completes the proof of the lemma. ###### Lemma 3.5 The kernel (4) is equal to $$nV(x)V(y)_0^1J_\nu (2\sqrt{ntx})J_\nu (2\sqrt{nty})𝑑t$$ (8) plus an error $`o_2(1)`$. Proof. Since the Laguerre kernel has Hilbert-Schmidt norm $`n^{1/2}`$ (the operator is a rank $`n`$ projection) multiplying (4) by constants which are $`1+O(n^1)`$ produces an error $`O_2(n^{1/2})`$. The constants we choose are $`n/(\sqrt{n(n+\nu )}c_{n1,\nu }c_{n,\nu })`$. It follows from this and Lemma 3.4 that with error $`o_2(1)`$ we can replace (4) by $$\frac{n}{2}V(x)V(y)_0^1[J_\nu (2\sqrt{N^{}tx})J_\nu (2\sqrt{Nty})+J_\nu (2\sqrt{N^{}ty})J_\nu (2\sqrt{Ntx})]𝑑t.$$ (9) Now we show that if we replace $`N`$ and $`N^{}`$ by $`n`$ in this kernel the error is $`o_2(1)`$. Let’s look at the error incurred in the integral involving the first summand when we replace $`N^{}`$ by $`n`$. It equals $$_N^{}^n\frac{dr}{\sqrt{r}}_0^1\sqrt{tx}J_\nu ^{}(2\sqrt{rtx})J_\nu (2\sqrt{Nty})𝑑t.$$ (10) If $`ntx<1`$ and $`nty>1`$ the inner integral is at most a constant times $$n^{\mu /21/2}x^{\mu /2}_0^{\mathrm{min}(1,\mathrm{\hspace{0.17em}1}/nx)}t^{\mu /2}𝑑t,$$ where $`\mu =\mathrm{max}(\nu ,\mathrm{\hspace{0.17em}0})`$ as before. This is bounded by a constant times $$n^{1/2}(nx)^{\mu /2}\mathrm{if}nx<1,n^{3/2}x^1\mathrm{if}nx>1.$$ This times $`V(x)V(y)`$ is $`O_2(n^1)`$ and so its eventual contribution to the Hilbert-Schmidt norm (because of the external factor $`n`$ and the fact that $`rn`$) is $`O(n^{1/2})`$. If $`ntx<1`$ and $`nty<1`$ the inner integral is at most a constant times $$n^{\mu 1/2}x^{\mu /2}y^{\mu /2}_0^{\mathrm{min}(1,\mathrm{\hspace{0.17em}1}/nx,\mathrm{\hspace{0.17em}1}/ny)}t^\mu 𝑑t.$$ By symmetry we may assume $`y<x`$. If $`nx<1`$ this is at most a constant times the outer factor which, when multiplied by $`V(x)V(y)`$, is $`O_2(n^{3/2})`$. If $`nx>1`$ the above is at most a constant times $$n^{3/2+\mu /2}x^{1+\mu /2}y^{\mu /2},$$ and this times $`V(x)V(y)`$ is $`O_2(n^1)`$. Thus the eventual contribution of the portion of the integral where $`ntx<1`$ is $`O_2(n^{1/2})`$. In the region $`ntx>1`$, if we replace the first Bessel function in (10) by the first term of its asymptotic expansion it is easy to see we incur in the end an error $`O_2(n^{1/4})`$. After this replacement the inner integral becomes a constant times $$n^{1/4}x^{1/4}_{1/nx}^1t^{1/4}\mathrm{cos}(2\sqrt{rtx}\alpha )J_\nu (2\sqrt{Nty})𝑑t$$ $$=n^{3/2}x^{1/4}_{1/x}^nt^{1/4}\mathrm{cos}(2\sqrt{rtx/n}\alpha )J_\nu (2\sqrt{ty})𝑑t.$$ Taking account of the external factor of $`n`$ in our kernel, and the $`r`$-integral, we see that we want to show that the Hilbert-Schmidt norm of the kernel $$V(x)V(y)n^1_{1/x}^nt^{1/4}\mathrm{cos}(2\sqrt{rtx/n}\alpha )J_\nu (2\sqrt{ty})𝑑t$$ tends to 0 as $`n\mathrm{}`$. But from the asymptotics of the Bessel function it is clear that $`n^1`$ times the integral is uniformly bounded by a constant times $`1+y^{\mu /2}`$ and tends to 0 whenever $`xy`$. Hence the dominated convergence theorem tells us that the product is $`o_2(1)`$. An analogous argument shows that replacing $`N^{}`$ by $`n`$ in the second summand of (9) and then replacing $`N`$ by $`n`$ in both summands leads to an error $`o_2(1)`$. This completes the proof. Remark. In the preceding lemmas our various kernels had the factor $`V(x)V(y)`$. It is easy to see from their proofs that the lemmas hold with this factor replaced everywhere by $`V_1(x)V_2(y)`$ as long as $`V_1`$ and $`V_2`$ are bounded and have sufficiently rapid decay at infinity. The next lemma is the first that will require some smoothness. We shall denote by $`J_n(x,y)`$ the kernel $`(\text{5})`$ without the external factor $`V(x)V(y)`$, and by $`J_n`$ the corresponding operator. As before we denote by $`M_V`$ multiplication by $`V`$ so that (5) is the kernel of the operator $`M_VJ_nM_V`$. ###### Lemma 3.6 For any Schwartz function $`W`$ the commutator $`[W,J_n]`$ has Hilbert-Schmidt norm which is bounded as $`n\mathrm{}`$. Proof. Write the kernel of the commutator as in formula (3). It has the form $$\frac{W(x)W(y)}{xy}(J_\nu (2\sqrt{nx})\sqrt{ny}J_\nu ^{^{}}(2\sqrt{ny}))$$ (11) plus a similar term with $`x`$ an $`y`$ interchanged. If $`nx>1`$ and $`ny>1`$ then the product $$J_\nu (2\sqrt{nx})\sqrt{ny}J_\nu ^{^{}}(2\sqrt{ny})$$ is $`O(x^{1/4}y^{1/4}),`$ and thus if $`\|xy\|`$ is bounded away from zero, or if we integrate over any bounded region the Hilbert-Schmidt norm (the square root of the integral of the square) is bounded. If $`\|yx\|<1,`$ and assuming $`x,y>1`$ we see that $`y/x`$ is bounded and thus the Hilbert-Schmidt norm can be estimated by the square root of $$_1^{\mathrm{}}_1^{\mathrm{}}\|\frac{W(x)W(y)}{xy}\|^2𝑑x𝑑y.$$ But this is known to be bounded by $`_{\mathrm{}}^{\mathrm{}}\|x\|\widehat{W}(x)\|^2dx.`$ Now suppose that $`nx<1`$ and $`ny>1.`$ Then the product of Bessel functions is $`O(n^{\mu +1/4}x^{\mu /2}y^{1/4}).`$ If $`\|xy\|`$ is bounded away from zero, then the resulting Hilbert-Schmidt norm is $`O(n^{1/4}).`$ If $`\|xy\|<1,`$ then it is also clear that the Hilbert-Schmidt norm is $`O(n^{1/4})`$. The other two cases $`xn>1,yn<1`$ and $`xn<1,yn<1`$ are handled in the same fashion and are left to the reader. ###### Lemma 3.7 For any bounded functions $`W_1`$ and $`W_2`$ we have $$\mathrm{tr}(A_nB_n)W_2J_nW_10$$ as $`n\mathrm{}`$. Proof. For convenience all kernels $`K(x,y)`$ in this proof will be replaced by their unitary equivalents $`2\sqrt{xy}K(x^2,y^2)`$. We denote by $`L_n(x,y)`$ this unitary equivalent of the Laguerre kernel (4) without the external $`V`$ factors and by $`J_n(x,y)`$ here the unitary equivalent of the Bessel kernel (5) without the $`V`$ factors. We also make the substitution $`tt^2`$ in the $`t`$ integrals. Thus in our present notation $$J_n(x,y)=4n\sqrt{xy}_0^1J_\nu (2\sqrt{n}xt)J_\nu (2\sqrt{n}yt)t𝑑t.$$ (12) We also denote by $`\stackrel{~}{J}_n(x,y)`$ the unitary equivalent of the first displayed operator of Lemma 3.4, without the $`V`$ factors. Thus $$\stackrel{~}{J}_n(x,y)=2\sqrt{n(n+\nu )}c_{n1,\nu }c_{n,\nu }$$ $$\times \sqrt{xy}_0^1J_\nu [(2\sqrt{N^{}}tx)J_\nu (2\sqrt{N}ty)+J_\nu (2\sqrt{N^{}}ty)J_\nu (2\sqrt{N}tx)]t𝑑t.$$ (13) If we set $`V_i(x)=V(x^2)W_i(x^2)`$ then we see that our trace equals $`\mathrm{tr}M_{V_1}(L_nJ_n)M_{V_2}J_n`$. We shall show that this goes to 0 in two steps, showing first that $`\mathrm{tr}M_{V_1}(L_n\stackrel{~}{J}_n)M_{V_2}J_n0`$ and then that $`\mathrm{tr}M_{V_1}(\stackrel{~}{J}_nJ_n)M_{V_2}0`$. First, the asymptotics of the Bessel functions gives $$J_\nu (z)=\sqrt{\frac{2}{\pi z}}\mathrm{cos}(z\alpha )+O\left(\frac{z^{1/2}}{<z>}\right),$$ where $`<z>=(1+z^2)^{1/2}`$. (This also uses $`\nu 1/2`$.) Hence $$J_n(x,y)=\frac{8}{\pi }n\sqrt{xy}_0^1\left[\frac{\mathrm{cos}(2\sqrt{n}tx\alpha )}{(2\sqrt{n}tx)^{1/2}}+O\left(\frac{(\sqrt{n}tx)^{1/2}}{<(\sqrt{n}tx)>}\right)\right]$$ $$\times \left[\frac{\mathrm{cos}(2\sqrt{n}ty\alpha )}{(2\sqrt{n}ty)^{1/2}}+O\left(\frac{\sqrt{n}ty)^{1/2}}{<(\sqrt{n}ty)>}\right)\right]tdt$$ $$=\frac{8}{\pi }n_0^1\mathrm{cos}(2\sqrt{n}tx\alpha )\mathrm{cos}(2\sqrt{n}ty\alpha )𝑑t$$ $$+O\left(\frac{\sqrt{n}}{<\sqrt{n}x>}+\frac{\sqrt{n}}{<\sqrt{n}x>}+\frac{\sqrt{n}}{<\sqrt{n}x>^{1/2}<\sqrt{n}y>^{1/2}}\right).$$ The last summand is at most a constant times the sum of the preceding two. Using this, a trigonometric identity and integrating we find that $$J_n(x,y)=O\left(\frac{\sqrt{n}}{<\sqrt{n}(xy)>}+\frac{\sqrt{n}}{<\sqrt{n}(x+y)>}+\frac{\sqrt{n}}{<\sqrt{n}x>}+\frac{\sqrt{n}}{<\sqrt{n}y>}\right).$$ (14) We consider first $`M_{V_1}(L_n\stackrel{~}{J}_n)M_{V_2}J_n`$. Lemma 3.4 tells us that the kernel of $`M_{V_1}(L_n\stackrel{~}{J}_n)M_{V_2}`$ is $`O_1(n^{1/2})`$ plus the unitary equivalent of the expression in part (i) with modified $`V`$ factors. (See the remark following Lemma 3.5.) In the proof of part (ii) it was stated that if we replace the Bessel functions in this expression by their first order asymptotics the error is $`O_1(n^{1/2})`$. (We shall go through the details for similar integrals below.) So we may replace $`L_n\stackrel{~}{J}_n`$ by a constant which is $`O(1)`$ times $$_0^1[x^3\mathrm{sin}(2\sqrt{N^{}}tx\alpha )\mathrm{cos}(2\sqrt{N}ty\alpha )+y^3\mathrm{sin}(2\sqrt{N^{}}ty\alpha )\mathrm{cos}(2\sqrt{N}tx\alpha )]t^3𝑑t.$$ (Recall the unitary equivalents we are using and the variable change $`tt^2`$.) If in the integrals we made the replacements $`N,N^{}n`$ we would incur an error $`O((x^4+y^4)/\sqrt{n})`$. Multiplying this by $`V_1(x)V_2(y)`$ times (14) and integrating is easily seen to give $`o(1)`$. After these replacements the integral becomes what may be written $$(x^3y^3)_0^1\mathrm{sin}(2\sqrt{n}tx\alpha )\mathrm{cos}(2\sqrt{n}ty\alpha )t^3𝑑t+y^3_0^1\mathrm{sin}(2\sqrt{n}t(x+y)2\alpha )t^3𝑑t$$ $$=O\left(\frac{x^3y^3}{<\sqrt{n}(xy)>}+\frac{x^3y^3}{<\sqrt{n}(x+y)>}+\frac{y^3}{<\sqrt{n}(x+y)>}\right)$$ $$=O\left(\frac{x^3y^3}{<\sqrt{n}(xy)>}+\frac{x^3+y^3}{<\sqrt{n}(x+y)>}\right).$$ Let us see why if we multiply this by $`V_1(x)V_2(y)`$ times (14) and integrate we get $`o(1)`$. First, $$\frac{x^3y^3}{<\sqrt{n}(xy>}\frac{\sqrt{n}}{<\sqrt{n}(xy)>}=O\left(\frac{x^2+y^2}{<\sqrt{n}(xy)>}\right)$$ goes to zero pointwise and this times $`V_1(x)V_2(y)`$ is bounded by a fixed $`L^1`$ function. Thus the integral of the product goes to zero. The term with $`<\sqrt{n}(x+y)>`$ instead of $`<\sqrt{n}(xy)>`$ is even smaller. Next consider $$\frac{x^3y^3}{<\sqrt{n}(xy)>}\frac{\sqrt{n}}{<\sqrt{n}x>}.$$ We may ignore the factor $`x^3y^3`$ since it may be incorporated into the $`V_i`$. After the substitutions $`xx/\sqrt{n},yy/\sqrt{n}`$ the integral in question becomes $$\frac{1}{\sqrt{n}}\frac{\|V_1(x/\sqrt{n})V_2(y/\sqrt{n})\|}{<xy><x>}𝑑y𝑑x.$$ Schwarz’s inequality shows that the $`y`$ integral is $`O(n^{1/4})`$ so our double integral is bounded by a constant times $$n^{1/4}_0^1𝑑x+n^{1/4}_1^{\mathrm{}}\frac{\|V_1(x/\sqrt{n})\|}{x}𝑑x=n^{1/4}+n^{1/4}_{1/\sqrt{n}}^{\mathrm{}}\frac{\|V_1(x)\|}{x}𝑑x=O(n^{1/4}\mathrm{log}n).$$ Again the term with $`<\sqrt{n}(x+y)>`$ instead of $`<\sqrt{n}(xy)>`$ is even smaller. Now we look at $`M_{V_1}(\stackrel{~}{J}_nJ_n)M_{V_2}J_n`$. To find bounds for $`\stackrel{~}{J}_n(x,y)J_n(x,y)`$ let us look first at the error incurred if in the first integral in (13) we replace $`N^{}`$ by $`n`$. The error in the integral together with the external factor $`\sqrt{xy}`$ equals $$_N^{}^n\frac{dr}{\sqrt{r}}_0^1x^{3/2}y^{1/2}J_\nu ^{}(2\sqrt{r}tx)J_\nu (2\sqrt{n}ty)t^2𝑑t.$$ Using the asymptotics of $`J_\nu (z)`$ and the fact that $$J_\nu ^{}(z)=\sqrt{\frac{2}{\pi z}}\mathrm{sin}(z\alpha )+O(z^{3/2}),$$ we can write the above as a constant times $$_N^{}^n\frac{dr}{\sqrt{r}}_0^1x^{3/2}y^{1/2}\left[\frac{\mathrm{sin}(2\sqrt{r}tx\alpha )}{(\sqrt{r}tx)^{1/2}}+O((\sqrt{n}tx)^{3/2})\right]$$ $$\times \left[\frac{\mathrm{cos}(2\sqrt{N}ty\alpha )}{(\sqrt{N}ty)^{1/2}}+O\left(\frac{(\sqrt{n}ty)^{1/2}}{<(\sqrt{n}ty)>}\right)\right]t^2dt.$$ We estimate the trace norm of $`V_1(x)V_2(y)`$ times this by taking the trace norm under the integral signs. Since the integrand is, for fixed $`r`$ and $`t`$, a function of $`x`$ times a function of $`y`$ its trace norm equals the product of the $`L^2`$ norms of its factors. In multiplying out we will have main terms and error terms and we must estimate norms of all products. Thus we compute (in each line there will be an integral corresponding to a main term and then and error term) $$\frac{x^3\|V_1(x)\|^2}{\sqrt{r}tx}𝑑x=O(n^{1/2}t^1),\frac{x^3\|V_1(x)\|^2}{(\sqrt{n}tx)^3}𝑑x=O(n^{3/2}t^3),$$ $$\frac{y\|V_2(y)\|^2}{\sqrt{N}ty}𝑑y=O(n^{1/2}t^1),\frac{y(\sqrt{N}ty)^1\|V_2(y)\|^2}{<\sqrt{N}ty>^2}𝑑y=O(n^1t^2).$$ Combining $`L^2`$ norms we see that the trace norm of the contribution to the integrand of all but the product of the main terms is $`O(n^{3/4}t^2)`$. Integrating over $`t`$ we are left with $`O(n^{3/4})`$ and integrating over $`r`$ gives $`O(n^{5/4})`$. If we combine this we the external factor in (13) which is $`O(n)`$ we are left with $`O(n^{1/4})`$. Since the operator norms of the $`J_n`$ are bounded the eventual contribution to the trace of the product will be $`O(n^{1/4})`$. Thus we are left with the main term, which is $$_N^{}^n(rN)^{1/4}\frac{dr}{\sqrt{r}}_0^1x\mathrm{sin}(2\sqrt{r}tx\alpha )\mathrm{cos}(2\sqrt{N}ty\alpha )t𝑑t.$$ If in this we replaced $`r`$ and $`N`$ by $`n`$ everywhere in the integrand the error would be $`O(n^{3/2})`$. If we multiply this by $`V_1(x)V_2(y)`$ and use the estimate (14) we find by dominated convergence that the product has trace tending to zero, even keeping in mind the extra factor $`O(n)`$ in (13). So we may make these replacements, which results in $$(nN^{})n^1_0^1x\mathrm{sin}(2\sqrt{n}tx\alpha )\mathrm{cos}(2\sqrt{n}ty\alpha )t𝑑t.$$ Now there is a second integral in $`\stackrel{~}{J}_n`$, which is obtained from the first by interchanging $`x`$ and $`y`$. Interchanging and adding gives what can be written $$(nN^{})n^1_0^1[(xy)\mathrm{sin}(2\sqrt{n}tx\alpha )\mathrm{cos}(2\sqrt{n}ty\alpha )+y\mathrm{sin}(2\sqrt{n}tx+2\sqrt{n}ty2\alpha )]t𝑑t$$ $$=O\left(\frac{n^{3/2}(xy)}{<\sqrt{n}(xy)>}+\frac{n^{3/2}(x+y)}{<\sqrt{n}(x+y>}\right).$$ If we multiply by $`V_1(x)V_2(y)`$ and use the estimate (14) we find again that the product has trace tending to zero, even keeping in mind the extra factor $`O(n)`$. Thus replacement of $`N^{}`$ by $`n`$ in (13) leads to an eventual error in the trace of $`o(1)`$. Similarly so does then the replacement of $`N`$ by $`n`$. Finally, $$\sqrt{n(n+\nu )}c_{n1,\nu }c_{n,\nu }=n(1+O(n^1)),$$ so the eventual error in the trace upon replacing the constant by $`n`$ is $`O(n^1)`$ times what is obtained by multiplying $`V_1(x)V_2(x)`$ by the square of (14) and integrating. Dominated convergence shows this also to be $`o(1)`$. This completes the proof of the lemma. ## 4 Completion of the proof Recall that $`A_n`$ is the Laguerre kernel (4) and $`B_n`$ is the compressed Bessel kernel (5), which we also denote in its operator version as $`M_VJ_nM_V`$. We shall show first that $$det(I+A_n)det(I+B_n)$$ as $`n\mathrm{}`$. (It is clear that $`A_n`$ is a finite rank operator and so it is trace class, and using the integral representation (8) for the compressed Bessel kernel and integrating over $`t`$ shows that $`B_n`$ is also trace class. Thus both determinants are defined.) This will follow from what we have already done once we know that the operators $`I+B_n`$ are uniformly invertible, which means that they are invertible for sufficiently large $`n`$ and the operator norms of their inverses are $`O(1)`$. ###### Lemma 4.1 The operators $`I+B_n`$ are uniformly invertible and $$(I+B_n)^1=IM_VJ_nM_{VU^1}+O_2(1).$$ (15) Proof. We replace the operator by its unitary equivalent $`M_{\stackrel{~}{V}}J_nM_{\stackrel{~}{V}}`$ where now $`J_n`$ is given by (12), or equivalently $$J_n(x,y)=\sqrt{xy}_0^{2\sqrt{n}}J_\nu (xt)J_\nu (yt)t𝑑t$$ and we set $`\stackrel{~}{V}(x)=V(x^2)`$. If we set $`H(x,y)=\sqrt{xy}J_\nu (xy)`$ with $`H`$ the corresponding operator (the Hankel transform), and denote now by $`P_n`$ multiplication by the characteristic function of $`(0,\mathrm{\hspace{0.17em}2}\sqrt{n})`$, then $`J_n=HP_nH`$ and so $`I+B_n`$ is unitarily equivalent to $`I+M_{\stackrel{~}{V}}HP_nHM_{\stackrel{~}{V}}`$. These operators will be uniformly invertible if $`I+P_nHM_{\stackrel{~}{V}^2}HP_n`$ are. Now it has recently been shown that the operator $`HM_{\stackrel{~}{V}^2}H`$ is of the form $`W(\stackrel{~}{V}^2)+K`$, where $`W(\stackrel{~}{V}^2)`$ denotes the Wiener-Hopf operator with symbol $`\stackrel{~}{V}(x)^2=V(x^2)^2`$ and $`K`$ is a compact operator on $`L^2(𝐑^+)`$. (Much less is needed for this than that $`\stackrel{~}{V}^2`$ be a Schwartz function.) Since $`1+V(x^2)^2=U(x^2)`$ is nonzero and, being even, has zero winding number it follows from general facts about truncations of Wiener-Hopf operators that the operators $`I+P_nW(\stackrel{~}{V}^2)P_n`$ are uniformly invertible. Then since $`K`$ is compact it follows that the $`I+P_n(W(\stackrel{~}{V}^2)+K)P_n`$ will be uniformly invertible if the limiting operator $`I+W(\stackrel{~}{V}^2)+K=I+HM_{\stackrel{~}{V}^2}H`$ is invertible. (For an exposition of the facts we used here see, for example, Chap.2 of .) However, since $`H^2=I`$ the inverse of $`I+HM_{\stackrel{~}{V}^2}H`$ is easily seen to be $`I+HWH`$ where $`W=(1+\stackrel{~}{V}^2)^1\stackrel{~}{V}^2`$. This establishes the first statement of the lemma. For the second statement we apply Lemma 3.6 and use the facts $`J_n^2=J_n`$ (which follows from $`H^2=I`$) and that $`V^2`$ is a Schwartz function to see that for any bounded function $`W`$ we have $$(I+B_n)(I+M_VJ_nM_W)=(I+M_VJ_nM_V)(I+M_VJ_nM_W)$$ $$=I+M_VJ_nM_{V+V^2W+W}+O_2(1)=I+M_VJ_nM_{V+UW}+O_2(1).$$ If we choose $`W=VU^1`$ and multiply both sides by $`(I+B_n)^1`$ we obtain the result. ###### Lemma 4.2 $`det(I+A_n)det(I+B_n)`$ as $`n\mathrm{}`$. Proof. If an operator $`C`$ is trace class then $$det(I+C)=\stackrel{~}{det}(I+C)e^{\text{tr }C}$$ where $`\stackrel{~}{det}`$ is the generalized determinant . (The generalized determinant is defined for any Hilbert-Schmidt operator.) Hence we can write $$\frac{det(I+A_n)}{det(I+B_n)}=det((I+A_n)(I+B_n)^1)$$ $$=det(I+(A_nB_n)(I+B_n)^1)=\stackrel{~}{det}(I+(A_nB_n)(I+B_n)^1)e^{\text{tr }(A_nB_n)(I+B_n)^1}.$$ It follows from Lemmas 3.4(ii) and 3.5 that $`A_nB_n0`$ in Hilbert-Schmidt norm, and therefore from the uniform invertibility of the $`I+B_n`$ that the same is true of $`(A_nB_n)(I+B_n)^1`$. Therefore from the continuity of the generalized determinant in Hilbert-Schmidt norm that we conclude that the generalized determinant above has limit 1. Thus it suffices to show that $`\text{tr }(A_nB_n)(I+B_n)^10`$ as $`n\mathrm{}.`$ Since $`A_nB_n0`$ in Hilbert-Schmidt the $`O_2(1)`$ term in (15) contributes $`o(1)`$ to the trace of the product. By Lemma 3.7 the term $`M_VJ_nM_{VU^1}`$ in (15) also contributes $`o(1)`$. That $`\text{tr }(A_nB_n)`$ itself is $`o(1)`$ follows easily from arguments already given—one can check that at each stage the traces of the error operators tend to zero. Finally, we can quote the main result of which gives the asymptotics of $`det(I+B_n)`$ or, more exactly the determinants of their unitary equivalents. The formula is $$det(I+B_n)\mathrm{exp}\{\frac{2n^{1/2}}{\pi }_0^{\mathrm{}}\mathrm{log}(U(x^2)dx\nu /2\mathrm{log}U(0)+\frac{1}{2\pi ^2}_0^{\mathrm{}}xS(x)^2dx\}$$ where $`S(x)=_0^{\mathrm{}}\mathrm{cos}(xy)\mathrm{log}(U(y^2)dy.`$ This gives ###### Theorem 4.3 Suppose $`U`$ is nowhere zero and $`U1`$ is a Schwartz function. Then (1) holds. As mentioned in the introduction this result was computed heuristically in using the Coulomb fluid approach. In the same paper the analogous result was also obtained for weights supported on the entire real line. These determinants involve Hermite polynomials. It is highly likely that the results here (and techniques) could also be extended to that case.
warning/0006/cond-mat0006462.html	ar5iv	text	# Uniaxial Lifshitz Point at 𝑂⁢(ϵ_𝐿²) ## I Introduction The Lifshitz point occurs in a variety of physical systems and has been extensively studied over the last twenty-five years . It appears in High-$`T_C`$ superconductivity , polymer physics , ferroelectric liquid crystals , etc.. In magnetic systems, the uniaxial Lifshitz critical behavior can be described by an axially next-nearest-neighbor Ising model (ANNNI). It consists of a spin-$`\frac{1}{2}`$ system on a cubic lattice ($`d=3`$) with nearest-neighbor ferromagnetic couplings and next-nearest-neighbor antiferromagnetic interactions along a single lattice axis . The competition gives origin to a modulated phase, in addition to the ferromagnetic and paramagnetic ones. In spite of having several modulated phases, it was shown recently that around the Lifshitz critical region, a simple field-theoretic setting can be defined for this ANNNI model . In general, the antiferromagnetic couplings can show up in $`m`$ directions. In that case, the system possesses the $`m`$-fold Lifshitz critical point. Here we are going to focus our attention in the uniaxial case ($`m=1`$), since some materials present this type of critical behavior. MnP was studied both theoretically and experimentally and displays the uniaxial behavior . Theoretical studies involving the uniaxial Lifshitz point have been put forth using analytical and numerical tools. Examples of the latter are high-temperature series expansion and Monte Carlo simulations . Conformal invariance calculations in $`d=2`$ (in the context of strongly anisotropic criticality) and $`ϵ`$-expansion techniques have been the main analytical tools available to dealing with this kind of system. From the field-theoretic point of view, the critical dimension of a scalar field theory describing the uniaxial Lifshitz critical behavior is found (by using the Ginzburg criterion) to be $`d_c=4.5`$. The expansion parameter is $`ϵ_L=4.5d`$, where $`d`$ is the space dimension of the system under consideration. (In the pure Ising model the expansion parameter is $`ϵ=4d`$). As is well known, the critical exponents for the Lifshitz point at one-loop approximation have the same dependence in $`ϵ_L`$ as those from the pure Ising model have in $`ϵ`$. Of particular importance is the effect of the mixing of the two momenta scales, i. e., along and perpendicular to the competing axis. One can choose a convenient symmetry point to fix the external momenta scale in the quartic and quadratic directions. Then, the one-loop integral contributing to the four-point function, needed to find out the fixed point, can be performed without any approximation. The choice which simplifies the referred integral is to set the external momenta scale in the quartic direction to zero. The solution of this integral yields a leading singularity and a regular term in $`ϵ_L`$. The leading singularity in $`ϵ_L`$ can be chosen to be the same as that obtained in $`ϵ`$ when solving the analogous one-loop integral for the Ising model (by absorbing a convenient angular factor in a redefinition of the coupling constant), even though the coefficients of the regular terms in $`ϵ_L`$ and $`ϵ`$ are slightly different and depend on $`m`$ in the Lifshitz case . Thus, although the multiplicative factors to be absorbed in the coupling constant are different in the two cases, the critical exponents have the same dependence on the expansion parameter. We can then ask ourselves if this happens to be true when we proceed to calculate the critical exponents in higher orders in the perturbative expansion. Going one step further to evaluate higher order integrals in the loop expansion would be highly desirable using the same line of reasoning. In this way, all the dependence in the external momenta is along the $`(d1)`$ directions, perpendicular to the competition axis. In the resulting $`ϵ_L`$-expansion, the leading singularities can be chosen equal to those coming from the theory without competition (see below). The nontrivial new features of this expansion around the usual quadratic field theory are the coefficients of the subleading singularities and of the regular terms, which are no longer rational numbers. This approach would allow to treat this system properly in a perturbative expansion at least for correlation functions along the perpendicular directions to the competition axis. A better comprehension of this procedure might shed light in the perturbative study of higher order derivative field theories. In this sense, the Lifshitz critical behavior seems to be the natural laboratory to study higher order field theories in a perturbative framework. We report on what we believe to be the first study of critical exponents at two-loop order for the uniaxial Lifshitz point. Using $`\lambda \varphi ^4`$ field theory and the expansion in powers of $`ϵ_L=4.5d`$ in the Lifshitz critical point, we give a detailed account of the calculation of the exponents $`\nu _{L2}`$ and $`\eta _{L2}`$, which by now can be viewed as a worked out example of a recent generalization obtained for anisotropic behaviors . In order to solve higher-loop integrals we introduce a constraint relating the loop momenta in internal and external subdiagrams along the competing axis. The results for these integrals are consistent with homogeneity of the Feynman integrals in the external quadratic momenta scale. The exponents $`\nu _{L2}`$ and $`\eta _{L2}`$ are associated with the directions perpendicular to the competition axis. (The exponents $`\nu _{L4}`$ and $`\eta _{L4}`$ associated with the competition axis are not going to be considered here and we shall analyse them in another work.) We then obtain the exponent $`\gamma _{L2}`$ via scaling relations. The paper is organized as follows. In section 2 we introduce the notation and calculate the relevant integrals for the determination of the critical exponents at two-loop level. We present the critical exponents $`\eta _{L2}`$, $`\nu _{L2}`$ and $`\gamma _{L2}`$ in section 3. In section 4 we discuss our results and compare the $`\gamma _{L2}`$ exponent with numerical estimates based on Monte Carlo methods and high-temperature series. ## II Calculation of higher loop integrals The Lifshitz critical behavior can be described using a modified $`\lambda \varphi ^4`$ field theory. The bare Lagrangian associated with the uniaxial critical behavior is given by: $$L=\frac{1}{2}\|_1^2\varphi _0\|^2+\frac{1}{2}\|_{(d1)}\varphi _0\|^2+\delta _0\frac{1}{2}\|_1\varphi _0\|^2+\frac{1}{2}t_0\varphi _0^2+\frac{1}{4!}\lambda _0\varphi _0^4.$$ (1) The competition is responsible for the appearance of the quartic term in the free propagator. The Lifshitz critical point is characterized by the values $`t_0=\delta _0=0`$. From now on, this is the case which interests us in this work. First, we are going to compute the renormalized coupling constant at the fixed point. In order to find out the factor that shall be absorbed in the coupling constant, we quickly review the one-loop contribution to the four point function . The relevant integral is: $$I_2=\frac{d^{d1}qdk}{\left((k+k^{^{}})^4+(q+p)^2\right)\left(k^4+q^2\right)}.$$ (2) The external momenta are $`k^{}`$ along the quartic (competing) direction and $`\stackrel{}{p}`$ along the $`(d1)`$ quadratic directions. We then choose a symmetry point which simplifies the integral at external momenta $`k^{}=0`$, $`p^2=1`$. Using Schwinger’s parameterization we get $`{\displaystyle \frac{d^{d1}qdk}{\left(k^4+(q+p)^2\right)\left(k^4+q^2\right)}}={\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑\alpha _1𝑑\alpha _2\left(2{\displaystyle _0^{\mathrm{}}}𝑑k\mathrm{exp}((\alpha _1+\alpha _2)k^4)\right)`$ (3) $`\times {\displaystyle }d^{d1}q\mathrm{exp}((\alpha _1+\alpha _2)q^22\alpha _2q.p\alpha _2p^2).`$ (4) The $`\stackrel{}{q}`$ integral is straightforward. It can be performed to give $`{\displaystyle }d^{d1}q\mathrm{exp}((\alpha _1+\alpha _2)q^22\alpha _2q.p\alpha _2p^2)`$ (5) $`={\displaystyle \frac{1}{2}}S_{d1}\mathrm{\Gamma }({\displaystyle \frac{d1}{2}})(\alpha _1+\alpha _2)^{\frac{d1}{2}}\mathrm{exp}({\displaystyle \frac{\alpha _1\alpha _2p^2}{\alpha _1+\alpha _2}}).`$ (6) The next step is to compute the $`k`$ integration, which is: $$2_0^{\mathrm{}}𝑑k\mathrm{exp}((\alpha _1+\alpha _2)k^4)=\frac{1}{2}(\alpha _1+\alpha _2)^{\frac{1}{4}}\mathrm{\Gamma }(\frac{1}{4}).$$ (7) Replacing equations (5), (7) into equation (3) together with the value $`p^2=1`$, one finds $`\left({\displaystyle \frac{d^{d1}qdk}{\left(k^4+(q+p)^2\right)\left(k^4+q^2\right)}}\right)_{p^2=1}={\displaystyle \frac{1}{4}}S_{d1}\mathrm{\Gamma }({\displaystyle \frac{d1}{2}})\mathrm{\Gamma }({\displaystyle \frac{1}{4}})`$ (8) $`\times {\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}d\alpha _1d\alpha _2\mathrm{exp}({\displaystyle \frac{\alpha _1\alpha _2}{\alpha _1+\alpha _2}})(\alpha _1+\alpha _2)^{(\frac{d1}{2}+\frac{1}{4})}.`$ (9) We can perform one of the integrals in the Schwinger parameters using a change of variables. Set $`v=\frac{\alpha _2}{\alpha _1+\alpha _2}`$, and $`u=\alpha _1v`$. The integral over $`u`$ can be done, and we are left with $`{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑\alpha _1𝑑\alpha _2\mathrm{exp}({\displaystyle \frac{\alpha _1\alpha _2}{\alpha _1+\alpha _2}})(\alpha _1+\alpha _2)^{(\frac{d1}{2}+\frac{1}{4})}`$ (10) $`=\mathrm{\Gamma }(2({\displaystyle \frac{d1}{2}}+{\displaystyle \frac{1}{4}})){\displaystyle _0^1}𝑑v(v(1v))^{(\frac{d1}{2}+\frac{1}{4})2}.`$ (11) Now we make the continuation $`d=4.5ϵ_L`$. One obtains a result in terms of Gamma functions with non integer arguments. A useful identity involving the expansion of Gamma functions around a small number is given by: $$\mathrm{\Gamma }(a+bx)=\mathrm{\Gamma }(a)\left[\mathrm{\hspace{0.17em}1}+bx\psi (a)+O(x^2)\right],$$ (12) where $`\psi (z)=\frac{d}{dz}ln\mathrm{\Gamma }(z)`$. This allows one to obtain the $`ϵ_L`$-expansion when the Gamma functions have non integer arguments. Replacing Eqs. (7), (8) into Eq. (6), we obtain: $$I_2=\frac{1}{2}\mathrm{\Gamma }(7/4)\mathrm{\Gamma }(1/4)S_{d1}\frac{1}{ϵ_L}(1+i_2ϵ_L),$$ (13) where $`i_2=1+\frac{1}{2}(\psi (1)\psi (\frac{7}{4}))`$. We absorb in the coupling constant a geometric angular factor, which is $`\frac{1}{2}\mathrm{\Gamma }(7/4)\mathrm{\Gamma }(1/4)S_{d1}`$, where $`S_d=[2^{d1}\pi ^{\frac{d}{2}}\mathrm{\Gamma }(\frac{d}{2})]^1`$ . Then the redefined integral is: $$\stackrel{~}{I}_2=\frac{I_2}{\frac{1}{2}\mathrm{\Gamma }(7/4)\mathrm{\Gamma }(1/4)S_{d1}},$$ (14) or $$\stackrel{~}{I}_2=\frac{1}{ϵ_L}(1+i_2ϵ_L).$$ (15) We suppress the tilde hereafter to simplify the notation. We have to keep in mind that we should divide out this factor for each loop integration. We now turn our attention to higher-loop integrals. In practice, we have to calculate the two-loop integrals $`I_{4SP}I_4`$, $`\frac{}{p^2}I_3\|_{SP}I_3^{}`$ and the three-loop integral $`\frac{}{p^2}I_5\|_{SP}I_5^{}`$, in order to find the fixed point at two-loop level and the critical exponents. The subscript $`SP`$ is used to denote our choice of the subtraction point. They are given by (see Figure 1): $$I_3=\frac{d^{d1}q_1d^{d1}q_2dk_1dk_2}{\left(q_1^2+k_1^4\right)\left(q_2^2+k_2^4\right)\left((q_1+q_2+p)^2+(k_1+k_2+k^{})^4\right)},$$ (16) $`I_5=`$ $`{\displaystyle \frac{d^{d1}q_1d^{d1}q_2d^{d1}q_3dk_1dk_2dk_3}{\left(q_1^2+k_1^4\right)\left(q_2^2+k_2^4\right)\left(q_3^2+k_3^4\right)\left((q_1+q_2p)^2+(k_1+k_2k^{})^4\right)}}`$ (18) $`\times {\displaystyle \frac{1}{(q_1+q_3p)^2+(k_1+k_3k^{})^4}}.`$ $`I_4=`$ $`{\displaystyle \frac{d^{d1}q_1d^{d1}q_2dk_1dk_2}{\left(q_1^2+k_1^4\right)\left((Pq_1)^2+(K^{}k_1)^4\right)\left(q_2^2+k_2^4\right)}}`$ (20) $`\times {\displaystyle \frac{1}{(q_1q_2+p_3)^2+(k_1k_2+k_3^{})^4}}.`$ In the first two integrals, $`\stackrel{}{p}`$ is the external momentum (associated with the two-point vertex) along $`(d1)`$ directions, whereas $`k^{}`$ is the external momentum along the competition axis. Inside the integral $`I_4`$, $`P=p_1+p_2`$, with $`p_1,p_2,p_3`$ being the external momenta (associated with the four-point vertex) along the quadratic directions, and $`K^{}=k_1^{}+k_2^{}`$, with $`k_1^{},k_2^{},k_3^{}`$ the external momenta along the quartic direction. The symmetry point is chosen as follows. We set all the external momenta at the competition axis equal to zero. For the four-point vertex, the external momenta along the quadratic directions are chosen as $`p_i.p_j=\frac{\kappa ^2}{4}(4\delta _{ij}1)`$. We fix the momentum scale of the two-point function through $`p^2=\kappa ^2=1`$. Now we can study the solution of the higher-loop integrals shown above. Consider the integral $`I_3`$. With our choice for the quartic external momenta it is given by: $$I_3=\frac{d^{d1}q_1dk_1}{(q_1^2+k_1^4)}\frac{d^{d1}q_2dk_2}{\left(q_2^2+k_2^4\right)\left((q_1+q_2+p)^2+(k_1+k_2)^4\right)}.$$ (21) First, we perform the integral over the internal bubble, i.e., we integrate over $`q_2`$ and $`k_2`$. In order to solve the internal bubble we demand that the loop momenta $`k_1`$ should be related to $`k_2`$. Note that we could have chosen the other way around, since the integral is symmetric under the exchange $`k_1k_2`$, $`q_1q_2`$. There are two issues which need to be emphasized here. First, the Lifshitz point condition eliminates the quadratic part of the momenta along the competition axis, for $`\delta _0=0`$. Second, the remaining quartic part of the integral mixes the two loop momenta ($`k_1,k_2`$) in different subdiagrams, yielding crossed terms which are difficult to integrate. Indeed, using Schwinger parameters and carrying out the integration over $`q_2`$ first, we obtain $`I_3(p,0)={\displaystyle \frac{1}{2}}S_{dm}\mathrm{\Gamma }({\displaystyle \frac{d1}{2}}){\displaystyle \frac{d^{d1}q_1dk_1}{q_1^2+(k_1^2)^2}}`$ (22) $`\times {\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}d\alpha _1d\alpha _2(\alpha _1+\alpha _2)^{\frac{(d1)}{2}}exp({\displaystyle \frac{\alpha _1\alpha _2}{\alpha _1+\alpha _2}}(q_1+p)^2){\displaystyle }dk_2e^{\alpha _1(k_2^2)^2}e^{\alpha _2((k_1+k_2)^2)^2}.`$ (23) In order to integrater over $`k_2`$, we have to expand the argument of the last exponential. This results in a complicated function of $`\alpha _1,\alpha _2,k_1`$ and $`k_2`$, which has no elementary primitive. Considering the remaining terms as a damping factor to the integrand, the maximum of the integrand will be either at $`k_1=0`$ or at $`k_1=2k_2`$. (The most general choice $`k_1=\alpha k_2`$ yields a hypergeometric function.) The constraint eliminates the crossed terms and is the simplest way to disentangle the two quartic integrals in the loop momenta. The choice $`k_1=2k_2`$ implies that $`k_1`$ varies in the internal bubble, but not in an arbitrary manner. Its variation is dominated by $`k_2`$ through this constraint, which eliminates the dependence on $`k_1`$ in the internal subdiagram. Integrating over $`k_2`$ yields a simple factor to the remaining parametric integrals (over the two Schwinger parameters) proportional to $`(\alpha _1+\alpha _2)^{\frac{1}{4}}`$. After solving the parametric integrals we find: $$I_3=\frac{d^{d1}q_1dk_1}{\left(q_1^2+k_1^4\right)[(q_1+p)^2]^{\frac{ϵ_L}{2}}}I_2.$$ (24) We use Schwinger’s parameterization again to solve this integral along the quartic direction. We obtain: $$I_3=I_2\frac{d^{d1}q_1}{[q_1^2]^{\frac{3}{4}}[(q_1+p)^2]^{\frac{ϵ_L}{2}}}.$$ (25) The difference with respect to the pure $`\varphi ^4`$ theory is that after performing the quartic integral, we get a exponent for the quadratic part of the momenta which is not an integer. At this point, one can use Feynman parameters to solve the momentum integrals. The dependence of the integral in the external momenta is proportional to $`(p^2)^{1ϵ_L}`$, in conformity with the homogeneity of the Feynman integrals in the external momenta scale. Deriving with respect to $`p^2`$ and setting $`p^2=1`$, we find : $$I_3^{}=\frac{1}{7ϵ_L}\left[1+\left(i_2+\frac{6}{7}\right)ϵ_L\right].$$ (26) The integral $`I_4`$ can be calculated in a similar fashion. First, we set the external momenta along the competing direction equal to zero. Therefore $`I_4=`$ $`{\displaystyle \frac{d^{d1}q_1dk_1}{\left(q_1^2+k_1^4\right)\left((Pq_1)^2+k_1^4\right)}}`$ (29) $`\times {\displaystyle }{\displaystyle \frac{d^{d1}q_2dk_2}{\left(q_2^2+k_2^4\right)[(q_1q_2+p_3)^2+(k_1+k_2)^4]}},`$ where we changed variables from $`k_2k_2`$. We set $`k_1=2k_2`$ in the internal bubble $`q_2,k_2`$ (as we did for $`I_3`$), and integrate over $`q_2,k_2`$. We then have $$I_4=I_2\frac{d^{d1}q_1dk_1}{\left(q_1^2+k_1^4\right)\left((Pq_1)^2+k_1^4\right)}\frac{1}{[(q_1+p_3)^2]^{\frac{ϵ_L}{2}}}.$$ (30) We use Schwinger’s parameterization to get rid of the $`k_1`$ integral. After performing the change of variables used to calculate the one-loop integral and a rescaling, we can solve one of the parametric integrals to get $$I_4=I_2_0^1𝑑z\frac{d^{d1}q_1}{(q_1^22zP.q_1+zP^2)^{\frac{7}{4}}[(q_1+p_3)^2]^{\frac{ϵ_L}{2}}}.$$ (31) In order to perform the integral over $`q_1`$, we make use of a Feynman parameter obtaining $`I_4=`$ $`{\displaystyle \frac{1}{2}}I_2\left(1{\displaystyle \frac{ϵ_L}{2}}\psi \left({\displaystyle \frac{7}{4}}\right)\right){\displaystyle \frac{\mathrm{\Gamma }(ϵ_L)}{\mathrm{\Gamma }\left(\frac{ϵ_L}{2}\right)}}{\displaystyle _0^1}𝑑yy^{\frac{3}{4}}(1y)^{\frac{1}{2}ϵ_L1}`$ (33) $`\times {\displaystyle _0^1}dz[yz(1yz)P^2+y(1y)p_3^22yz(1y)p_3.P]^{ϵ_L}.`$ There is a subtlety that needs to be analyzed with care. Here we proceed in complete analogy to the pure $`\varphi ^4`$ theory . The integral over $`y`$ is singular at $`y=1`$ when $`ϵ_L=0`$. We add and subtract the value of the integrand in the last integral at the point $`y=1`$ $`[yz(1yz)P^2+y(1y)p_3^22yz(1y)p_3.P]^{ϵ_L}=\left[z(1z)P^2\right]^{ϵ_L}`$ (34) $`ϵ_L\mathrm{ln}\left\{{\displaystyle \frac{\left[yz(1yz)P^2+y(1y)p_3^22yz(1y)p_3.P\right]}{z(1z)P^2}}\right\}+O(ϵ_L^2).`$ (35) As $`y1`$ the logarithm is zero when $`ϵ_L=0`$, leading to a well defined result for the integral over $`y`$. The coefficient of the integral is proportional to $`\frac{1}{ϵ_L}`$, which cancels the $`ϵ_L`$ in front of the logarithm. The logarithm contributes only to the order $`ϵ_L^0`$ and can be neglected. We then find $$I_4=\frac{1}{2ϵ_L^2}\left[1+3i_2ϵ_L\right].$$ (36) Finally, let us describe the computation of the three-loop integral $`I_5`$. At zero external momenta along the competition axis, this integral reads: $`I_5=`$ $`{\displaystyle \frac{d^{d1}q_1dk_1}{\left(q_1^2+k_1^4\right)}\frac{d^{d1}q_2dk_2}{\left(q_2^2+k_2^4\right)\left((q_1+q_2p)^2+(k_1+k_2)^4\right)}}`$ (40) $`\times {\displaystyle }{\displaystyle \frac{d^{d1}q_3dk_3}{\left(q_3^2+k_3^4\right)((q_1+q_3p)^2+(k_1+k_3)^4)}}.`$ The integral is symmetric in $`q_2q_3`$, $`k_2k_3`$. As the loop momenta are dummy variables, the two internal bubbles are really the same. That is why we do not need more than one relation among the loop momenta in the external and internal bubbles, even though this is a three-loop diagram (with two internal bubbles). As a matter of fact, we can solve the two integrals independently in the following way. In the internal bubble $`q_3,k_3`$ we set $`k_1=2k_3`$, as well as $`k_1=2k_2`$ over the other internal bubble $`q_2,k_2`$. Apparently we have two different relations among the loop momenta, but one of them is artificial. This means that the two internal bubbles give the same contribution, i.e. the integration over the internal bubbles is proportional to $`I_2^2`$. Thus $$I_5=I_2^2\frac{d^{d1}q_1dk_1}{\left(q_1^2+k_1^4\right)[(q_1+p)^2]^{ϵ_L}}.$$ (41) We integrate first over $`k_1`$ and proceed analogously as before, to find the following result: $$I_5^{}=\frac{4}{21ϵ_L^2}\left[1+\left(2i_2+\frac{8}{7}\right)ϵ_L\right].$$ (42) We stress once again that the constraint preserves the physical principle of homogeneity in all the higher-loop Feynman integrals in the external momenta scale, which is consistent with scaling theory. With the integrals calculated in this way, we can find out the exponents as is going to be shown in the next section. ## III Critical exponents To compute the critical exponents associated to the ferromagnetic planes at the Lifshitz critical point, one may use the standard field-theoretic approach . This is possible since no new renormalization constants need to be introduced in this caseThis is not valid in the calculation of the critical exponents along the competition axis.. From the results for $`I_3^{}`$, Eq. (26), and $`I_5^{}`$, Eq. (42), we note that these integrals do not have the same leading singularities as in pure $`\varphi ^4`$ theory. As an approximation, we introduce a weight factor for the two point vertex function, in order to identify the leading singularities with the ones appearing in the pure $`\varphi ^4`$ field theory. This factor is $`\frac{7}{8}`$ for the integrals above. This approximation is suitable when one considers the generalization for the $`m`$-fold case with $`m8`$. In this way, we have a smooth transition from the Isinglike case ($`m=0`$) to the general Lifshitz anisotropic critical behavior ($`m8`$) . The bare and renormalized quantities are related through $`\varphi _0=Z_\varphi ^{1/2}\varphi `$ and $`u_0=Z_\varphi ^2Z_uu`$, where $`\varphi _0`$ and $`\lambda _0\kappa ^{ϵ_L}u_0`$ are the bare parameters in Eq. (1). As usual, $`Z_\varphi 1+\delta _\varphi `$ and $`Z_u1+\delta _u`$ are the wave-function and coupling constant renormalization constants, respectively. In addition, we introduce the composite field renormalization constant $`Z_{\varphi ^2}`$, and also $`\overline{Z}_{\varphi ^2}=Z_{\varphi ^2}Z_\varphi 1+\overline{\delta }_{\varphi ^2}`$. We have $`\beta _u=ϵ_L\left({\displaystyle \frac{\mathrm{ln}u_0}{u}}\right)`$ (43) $`\gamma _\varphi =\beta _u{\displaystyle \frac{\mathrm{ln}Z_\varphi }{u}}`$ (44) $`\overline{\gamma }_{\varphi ^2}=\beta _u{\displaystyle \frac{\mathrm{ln}\overline{Z}_{\varphi ^2}}{u}}`$ (45) The fixed point $`u^{}`$ is given by the solution of the equation $`\beta _u(u^{})=0`$, and the critical exponents by the relations $`\eta _{L2}=\gamma _\varphi (u^{})`$ and $`\nu _{L2}^1=2\eta _{L2}\overline{\gamma }_\varphi ^2(u^{})`$. In case the order parameter has a $`O(N)`$ symmetry, the formulas relating the integrals computed in section 2 and the renormalization constants defined above are: $`\delta _\varphi =B_2u^2I_3^{}+\left(2B_3I_2I_3^{}B_3I_5^{}\right)u^3+O(u^4)`$ (46) (47) $`\delta _u=3A_1uI_2+3\left(6A_1^2I_2^2A_2^{(1)}I_2^22A_2^{(2)}I_4\right)u^2+O(u^3)`$ (48) (49) $`\overline{\delta }_{\varphi ^2}=C_1uI_2+\left(C_2I_2^2C_1I_4\right)u^2+O(u^3)`$ (50) where $`A_1=(N+8)/18`$, $`A_2^{(1)}=(N^2+6N+20)/108`$, $`A_2^{(2)}=(5N+22)/54`$, $`B_2=(N+2)/18`$, $`B_3=(N+2)(N+8)/108`$, $`C_1=(N+2)/6`$, and $`C_2=(N+2)(N+8)/36`$. With this information we compute the fixed point at two-loop level. We expand the dimensionless bare coupling constant $`u_0`$ in terms of the renormalized coupling $`u`$ and the $`ϵ_L`$ parameter. Using Eqs. (43) and (46), we find the fixed point for the $`O(N)`$ symmetric theory in the following form: $$u^{}=\frac{6}{8+N}ϵ_L\left\{1+ϵ_L\left[\left(\frac{4(5N+22)}{(8+N)^2}1\right)i_2\frac{(2+N)}{(8+N)^2}\right]\right\}.$$ (51) Therefore, the exponents $`\eta _{L2}`$ and $`\nu _{L2}`$ are given by: $`\eta _{L2}={\displaystyle \frac{1}{2}}ϵ_L^2{\displaystyle \frac{2+N}{(8+N)^2}}`$ (52) (53) $`+ϵ_L^3{\displaystyle \frac{(2+N)}{(8+N)^2}}\left[\left({\displaystyle \frac{4(5N+22)}{(8+N)^2}}{\displaystyle \frac{1}{2}}\right)i_2+{\displaystyle \frac{1}{7}}{\displaystyle \frac{(2+N)}{(8+N)^2}}\right].`$ (54) $`\nu _{L2}={\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{4}}ϵ_L{\displaystyle \frac{2+N}{8+N}}`$ (55) (56) $`+{\displaystyle \frac{1}{8}}{\displaystyle \frac{(2+N)}{(8+N)^3}}\left[2(14N+40)i_22(2+N)+(8+N)(3+N)\right]ϵ_L^2.`$ (57) Now using Fisher’s law along directions perpendicular to the competing axis, namely $`\gamma _{L2}=\nu _{L2}(2\eta _{L2})`$, the exponent $`\gamma _{L2}`$ can be written as: $`\gamma _{L2}=1+{\displaystyle \frac{1}{2}}ϵ_L{\displaystyle \frac{2+N}{8+N}}`$ (58) (59) $`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{(2+N)}{(8+N)^3}}\left[12+8N+N^2+4i_2(20+7N)\right]ϵ_L^2.`$ (60) Previous results in the literature only yielded the exponent $`\eta _{L2}`$ at $`O(ϵ_L^2)`$ and the exponent $`\nu _{L2}`$ at $`O(ϵ_L)`$. For this uniaxial case, our results express the critical exponents in a higher order in $`ϵ_L`$ compared to earlier investigations. A detailed comparison with other methods is provided in the next section. ## IV Discussion First of all, our result for the exponent $`\eta _{L2}`$ is in agreement with Mukamel’s calculation at $`O(ϵ_L^2)`$. Therefore, our method is equivalent to integrating over the momentum shell as was done in his work using the Landau-Ginzburg-Wilson Hamiltonian approach. For the ANNNI model, $`\gamma _{L2}=1.4\pm 0.06`$ is the former Monte Carlo output , whereas the best estimates from the high-temperature series is $`\gamma _{L2}=1.62\pm 0.12`$ . Note that we use the subscript $`\gamma _{L2}`$ instead of $`\gamma _L`$, since it was shown recently that the exponents parallel and perpendicular to the competition axis obey independent scaling laws . Our two-loop calculation obtained from the $`ϵ_L`$-expansion via the scaling law (when neglecting $`O(ϵ_L^3)`$ terms) in three dimensions yields $`\gamma _{L2}=1.45`$. This agrees (within the error bar) with the former Monte Carlo result, the difference being of order $`10^2`$. Nevertheless, the most recent high-precision numerical Monte Carlo estimate for the ANNNI model yielded $`\gamma _{L2}=1.36\pm 0.03`$ . In order to figure out how to extract the best numerical results from the $`ϵ_L`$-expansion when the $`ϵ_L`$ parameter is not small (which is the case for $`d=3`$), a comparison with the Ising model is worthwhile. For the exponent $`\gamma `$ in three dimensions, the $`ϵ`$-expansion gives a contribution of 0.167 at $`O(ϵ)`$ and 0.077 at $`O(ϵ^2)`$ . The $`O(ϵ)`$ contributes with $`13\%`$ and the order $`O(ϵ^2)`$ with $`6\%`$ to the value of $`\gamma `$ ($`1.24`$), respectively. For the uniaxial Lifshitz case, the contributions for the $`\gamma _{L2}`$ index are $`0.25`$ ($`17\%`$) and $`0.196`$ ($`14\%`$). The very close values of the contributions of first and second order to $`\gamma _{L2}`$ (as the $`ϵ_L`$ parameter is 1.5 not being a small number), indicates that neglecting $`O(ϵ_L^3)`$ could be a dangerous step in obtaining the exponent $`\gamma _{L2}`$ via scaling relations in a more accurate way. Indeed, had we replaced the numerical values obtained for $`\nu _{L2}=0.73`$, $`\eta _{L2}=0.04`$ and $`d=3`$ directly into the scaling law, we would have obtained $`\gamma _{L2}=1.43`$. As argued in for the other critical exponents $`\alpha _{L2}`$ and $`\beta _{L2}`$, whenever $`ϵ_L>1`$ one should use the numerical values of $`\nu _{L2}`$, $`\eta _{L2}`$ obtained from the $`ϵ_L`$-expansion for fixed values of $`(N,d,m)`$ in order to obtain the numerical values of the other exponents via scaling laws. Therefore, provided we give this new interpretation to the numerical output of the $`ϵ_L`$-expansion when $`ϵ_L>1`$, we consider that the agreement between the numerical (Monte Carlo) and analytical ($`ϵ_L`$-expansion) results is remarkable for $`d=3`$. The numerical value obtained here for the correlation length exponent is $`\nu _{L2}=0.73`$. The experimental value of this critical index is still lacking. We hope our result sheds some light towards its experimental determination. The extension of the present method to the calculation of critical exponents for the $`m`$-fold ($`m8`$) case reduces to the Ising-like case when $`m=0`$ and to the present case when $`m=1`$ . An interesting open question is the calculation of the critical exponents $`\nu _{L4}`$ and $`\eta _{L4}`$ using the $`ϵ_L`$-expansion at two-loop level. The approach followed here is not suitable to computing these critical exponents (parallel to the competition axis), since our choice of the symmetry point prevents a proper treatment in this direction. The possibility of devising another symmetry point to deal with these exponents seems to be feasible, and will be reported elsewhere. In recent articles, some authors studied an alternative field-theoretic approach based on coordinate space calculations. In the first paper they recovered the results of reference for the cases $`m=2,6`$ analytically, but only could get the exponents numerically for the $`m=1`$ case, working entirely in coordinate space. It is worth emphasizing that they computed the fixed point only at one-loop order (see equation (82) in the mentioned paper). In the second paper, they computed the critical exponents at second order in perturbation theory by making use of a hybrid mechanism, going to coordinate or momentum space according to the necessity through a scaling function related to the free propagator in coordinate space. They obtained the exponents, whose coefficients of each power of $`ϵ_L`$ are integrals to be performed numerically. The very similar values obtained for the exponents using their method or ours confirms that momentum and coordinate space calculations should give the same results, since either our approximation or the numerical approximation made by them is responsible for a rather small deviation in the two results when compared to the above numerical values. In conclusion, we have found a way to perform two- and three-loop integrals for the uniaxial Lifshitz point, needed to calculate universal properties at directions perpendicular to the competition axis. The constraint in the loop momenta at the competition direction is the key ingredient to carry out the calculations. In this approximation, the loop momenta along the competition axis are not conserved when one uses this constraint. However the momenta along the $`(d1)`$ directions are conserved. Momentum non-conservation along the competing direction as a higher-order effect does not seem to affect the critical exponents considered here in a significant way, as indicated by the comparison of our study with the available numerical data for the $`d=3`$ case. It might be interesting to study general field theories including higher order derivative terms in this new framework. Topics including the extension of the present method to the region out of the Lifshitz point ($`\delta _00`$) and two-loop calculations using a modified symmetry point along the competition axis are in development. Acknowledgments The authors acknowledge support from FAPESP, grant numbers 00/03277-3 (LCA), 98/06612-6 and 00/06572-6 (MML), and Élcio Abdalla for kind hospitality at the Departamento de Física Matemática da Universidade de São Paulo. LCA would like to thank Marcelo Gomes and Adilson J. da Silva for useful discussions. MML would like to thank Warren Siegel, Peter van Nieuwenhuizen, Victor O. Rivelles and Nathan Berkovits for helpful conversations. He also acknowledges kind hospitality at C. N. Yang Institute for Theoretical Physics (SUNYSB) where this work has started. Figure captions * Figure 1. Feynman graphs corresponding to the integrals: (a) $`I_2`$;(b) $`I_3`$; (c) $`I_4`$; (d) $`I_5`$. The broken lines in the graphs (b), (c), and (d) define the “internal” bubbles in each case. The momenta $`q_i,k_i`$ refer to the loop momenta in the quadratic and quartic directions, respectively. The labels $`p_i,k_i^{}`$ denote the external momenta in the quadratic and quartic directions, respectively.
warning/0006/cond-mat0006446.html	ar5iv	text	# Spectral Properties of Coupled Bose-Einstein Condensates ## I introduction After the recent observation of Bose-Einstein condensation in dilute atomic gases realized by confining a macroscopic population of atoms in a potential trap , a large amount of work has been devoted to construct devices where two condensates are trapped in a double-well potential. The interaction of such coupled Bose-Einstein condensates (BEC) gives rise to quantum phenomena such as coherent tunneling and interference effects that have been the subject of a huge number of studies, both theoretic and experimental . The present paper is focussed on the dynamical aspects inherent in the interaction of two condensates trapped in two identical wells. Such a dynamics has been studied thoroughly in Refs. within the minimal interaction scheme $$\{\begin{array}{cc}& i\mathrm{}\dot{\psi }_1=\left[\frac{\mathrm{}^2}{2m}\mathrm{}v+U\|\psi _1\|^2\right]\psi _1T\psi _2\hfill \\ & \\ & i\mathrm{}\dot{\psi }_2=\left[\frac{\mathrm{}^2}{2m}\mathrm{}v+U\|\psi _2\|^2\right]\psi _2T\psi _1\hfill \end{array}$$ (1) where the classical fields $`\psi _j(𝐫,t)`$ obey two coupled Gross-Pitaevskii equations (GPE), and $`U`$, $`v`$, $`T`$, describe the interatomic scattering, the external potential and the tunneling amplitude, respectively. Fields $`\psi _j(𝐫)`$ (often called the wave functions of the condensate ) are defined as the expectation value $`\psi _j(𝐫,t)=\widehat{\mathrm{\Psi }}_j(𝐫,t)`$ of the field operators $`\widehat{\mathrm{\Psi }}_j`$ within the many-body quantum theory of BEC’s . The negligible space dependence of $`\psi _j`$ through the condensates allows one to set $`\mathrm{}\psi _j0`$ so that Eqs. (1) reduce to a hamiltonian system with two complex degrees of freedom where the nonlinear cubic terms provide the system with an ample variety of interwell processes. Before discussing the goals of this paper, it is useful to briefly review the relevant traits of the dynamics issued from the space-independent form of Eqs. (1). Due to the assumption that the bosonic wells are identical, the associated model Hamiltonian $`(\psi _1,\psi _2)=U(\|\psi _1\|^4+\|\psi _2\|^4)v𝒩T(\psi _1\psi _2^{}+\psi _2\psi _1^{})`$ ($`𝒩=\|\psi _1\|^2+\|\psi _2\|^2`$) exhibits a permutational symmetry (PS) realized by the exchange of the dynamical variables $`\psi _1`$, $`\psi _2`$. Physically, this is represented by the population exchange, $`\|\psi _j\|^2`$ being the boson number of the $`j`$-th well up to a volume factor. The structure of the model phase space $`𝒫`$ (this is two-dimensional because $`𝒩`$ is a constant of motion) reflects in a nontrivial way the presence of the PS. For energies $`E<E_{}`$, ($`E_{}`$ is a critical value of the energy depending on the model parameters) the orbits are placed concentrically around the energy minimum, and cover a region $`C_0𝒫`$ that has essentially the structure of the harmonic oscillator phase space. This can be shown to entail the oscillations of the two condensates’ populations around the common value $`𝒩/2`$. Also, for $`E<E_{}`$ the PS maps each orbit in itself. The remaining part of $`𝒫`$, filled by orbits with $`E>E_{}`$, is formed by two (spatially) disjoint components $`C_+`$, $`C_{}`$ such that $`𝒫`$ $`C_+C_{}C_0`$. In this case any given energy value $`E`$ is associated with two distinct orbits $`\gamma _+C_+`$, $`\gamma _{}C_{}`$ such that $`\gamma _\pm \gamma _{}`$ under the PS action. The presence of two energy maxima (located in $`C_+`$ and $`C_{}`$ symmetrically) causes such a structure. The remarkable feature is that, when $`E`$ crosses $`E_{}`$ from below, the system undergoes a symmetry breaking (SB) phenomenon (governed by a bifurcation mechanism ) since the system ‘must’ choose to evolve either along $`\gamma _+`$ or along $`\gamma _{}`$. Dynamically, this entails the emergence of a permanent gap between the condensates’ populations. Such an effect is also called a self-trapping effect in that, within a finite range of $`E`$, the system never leaves the region $`C_+`$ ($`C_{}`$) where was initially placed to go in $`C_{}`$ ($`C_+`$). The investigation at the quantum level of the scenario just described has been performed in Ref. through the model Hamiltonian $`H=U(n_1^2+n_2^2)vNT(a_1a_2^++a_2a_1^+),`$ (2) which represents the quantum counterpart of the Hamiltonian $`(\psi _1,\psi _2)`$ for the Laplacian-free Eqs. (1). A simple way to obtain $`H`$ relies on the fact that for low numbers of bosons per well (a realizable experimental situation) one can replace the condensate wave functions $`\psi _j`$’s with the raising (lowering) operator $`a_i`$, ($`a_i^+`$) obeying the canonical commutators $`[a_i,a_j^+]=\delta _{ij}`$, $`i=1,2`$. Since $`[N,H]=0`$, the total number of bosons $`N:=n_1+n_2`$, $`n_i:=a_i^+a_i`$, is a constant of motion (to simplify, we shall denote its eigenvalue by $`N`$ as well). The rigorous derivation of $`H`$ is effected in Ref. by using a mode dependent form of the field operator $`\widehat{\mathrm{\Psi }}_j`$ within the many-body quantum theory of BEC’s. It is important to recall as well that model (2) has been studied also in Ref. from the viewpoint of dynamical system theory. In this paper we investigate the spectral properties of model (2) by combining the use of the PS and of a further symmetry involving the change $`TT`$. The latter will be called odd symmetry (OS) for recalling its basic role in determining the structure of the energy spectrum when the total boson number $`N`$ is odd. Such symmetries are used extensively to show that: (i) after recovering the known nondegeneracy of the Hamiltonian spectrum (see Ref. ), each energy eigenstate is either symmetric or antisymmetric under the PS action, (ii) the doublets (pairs of close energy levels occurring in the energy spectrum when the model parameters range in a suitable interval) always pair a symmetric eigenstates with an antisymmetric one; we shall show how this feature plays a basic role in the classic limit, (iii) the separation mechanism causing the splitting of the energy levels (the splitting effect has been observed numerically in Refs. ) can be explained in a purely analytic way, (iv) the fact that total boson number is even/odd dramatically influences the eigenvalues’ parity under the OS. As to point (i), we wish to emphasize that the main consequence of the nondegeneracy is to prevent the SB phenomenon occurring in $`𝒫`$ as well as the ensuing self-trapping of the system on a specific orbit of the two ones associated with a given energy $`E>E_{}`$. This reflects the intrinsic tunneling effect due to the quantum nature of the system. For attaining results (i)-(iv) we first reformulate the initial quantum problem of two coupled (identical) wells through the Schwinger realization of the spin algebra in terms of two-boson operators. This allows one to reconstruct Hamiltonian (2) within its dynamical algebra (this is introduced in Appendix A). The form of $`H`$ thus obtained can be interpreted in terms of a bosonic model defined on a nonhomogeneous linear lattice with one effective boson (single-boson picture). Such a picture has been derived in Ref. and is reviewed in Sec. II together with the related formal background. Thanks to its group-theoretic character such a formulation is applicable to many-well systems with any boson number $`N`$. In Sec. III the diagonalization of $`H`$ is faced in a systematic way by making explicit the PS and OS action on the components of the energy eigenstates. The resulting characterization of the eigenstates leads to identify the recursive expression of the eigenvalue equation of $`H`$ both for even $`N`$ and for odd $`N`$. This, in turn, enables us to recognize explicitly a hidden parameter able to control in an analytic way the level distance of the doublets. Such a parameter seems to suggest an alternative procedure to evaluate perturbatively the energy levels. A similar mathematical problem was analysed in Ref. as to the problem of the dynamical tunneling through a separatrix, where the splittings of the doublets were traced by using the standard perturbation method. ## II spin picture of the two-well model A convenient way to study the spectrum of $`H`$ consists in reformulating $`H`$ by means of the Schwinger picture of spin operators. The latter is a two-boson realization of the spin operators $$J_1=\frac{a_1a_2^++a_2a_1^+}{2},J_2=\frac{a_1a_2^+a_2a_1^+}{2i},J_3=\frac{n_2n_1}{2},$$ satisfying the commutators $`[J_a,J_b]=iϵ_{abc}J_c`$ of the algebra su(2), where $`a,b,c=1,2,3`$, and $`ϵ_{abc}`$ is the totally antisymmetric tensor . The generators $`J_a`$ of su(2) commute with the Casimir operator $`CJ_1^2+J_2^2+J_3^2J_4(J_4+1)`$ which, in the Schwinger picture, leads to the identification $`J_4(n_2+n_1)/2`$. Consistently, one can check that $`[J_4,J_a]0`$. Therefore, su(2) can be used to rewrite Hamiltonian (2) which takes the nonlinear form $`H=2[UJ_4^2vJ_4+UJ_3^2TJ_1].`$ (3) The spin picture embodies explicitly in $`H`$ the dimension $`(2J+1)`$ of the Hilbert space $`(N)`$ of physical states of $`H`$, where the eigenvalue $`J`$($`=N/2`$) of $`J_4`$ is the index of the spin representation . The standard basis $`_N=\{\|J,m,\|m\|J=N/2\}`$ ($`J_3\|J,m=m\|J,m`$) of $`(N)`$ is related to the number operator states through $`\|J,m=\|n_1,n_2`$, where $`J=(n_2+n_1)/2`$, $`m=(n_2n_1)/2`$. Thus $`H`$ and $`J_a`$ can be seen as $`(2J+1)\times (2J+1)`$ matrices. The fact that $`2J_4N`$ is a constant of motion due to $`[H,N]=0`$ is ensured by $`[J_4,J_a]0`$. An important consequence made explicit by the spin picture is that the nonlinear term $`J_3^2`$ arising in Eq. (3) prevents the standard procedure of diagonalization of $`H`$ via a unitary trasformation of the group SU(2). Such a procedure works only for matrices that can be written as linear combinations of the algebra generators in that, by construction, they can be reduced to one of the generators by some appropriate group transformations . Since the diagonalization is greatly simplified by such a reduction, the latter practically identifies with the diagonalization itself. When a matrix $`𝒪`$ contains nonlinear terms of its generating algebra (as in the case of Hamiltonian (3) with respect to su(2)), resorting to a larger algebraic structure $`𝒜`$ enables one to express $`𝒪`$ as a linear combination of generators of $`𝒜`$. Such an enlarged algebra is called a dynamical algebra for $`𝒪`$. About Hamiltonian $`H`$ the problem is solved by representing the algebra su$`__N`$(2) \[subscript $`N(=2J)`$ recalls the algebra link with the total boson number\] as a sub-algebra of $`𝒜`$ = su($`M`$) for a suitable values of $`M`$. In Appendix A we show that $`𝒜`$ $``$ su$`__1`$($`N`$+1). The realization in $`𝒜`$ of the su$`__N`$(2) generators occurring in $`H`$ reads $$J_3=\mathrm{\Sigma }_q(J+1q)n_q,J_+=\mathrm{\Sigma }_q[q(2J+1q)]^{\frac{1}{2}}b_q^+b_{q+1},$$ $`J_{}(J_+)^+`$ ($`J_\pm =J_1\pm iJ_2`$), where $`b_q^+`$ ($`b_q`$) are raising (lowering) bosonic operators, $`n_q=b_q^+b_q`$, and $`q[1,N+1]`$. Recalling that the constraint $`N_b:=\mathrm{\Sigma }_qn_q1`$ must be accounted for (see Appendix A) we shall call the realization just obtained single-boson picture. This furnishes the simplest way to represent in a linear form the nonlinear term of $`H`$. In fact, Hamiltonian (3) becomes $`H=C+2U\mathrm{\Sigma }_q[m^2(q)n_q\tau R(q,J)(b_q^+b_{q+1}+h.c.)],`$ (4) with $`\tau :=T/U`$, $`R(q,J):=[(J+1/2)^2(m(q)1/2)^2]^{1/2}`$, $`m(q):=J+1q`$, and $`C=2[UJ^2vJ]`$, where the quadratic term has been re-expressed as a linear combination of number operators $`n_q`$. A valuable alternative form of the original model (2) is offered by Eq. (4) which recasts the two-well dynamics in terms of the dynamics of a single boson on a linear nonhomogeneous lattice. In the single-boson picture physical states are expressed as $$\|\mathrm{\Psi }=\mathrm{\Sigma }_q\psi _qb_q^+\|0,\mathrm{},0$$ (5) with the normalisation condition $`\mathrm{\Sigma }_q\|\psi _q\|^2=1`$ (see Appendix A). The system dynamics thus takes place on a hypersphere inside $`𝐂^{N+1}`$, where it can be interpreted in a classical form. This matches the fact that states (5) can be shown to be su$`__1`$($`N`$+1) coherent states . Based on Eq. (5), the Schrödinger problem $`i_t\|\mathrm{\Psi }=H\|\mathrm{\Psi }`$ can be rewritten as a set of equations of motion $`i\dot{\mathrm{\Psi }}_m=2Um^2\mathrm{\Psi }_mT\left[R_{m+1}^{_{(J)}}\mathrm{\Psi }_{m+1}+R_m^{_{(J)}}\mathrm{\Psi }_{m1}\right],`$ (6) where we have introduced $`\mathrm{\Psi }_m:=\psi _q`$, and $`R_m^{_{(J)}}:=R(q,J)`$ with $`m:=J+1q`$, to join the present formalism with the spin picture basis $`_N`$ where states are labeled by $`m[J,J]`$. It is worth noting how Eqs. (6) can be derived as well from the effective Hamiltonian $`H=C+\mathrm{\Sigma }_m[2Um^2\|\mathrm{\Psi }_m\|^2TR_m^{_{(J)}}(\mathrm{\Psi }_m^{}\mathrm{\Psi }_{m1}+c.c.)],`$ (7) representing the energy expectation value $`\mathrm{\Psi }\|H\|\mathrm{\Psi }H`$, provided the Poisson structure $`\{\mathrm{\Psi }_m,\mathrm{\Psi }_{\mathrm{}}^{}\}=\delta _m\mathrm{}/i\mathrm{}`$ is assumed (see commment ). The time evolution of the dynamical variables (the components of $`\|\mathrm{\Psi }`$) is determined once the initial condition $`\|\mathrm{\Psi }(0)`$ at $`t=0`$ has been assigned. Upon denoting by $`X_m`$ the components of the $`H`$ eigenstates $`\|X`$, one can retrieve the secular equation $`EX_m=2Um^2X_mT\left[R_{m+1}^{_{(J)}}X_{m+1}+R_m^{_{(J)}}X_{m1}\right],`$ (8) from Eq. (6). $`X_m`$’s can be shown to be real. The procedure relying on the dynamical algebra construction has led to interpret model (2) as a lattice model with one boson via Eq. (4). This procedure has been implemented as well for illustrative purposes since it shows clearly how the nonlinearity occurring in the matrix form of $`H`$ is transferred to the coefficients of the linear combination in $`𝒜`$. We point out that such a simplification also works for a linear chain of $`S`$ interacting wells with $`N`$ bosons whose Hamiltonian can be written via the generators of su$`__N`$($`S`$). The latter, in fact, can be always immersed within an algebra su($`M`$) with $`M`$ sufficiently larger than $`S`$; moreover the two-boson realization of su($`M`$) can be obtained for any $`M`$. The component form of secular equation (8) represents, at the operational level, a basic intermediate result. The latter, as shown in Sec. III, is used to characterize explicitly (in the same spirit of Bloch’s theorem for electronic wave functions) the inner spectrum structure as well as the structure of the energy eigenstates. ## III spectrum structure In order to investigate the spectrum structure, we consider first the effects of the PS and the OS on the energy eigenstates, and make explicit how such two symmetries strongly characterize the eigenstate components. Then, we employ the results of such an analysis to recast the eigenvalue equation related to equation (8) in a recursive form, and identify the parameters able to control the splitting of the energy levels. The case with half-integer $`J`$ and integer $`J`$ are treated separately. The PS is realized via the action of the unitary transformation $`U_1:=\mathrm{exp}[i\pi J_1]`$ which takes $`J_3=(n_2n_1)/2`$ into $`U_1J_3U_1^+=J_3`$. This matches the effect of the PS classical action which involves the exchange of populations $`n_1`$ and $`n_2`$. Let us introduce the hermitian operator $`P:=`$ $`\mathrm{exp}[i\pi J]U_1`$ whose action on the states $`\|m`$ (we drop the representation index $`J`$ in $`\|J,m`$ since it is fixed) is deducible from the equation $`U_1\|m=\mathrm{exp}[i\pi J]\|m`$. Owing to $`[H,U_1]=0`$, $`P`$ can be diagonalized together with $`H`$. The action of $`P`$ on the standard basis, $`P\|m=\|m`$, implies that, for a generic state $`\|\mathrm{\Psi }`$, $$P:\|\mathrm{\Psi }=\mathrm{\Sigma }_m\mathrm{\Psi }_m\|mP\|\mathrm{\Psi }=\mathrm{\Sigma }_m\mathrm{\Psi }_m\|m.$$ In particular, the $`P`$ action on an eigenstate $`\|X`$ entails the component transformation $`X_m\sigma X_m`$, where $`\sigma `$ is not fixed, in Eqs. (8). Actually, these remain unchanged since $`R_{m+1}^{_{(J)}}R_m^{_{(J)}}`$ for each $`m`$. The fact that $`P^2𝐈`$ fixes $`\sigma `$ showing how the allowed eigenvalues for $`P`$ are $`\sigma =\pm 1`$. Since $`P:X_m\pm X_m`$, each eigenstate has thus a definite symmetry character under $`mm`$. This fact suggests to reorganize the vectors basis in terms of vectors $`\|m,\pm =(\|m\pm \|m)/\sqrt{2}`$ that are either symmetric or antisymmetric. The new basis allows one to define in $`(N)`$ a symmetric (antisymmetric) subspace $`^+(N)`$ ($`^{}(N)`$) whose vectors have components we will denote by $`\mathrm{\Phi }^+`$ ($`\mathrm{\Phi }^{}`$) such that $`P\mathrm{\Phi }^\pm =\pm \mathrm{\Phi }^\pm `$. When the description in the new basis is adopted then Eqs. (8) for the eigenvector components can be written in the matrix form $$E\left[\begin{array}{c}X^+\\ X^{}\end{array}\right]=\left[\begin{array}{ccc}S_J(T)& & 0\\ & & \\ 0& & A_J(T)\end{array}\right]\left[\begin{array}{c}X^+\\ X^{}\end{array}\right],$$ (9) where $`S_J(T)`$ ($`A_J(T)`$) is the sub-matrix associated with symmetric (antisymmetric) sector, and $`0`$ represents the zero-matrix. The matrix equation (9) is separable in two independent equations for $`X^+`$ and $`X^{}`$: Their explicit form which depends on the representation index $`J`$ is displayed in the sequel. The odd symmetry (denoted by OS) is obtained by combining the action $`U_3J_1U_3^+=J_1`$ of $`U_3:=\mathrm{exp}[i\pi J_3]`$ on $`TJ_1`$ in $`H`$ with the change $`TT`$ which restores the initial form of $`H`$. Since $`J`$ can be either integer or half-integer, for considering the two cases separately it is convinient to introduce the matrix $`L_J(T):=L_m\mathrm{}`$ $`L_m\mathrm{}=2Um^2\delta _{m,\mathrm{}}T[R_m^{_{(J)}}\delta _{m,\mathrm{}+1}+R_{\mathrm{}}^{_{(J)}}\delta _{m+1,\mathrm{}}],`$ (10) where $`m,\mathrm{}=1/2,3/2,\mathrm{},J`$ if $`J`$ is half-integer, and $`m,\mathrm{}=1,2,\mathrm{},J`$ if $`J`$ is integer. Let us start with the half-integer case. In Eq. (9), the sub-matrices $`S__J(T)`$ and $`A__J(T)`$ coincide with the matrix $`L__J(T)`$ up to the quantity $`\eta TR_{1/2}^{_{(J)}}`$ which must be added to the matrix element $`L_{^{\frac{1}{2}\frac{1}{2}}}`$ with $`\eta =1(+1)`$ in the antisymmetric (symmetric) case. Representing $`U_3`$ in the basis ($`\{\|m,\pm \}`$) entails $$U_3=\left[\begin{array}{ccc}0& & D\\ & & \\ D& & 0\end{array}\right]$$ (11) in which $`D=Diag(i,i,\mathrm{})`$. The action of $`U_3`$ on any vector takes its symmetric components into the antisymmetric ones and vice versa, namely $`PU_3\mathrm{\Phi }^\pm =U_3\mathrm{\Phi }^\pm `$ if $`P\mathrm{\Phi }^\pm =\pm \mathrm{\Phi }^\pm `$. The structure of energy spectrum is reconstructed through the following three observations: i) The secular equation for $`X^+`$ and $`X^{}`$ derived from Eqs. (9) can be written in the reduced form $`EC_m=2Um^2C_mT\left[R_{m+1}^{_{(J)}}C_{m+1}+R_m^{_{(J)}}C_{m1}\right],`$ (12) with $`C=X^+,X^{}`$, for $`1/2<mJ`$, whereas $`0=[2U(1/2)^2\eta TR_{1/2}^{_{(J)}}E]C_{1/2}TR_{3/2}^{_{(J)}}C_{3/2},`$ (13) holds for $`m=1/2`$, where $`\eta 1(+1)`$ in the antisymmetric (symmetric) case. For a given $`T0`$, Eqs. (12) and (13) show that the eigenvalues $`E`$ can be seen as a set of functions $`E_a(T,\eta )`$ of $`\eta `$ with $`1/2aN/2`$, defined implicitly. The energy eigenvalue $`E_a(T,+1)`$ of each symmetric eigenstate can be derived from that $`E_a(T,1)`$ of an antisymmetric eigenstate by moving $`\eta `$ from $`1`$ to $`+1`$, and vice versa. ii) We consider here the problem of ordering the set of eigenvalues $`E_a(T,\pm 1)`$. For $`T0`$ (this eliminates the $`\eta `$ dependence) Eqs. (8) are solvable; the resulting eigenvalues are doubly degenerate since the equations for both the symmetric and the antisymmetric components are identical. Explicitly, $`T0`$ entails $`E_a(T,+1),E_a(T,1)2Um^2`$ for some $`m`$ which shows how the label $`a`$ can be identified with $`m[1/2,J]`$. The order induced by (positive) $`m`$ on the set $`\{2Um^2:\|m\|J\}`$ is inherited both by the symmetric eigenvalues $`\{E_a(T,+1)\}`$ and by the antisymmetric eigenvalues $`\{E_a(T,1)\}`$ as proven by the limit $`T0`$. This also implies that, for sufficiently small $`T`$, $`E_a(T,\pm 1)E_b(T,\pm 1)`$ if $`ab`$. iii) Implementing the action of $`U_3`$ whose matrix form is given by Eq. (11) on Eq. (9) leads to the equation $$E\left[\begin{array}{c}\stackrel{~}{X}^{}\\ \stackrel{~}{X}^+\end{array}\right]=[\begin{array}{ccc}A_J(T)& & 0\\ & & \\ 0& & S_J(T)\end{array}]\left[\begin{array}{c}\stackrel{~}{X}^{}\\ \stackrel{~}{X}^+\end{array}\right],$$ (14) where $`\stackrel{~}{X}^{}=DX^\pm `$, and $`P\stackrel{~}{X}^\pm =\pm \stackrel{~}{X}^\pm `$. The substitution $`TT`$ entails that \[we use the simplified notation $`E_a^\pm (T):=E_a(T,\pm 1)`$\] the set of the symmetric (antisymmetric) eigenvalues $`\{E_a^+(T)\}`$ ($`\{E_a^{}(T)\}`$) coincides with the set of the antisymmetric (symmetric) ones $`\{E_b^{}(T)\}`$ ($`\{E_b^+(T)\}`$), where $`1/2a,bJ`$. Notice that, in general, $`E_a^\pm (T)E_b^{}(T)`$, where not necessarily $`b`$ coincides with $`a`$. Nevertheless, for $`T0`$ $`E_a^\pm (0)E_b^{}(0)2Ua^2=2Ub^2`$ implies that $`b=a`$, as pointed out at point ($`\mathrm{𝐢𝐢}`$). As a consequence of points ($`𝐢`$)-($`\mathrm{𝐢𝐢𝐢}`$), we find that the symmetric eigenvalues are associated with the antisymmetric ones through the formula $$E_a^\pm (T)=E_a^{}(T)(1/2aJ).$$ (15) Also, since the eigenvalues equation can be cast in an iterative form via the recurrence equation $`d_m(E)=(2Um^2E)d_{m+1}(E)T^2[R_{m+1}^{_{(J)}}]^2d_{m+2}(E),`$ (16) which starts from $`0=\left[U/2E+\eta TR_{1/2}^{_{(J)}}\right]d_{\frac{3}{2}}(E)T^2[R_{1/2}^{_{(J)}}]^2d_{\frac{5}{2}}(E),`$ (17) and terminates with $`d_J(E)=2UJ^2E`$, consistently with (iii) one finds that the eigenvalues cannot be even functions of $`T`$. For integer $`J`$, the dimension of matrix $`S__J(T)`$ changes from that of matrix $`A__J(T)`$. In the antisymmetric case one finds $`A_J(T)=L_J(T)`$, while in the symmetric case ($`S_J(T):=S_{mn}`$), where the indices runs over $`0,1,\mathrm{},J`$, one finds $`S_{01}=S_{10}=TR_{\mathrm{\hspace{0.17em}1}}^{_{(J)}}`$, $`S_{mn}=L_{mn}`$ for $`m,n1`$, and $`S_{mn}=0`$ otherwise. Eqs. (12) still hold for integer $`J`$ provided $`2mJ`$. The two special cases are those corresponding to $`m=0,1`$ $$0=EC_0T\sigma R_{\mathrm{\hspace{0.17em}1}}^{_{(J)}}C_1,$$ (18) $$0=(2UE)C_1T\left[R_2^{_{(J)}}C_2+\sigma R_0^{_{(J)}}C_0\right],$$ (19) with $`C=X^+,X^{}`$. The parameter $`\sigma `$ must be set equal to one in the symmetric case ($`C=X^+`$), while in the antisymmetric case ($`C=X^{}`$) the expected elimination of the component $`X_0^{}`$ follows from setting $`\sigma =0`$. Hence the dimensions of Hilbert sub-spaces are such that dim $`^{}(N)`$= dim $`^+(N)1`$, while the secular equation for $`X^{}`$ will have a degree diminished of one. Explicitly, one has $`0=[{\displaystyle \frac{E(2UE)}{T^2}}+\sigma ^2[R_{\mathrm{\hspace{0.17em}1}}^{_{(J)}}]^2]d_2(E)E[R_2^{_{(J)}}]^2d_3(E).`$ (20) In the symmetric case a $`(J+1)`$-th degree equation for $`E`$ issues from (20) through formula (16). Comparing the eigenvalue equations for the symmetric ($`\sigma =1`$) and antisymmetric ($`\sigma =0`$) states shows that each, but one, symmetric eigenvalue merges to an antisymmetric one when $`\sigma `$ goes from $`0`$ to $`+1`$. Due to the diversity of the secular equation with $`\sigma =1`$ from that with $`\sigma =0`$, even in the case with integer $`J`$, the energy spectrum is constituted by $`2J`$ nondegenerate eigenvalues $`\{E_a^\pm (T):1aJ\}`$ that for $`T0`$ form $`J`$ pairs $`E_a^\pm (T)2Ua^2`$ and a single one $`E_0(T)`$ which goes to zero in the same limit. Also, due to the quadratic dependence of Eqs. (16), (20) on $`T`$, Eq. (15) must be replaced with $$E_a^\pm (T)=E_a^\pm (T)(1aJ),$$ (21) which, contrary to what happens with half-integer $`J`$, maps each eigenvalue in itself under $`TT`$. In addition, of course, one must consider $`E_0(T)=E_0(T)`$ as well. Figs. 1 illustrate the spectrum structure dependence on $`T/U`$ for $`N=6,7`$ (see also Fig. 2). ## IV discussion The interesting feature disclosed by the above analysis is the possibility to recognize both in the half-integer and in the integer case two inner parameters ($`\eta `$ and $`\sigma `$) that control, in a way independent of $`T`$, the level splitting generating the doublets. The limit $`T0`$ causes a coalescence of doublet levels such that $`E_a^\pm (T)2Ua^2`$ which suggests $`T`$ as a possible perturbative parameter for evaluating the level splitting. On the other hand, Fig. 2 clearly shows that each eigenvalue $`E_a^+(T)`$ remains close to its partner $`E_a^{}(T)`$ on a finite range $`I_a(T)`$ of $`T`$ indexed by the eigenvalue label $`a`$. In the half-integer case, this implies that, inside $`I_a(T)`$, indeed $`\eta `$ represents a good perturbative parameter (recall that $`E_a^\pm (T)=E_a(T,\eta )`$ with $`\eta =\pm 1`$) which allows one to evaluate $`E_a(T,\pm 1)`$ by perturbing $`E_a(T,\eta )`$, e. g., around $`\eta =0`$. For integer $`J`$, where the level separation is controlled by $`\sigma [0,1]`$, one can observe an effect similar to that showed in Fig. 2: the symmetric eigenvalue $`E_a(T,1)=E_a^+(T)`$ remains close to its antisymmetric partner $`E_a(T,0)=E_a^{}(T)`$ on a finite range. Because the function $`E_a(T,\sigma )`$ joins analitically $`E_a^{}(T)`$ to $`E_a^+(T)`$ then $`\sigma `$ can be assumed as the perturbative parameter for the present case. The actual size of the range $`I_a(T)`$ can be evinced roughly from Fig. 2, where the level separation strongly diverges only when $`E_a^\pm (T)`$ cross $`EE_{}:=NT`$ (recall that $`E_{}`$ is the energy critical value defined in the introduction; its derivation can be found in Ref. ). This rule seems to be motivated from the insensitivity of Eq. (17) (Eq. (20)) from the parameter $`\eta `$ ($`\sigma `$) in the terms $$\eta TR_{1/2}^{_{(J)}}E(\sigma ^2[R_{\mathrm{\hspace{0.17em}1}}^{_{(J)}}]^2E^2/T^2),$$ for suitable values of $`E`$, $`T`$ and of the other coefficients. Concerning the classical limit effected through $`N\mathrm{}`$, numerical simulations with large boson numbers $`N`$ show that the series of doublets becomes degenerate (coalescence of the doublet levels) thus restoring the conditions that allow for the SB effect. How recovering the latter is briefly illustrated via the following comparison between the classical and the quantum behavior of the two-well system. Classically, at a given energy $`E>E_{}`$, the system described in $`𝒫`$ evolves either on $`\gamma _{}`$ or on $`\gamma _+`$ (see Sec. I). Trajectories $`\gamma _{}`$ and $`\gamma _+`$ are such that the populations’ gap $`n_2n_1`$ weakly oscillates around opposite values $`\mu `$ and $`+\mu `$, respectively. Quantally, for $`E>E_{}`$, the combination $`\|C_a^\pm =\|X_a^+\pm \|X_a^{}`$ of the symmetric/antisymmetric eigenstates $`\|X_a^\pm `$ of the $`a`$th doublet can be shown to provide opposite nonvanishing expectation values $`J_3=\pm \chi `$ of $`J_3=(n_2n_1)/2`$. This fact is caused by the eigenstate structure and was discussed in Ref. . Then, states $`\|C_a^\pm `$ can be associated naturally to a pair of isoenergetic orbits $`\gamma _{}`$, $`\gamma _+`$ that have $`\mu 2\chi `$. Increasing $`N`$, the time-dependent mixed state $$\|\mathrm{\Psi }=e^{itE_a^+(T)/\mathrm{}}\|X_a^++e^{itE_a^{}(T)/\mathrm{}}\|X_a^{}$$ (22) (satisfying the Schrödinger problem of $`H`$) exhibits a sort of temporary self-trapping effect (i. e. the localization either on $`\gamma _{}`$ or on $`\gamma _+`$) which is repeated periodically and has a duration increasing with $`N`$. In fact, because of the oscillations of the factor $`\mathrm{exp}\{[it(E_a^+(T)E_a^{}(T)]\}`$ between $`+1`$ and $`1`$ entailing $`\|\mathrm{\Psi }\|C_a^+`$ and $`\|\mathrm{\Psi }\|C_a^{}`$, respectively, the system stays in a quantum state involving the localization on $`\gamma _\pm `$ for intervals of the order of the period $`\mathrm{\Delta }t=2\pi \mathrm{}/[E_a^+(T)E_a^{}(T)]`$ that increase when the level separation is reduced. The system remains definitively in the classical-like states (full emergence of the SB effect inducing the self-trapping) when the tunneling time from $`\|C_a^+`$ to $`\|C_a^{}`$ diverges, namely for $`[E_a^+(T)E_a^{}(T)]0`$ (coalescence of doublet levels induced by $`N\mathrm{}`$). ## V conclusions In Sec. I we have reviewed the dynamics of coupled Bose condensates described by Eqs. (1) (in the approximation with zero Laplacian terms) emphasizing the SB phenomenon that occurs in the phase space when increasing the energy over the critical value $`E_{}`$. Such a phenomenon (and the ensuing self-trapping effect on isoenergetic orbits $`\gamma _\pm 𝒫`$ where $`n_2n_1`$ oscillates around nonzero values $`\pm \mu `$) has prompted the study of model (2) which represents the quantum counterpart of model (1) within the two-mode approximation of the condensate field operator. One of the purposes of the present work was to investigate the dynamical behavior corresponding, at the quantum level \[through model (2)\], both to the SB effect and to the related self-trapping. An aspect we have particularly deepened is the quantum counterpart (level splitting) of the bifurcation mechanism generating pair of isoenergetic trajectories $`\gamma _\pm `$ when $`E`$ crosses $`E_{}`$. The formal set-up for studying the energy spectrum has been developed in Sec. II by recasting model (2) for the boson modes $`a_1`$, $`a_2`$ into the matrix form (3) within the spin formulation à la Schwinger. The spin form of Hamiltonian (3) makes easily viable the derivation of the secular equation (8). The latter has been achieved by using the dynamical algebra method (this enacts systematically the reduction of Hamiltonian nonlinearities) whose application is described in Appendix A. In addition to supplying equation (8), the use of the dynamical algebra method has shown the implicit link of the secular equation with the effective-bosons model (4). Such a model reformulates the two-well dynamics of Hamiltonian $`H`$ with $`N`$ bosons in a noticeably simplified form that consists of a single boson hopping on a nonhomogenous lattice (single boson picture). The interest for the single boson picture and the underlying formal construction is motivated by the possibility of extending it to more structured models such as a chain model of $`S`$ condensates with $`N`$ boson. The case $`S=3`$, which raises interest owing to its nonintegrable dynamical character, is presently under investigation . Sec. III has been devoted to make explicit the structure of the energy spectrum based on the symmetries of $`H`$. Upon introducing the permutational symmetry (PS) and recognizing the further odd symmetry (OS), we have employed them to characterize the Hilbert space of $`H`$ as well as the energy eigenstates, both for half-integer $`J`$ (odd number of boson $`N`$) and for integer $`J`$ (even number of boson $`N`$). The $`N`$-dependent form of the secular equation obtained in the two cases has led to the central result of Sec. III, namely to recognize the possibility of introducing in a natural way inner parameters that control the level splitting in the energy spectrum doublets. Such parameters –$`\eta `$ \[$`\sigma `$\] is defined in Eq. (17) \[Eq. (20)\] for half-integer \[integer\] $`J`$– have shown that the splitting originating the doublets can be generated explicitly in an analytic way. This fact combined with the doublets’ structure exhibited in Fig. (2) suggests that parameters $`\eta `$, $`\sigma `$ can be used as perturbative parameters in approximating the doublet levels inside the regions of the $`T/U`$ axis where the levels keep close. This approach may be preferable than the standard perturbative approach depending on the natural parameter $`T/U`$: this, in fact, can be shown to require higher and higher powers of $`T/U`$ when approximating levels far from the ground-state . We emphasize the fact that generating the level splitting via the inner parameters $`\eta `$, $`\sigma `$ can be interpreted as the quantum form of the bifurcation mechanism issuing pairs of orbits $`\gamma _\pm `$. In general, our construction should furnish the quantum framework for describing the bifurcation effects of any system whose Hamiltonian (in the critical regions of its phase space) has locally the same form of $`H`$. In this sense, both the quantum phase models and spin models in the mesoscopic system physics promise interesting applications. We notice as well that the matrix/algebraic analysis underlying the study of ‘quantum’ bifurcation effects gives a valuable, both formal and practical, tool for characterizing quantally the chaos onset in the model with $`S=3`$. The classical limit $`N\mathrm{}`$ has been commented in Sec. IV, where the symmetry breaking (SB) effect inherent in the classical two-well dynamics is recovered from the quantum scenario via superposition (22) of the symmetric and antisymmetric states of each doublet. In fact, the coalescence of the doublet levels caused by $`N\mathrm{}`$ (revealed by numerical simulations) leads, through a limiting process, to inhibit the oscillations of the system state $`\|\mathrm{\Psi }`$ between $`\|C_a^+`$ and $`\|C_a^{}`$ which ends up by coinciding with one of such states. This realizes the localization interpreted classically as the self-trapping effect. We conclude by illustrating a possible reformulation of Eq. (8) in a continuous form valid for large $`J`$ directed to extimate the low part of energy spectrum. Setting $`X_m`$$`Y_m/[(m+J)!(Jm)!]^{1/2}`$ in Eq. (8) provides $`(2Um^2E)Y_m=T[(Jm)Y_{m+1}+(J+m)Y_{m1}]`$, which reduces to the second order equation ($`\dot{F}:=dF/dx`$) $$\ddot{F}2xs\dot{F}+(2J+1sR+E/T)F/R=0,$$ (23) where $`s:=\pm 1`$ and $`x:=m/\sqrt{JR}`$ with $`R^2:=1+2UJ/T`$, while $`Y(m)\mathrm{𝑒𝑥𝑝}[\alpha m^2/J]F(m/\sqrt{JR})`$, with $`\alpha :=(sR1)/2`$, is the dependence of the (rescaled) components $`Y_m`$ on (the continuous variable) $`m`$. The assumption that $`F`$ identifies with the $`n`$-th Hermite polynomial gives the equation $`2n[2J+1R+E/T]/R`$ for $`s=+1`$ which, in turn, supplies a set of eigenvalues. Their dependence on $`T`$ is compared in Fig. 3 with the lowest part of the spectrum of $`N=20`$ bosons. This result as well as the results/observations discussed above have prompted further work that is presently in progress. ###### Acknowledgements. We are indebted to E. Guadagnini and M. Rasetti for stimulating discussions. The financial support of I.N.F.M. (Italy) and of M.U.R.S.T. (within the Project Sintesi) is gratefully acknowledged. ## A The type of enlarged algebra we deal with is $`𝒜`$ = su($`M`$) \[the latter can be viewed as the generalized version of the spin algebra with $`M^21`$ generators\]. Selecting an appropriate value of $`M`$ allows one to rewrite the nonlinear Hamiltonian $`H`$ in terms of a linear combination of generators of $`𝒜`$. This furnishes $`𝒜`$ with the status of dynamical algebra for $`H`$. To construct explicitly $`𝒜`$ it is useful to consider the two-boson form of the su(M) generators $`E_{ij}:=b_i^+b_j,(ij)`$, $`H_{ij}:=(b_i^+b_ib_j^+b_j)/2`$ that satisfy the commutators $$[E_{ij},E_{kl}]=\delta _{jk}E_{il}\delta _{il}E_{kj},$$ $$[E_{ij},H_{kl}]=(\delta _{jk}E_{ik}\delta _{jl}E_{il}+\delta _{il}E_{lj}\delta _{ik}E_{kj})/2,$$ with $`1i,j,l,kM`$. Within the present realization of su($`M`$), the representation theory of semi-simple Lie groups states that the eigenvalue $`Q`$ of the invariant operator $`N_b=\mathrm{\Sigma }_ib_i^+b_i`$ ($`[N_b,E_{ij}]=0`$) selects a specific representation of su($`M`$) \[$`N_b`$ can be viewed as the total particle number relatively to the creation (destruction) processes caused by $`b_i^+`$ ($`b_i`$)\]. In fact, the dimension of the Hilbert space basis $`B(M,Q)=\{\|n_1,\mathrm{},n_M,Q=\mathrm{\Sigma }_{^{i=1}}^_Mn_i\}`$ is given by $$dimB(M,Q)=(Q+M1)!/[(M1)!Q!],$$ where the states of the basis $`B(M,Q)`$ are defined as $`\|n_1,\mathrm{},n_M=_{i=1}^_M\|n_i`$ and the number operator states $`\|n_j`$ fulfil the equations $$b_i\|n_i=\sqrt{n_i}\|n_i1,b_i^+\|n_i=\sqrt{n_i+1}\|n_i+1.$$ The simplest realization of $`𝒜=`$ su$`__Q`$($`M`$) is achieved by setting $`Q=1`$ (single-boson picture); combining this fact with the requirement of preserving the dimension $`N+1`$ of $`(N)`$, entails $`dimB(M,1)=MN+1`$ thus providing $`𝒜=`$ su$`__1`$($`N`$+1) as a dynamical algebra for $`H`$. The states of the related basis $`\{\|q=b_q^+\|0,\mathrm{},0,q[1,N+1]\}`$ are in a one-to-one correspondence with the states $`\|J;m`$, $`\|m\|J=N/2`$, of the su$`__N`$(2) standard basis ($`J_3\|J;m=m\|J;m`$) via the index map $`m=J+1q`$.
warning/0006/hep-ph0006262.html	ar5iv	text	# NTUA-78/00 Strongly-interacting Fermions from a higher-dimensional Unified Gauge Theory ## 1 Introduction Efforts to describe the theory hidden behind not only the Higgs mechanism but also the fermion-generation pattern has led to various unified gauge theories. These typically contain fundamental scalar particles with masses around the weak scale and usually below 1 TeV. While a space-time global symmetry like supersymmetry can stabilize scalar masses at such low scales, it cannot explain why nature has chosen this particular value for the weak scale. Since unified field theories use a gauge symmetry to explain the hierarchy between the strong-interactions scale and the unification scale $`\mathrm{\Lambda }_{GUT}`$, it is reasonable to imagine that the hierarchy between the weak scale and $`\mathrm{\Lambda }_{GUT}`$ is also due to a local symmetry. This leads one to replace the perturbative Higgs sector with a non-perturbative effective one providing a dynamical symmetry breaking mechanism based on high-colour representations of ordinary QCD or on new local symmetries like the ones introduced in technicolour or top-colour scenaria. An effort to construct a unified theory which avoids fundamental scalars near the weak scale and in parallel addresses the strong CP problem appeared recently . It introduces new fermions with interchanged weak quantum-number assignments, referred to as katoptrons, which, in contrast to the hitherto known mirror fermions, interact strongly with each other under a new “horizontal” gauge interaction. While being unified with the other SM interactions at a unification scale $`\mathrm{\Lambda }_{GUT}`$ consistent not only with flavour-changing neutral-current and proton-lifetime bounds but also with small SM-neutrino masses , the katoptron horizontal interaction becomes naturally strong around the weak scale. Other dynamical symmetry breaking models can hardly offer such a unification perspective which tackles the hierarchy problem. This approach is also to be distinguished from models using weak horizontal interactions between SM fermions in attempts to understand family mixing and weak CP violation , since only the katoptrons transform under the new generation symmetry. On the other hand, generating mass for the SM fermions depends on the breaking of this new gauge symmetry, which is denoted here by $`SU(3)^{}`$ (and introduced under the notation $`SU(3)_{2G}`$ in ). While fermion composite operators which can break $`SU(3)^{}`$ dynamically exist in the theory , it has still to be shown that they assume the values needed to reproduce correctly the SM fermion-mass spectrum. The approach advocated here has some clear phenomenological advantages over previous dynamical symmetry breaking models. The $`T`$ (or alternatively the $`\mathrm{\Delta }\rho `$) parameter is naturally small, since the operators responsible for the splitting between the top- and bottom-quark masses are $`SU(2)_L\times U(1)_Y`$ invariant, in contrast with technicolour for instance, where the relevant extended-technicolour operators violate the electroweak symmetry. Moreover, within the present framework the $`S`$ parameter can also be kept small, mainly for three reasons. First, the strong gauge group is broken, second Majorana katoptron-neutrinos contribute negatively to $`S`$, and third, vertex corrections can have a large negative effect on this parameter. Interestingly enough, these vertex corrections are consistent with the large deviation of the experimentally measured weak right-handed bottom-quark coupling $`g_R^b`$ from its SM value . This is in contrast with technicolour theories for instance, which can hardly provide a cancellation mechanism for the $`S`$ parameter, since there are no sizable contributions to $`g_R^b`$ and the technicolour group is left intact. The purpose of the present study is to show explicitly how such a model could result from a more fundamental theory. In particular, the 10-dimensional $`E_8\times E_8^{}`$ model is one of the very few to possess the advantage of anomaly freedom and is extensively used in efforts to describe quantum gravity along with the observed low-energy interactions in the heterotic-string framework . It will be therefore used in the following as a starting point for our investigation. The results obtained indicate that it constitutes a very solid basis indeed for the understanding of our world. ## 2 Going from 10 to 4 dimensions ### 2.1 The use of coset spaces As promised above, we start with the gauge group $`G=E_8\times E_8^{}`$ defined in 10 space-time dimensions. The exceptional Lie group $`E_8`$ has the unique property of having its 248-dimensional fundamental and adjoint representations identical. Therefore, spin-1/2 and spin-1 particles are both taken to transform like $`(\mathrm{𝟐𝟒𝟖},\mathrm{𝟏})`$ and $`(\mathrm{𝟏},\mathrm{𝟐𝟒𝟖})`$ under the above group structure, and the theory is at this level supersymmetric and anomaly-free. The spin-1/2 fields are taken to be Weyl-Majorana fermions. By use of the 10-dimensional chirality operator, the $`(\mathrm{𝟐𝟒𝟖},\mathrm{𝟏})`$ and $`(\mathrm{𝟏},\mathrm{𝟐𝟒𝟖})`$ fermion representations are defined to be left- and right-handed respectively. Making connection with our 4-dimensional world leads one to consider 10-dimensional space-times of the form $`M^4\times B`$, where $`M^4`$ is the usual Minkowski space and $`B`$ is a 6-dimensional compactified manifold the structure of which has to be determined. Coset-space dimensional reduction (CSDR) of higher-dimensional gauge theories - provides a very elegant method of analyzing in detail the resulting 4-dimensional models. In the CSDR approach, the manifold $`B`$ is taken to be a coset space $`S/R`$, where $`S`$ and $`R`$ are compact Lie groups and $$\mathrm{dim}(B)=\mathrm{dim}(S)\mathrm{dim}(R).$$ (1) The group $`S`$ can be considered as an $`R`$-bundle over $`B`$. The group $`R`$ is taken to be a subgroup of both $`S`$ and $`G`$. The fact that $`R`$ is not trivial is a necessary condition for a non-trivial topology for the manifold $`B`$, something which is needed for the survival of chiral fermion fields in 4 dimensions. The fact that the 4-dimensional Lagrangian is independent of the extra coordinates is then guaranteed by gauge invariance. Embedding the symmetry $`R`$ in $`G`$ gives an interesting geometrical content to some of the gauge symmetries of the theory. Dimensional reduction from 10 to 4 dimensions is thus accompanied with rank reduction resulting to a surviving gauge symmetry $`HG`$, which is the centralizer $`C_G(R)`$ of the group $`R`$ in $`G`$. Furthermore, the extra compactified dimensions used in CSDR offer a natural framework for the unification of gauge and scalar fields. The latter have interaction potentials which can lead to spontaneous symmetry breaking, leading us from a unified gauge theory to the SM. Recently, higher-dimensional theories were considered beyond the classical level and were given a quantum meaning in the sense of the Wilson renormalization group in agreement with the treatment involving massive Kaluza-Klein excitations . The CSDR approach can therefore be exploited in the study of higher-dimensional unified quantum field theories independently of more general frameworks like string theory. Very strict rules , , determine which fields finally survive, i.e. remain massless, after this process, since gauge transformations have to be compensated by the action of the symmetry group $`S`$. These rules are a guiding light for model-building, ruling out groups that lead to unacceptable phenomenologies. It will be seen for instance that katoptrons surviving at low energies can be obtained only by coset spaces which are non-symmetric . In connection with initial $`E_8`$ groups, these lead interestingly enough to an $`E_6`$ unification group. In particular, one has to decompose the adjoint representations of the groups $`G`$ and $`S`$ under $`R\times H`$ and $`R`$ respectively according to $`\mathrm{adj}(G)`$ $`=`$ $`(\mathrm{adj}(R),1)+(1,\mathrm{adj}(H))+{\displaystyle \underset{i}{}}(r_i,h_i)`$ $`\mathrm{adj}(S)`$ $`=`$ $`\mathrm{adj}(R)+{\displaystyle \underset{i}{}}s_i.`$ (2) The only spin-1 fields surviving are the ones transforming under the adjoint representation of $`H`$. The spin-0 fields that appear after dimensional reduction, even though initially absent, are the ones transforming like $`h_i`$ under $`H`$, and only for those $`i`$’s for which $`r_i=s_i`$. As regards spin-1/2 fields, one decomposes the fermion representation $`F`$ of $`G`$ and the spinor representation of $`SO(6)`$ under $`R\times H`$ and $`R`$ respectively according to $`F`$ $`=`$ $`{\displaystyle \underset{i}{}}(r_i,h_i)`$ $`\sigma `$ $`=`$ $`{\displaystyle \underset{i}{}}\sigma _i.`$ (3) The only fermion fields surviving are the ones transforming like $`h_i`$ under $`H`$, and only for those $`i`$’s for which $`r_i=\sigma _i`$. When studying $`E_8\times E_8^{}`$ models , it is customary to identify the origin of the fields transforming under $`E_8^{}`$ with some obscure “hidden” or “shadow” world that interacts only gravitationally with ours. The philosophy here is different, because katoptron fermions originate from this new world. Since these fermions will finally assume the role of a dynamical Higgs sector, they should have quantum-number assignments similar (but not identical) to the ones of their SM partners. A way to achieve this goal is to make use of a discrete abelian subgroup of $`G`$ consisting of two elements, which we denote by $`Z_2^{E_6}`$. The action of its non-trivial group element corresponds to an outer automorphism that interchanges the $`E_6`$ subgroups of the two $`E_8`$’s. (Analogously, $`Z_2^{E_8}`$ interchanges the two $`E_8`$’s.) Making use of $`Z_2^{E_6}`$ has the effect of reducing further the rank of the surviving symmetry $`H`$ in a manner analogous to the construction in Ref., as will be seen shortly. In the following, a particular 6-dimensional non-symmetric coset space is analyzed and shown to lead to an acceptable phenomenology. ### 2.2 CSDR with $`S=Sp(4)`$, $`R=(SU(2)\times U(1))_{\mathrm{non}\mathrm{max}}`$ We consider a Lie group $`R=SU(2)\times U(1)`$ embedded non-maximally into $`S=Sp(4)`$ and into $`E_8G`$, i.e. into the exceptional group under which the Weyl-Majorana fermions of the model are left-handed. The Euler characteristic of $`Sp(4)/(SU(2)\times U(1))_{\mathrm{non}\mathrm{max}}`$ is equal to $`\chi =4`$, and a priori the number of copies of the fermion representations is, according to the index theorem, equal to $`\|\chi /2\|=2`$. Compactifying on $`B_0=S/R`$ leads to the following decompositions of the adjoint and spinor representations of $`SO(6)`$ and $`Sp(4)`$ under $`R`$ respectively: $`SO(6)(SU(2)\times U(1))_{\mathrm{non}\mathrm{max}},`$ $`\mathrm{𝟒}=(\mathrm{𝟏},0)+(\mathrm{𝟏},2)+(\mathrm{𝟐},1)`$ $`\overline{\mathrm{𝟒}}=(\mathrm{𝟏},0)+(\mathrm{𝟏},2)+(\mathrm{𝟐},1)`$ $`Sp(4)(SU(2)\times U(1))_{\mathrm{non}\mathrm{max}},`$ $`\mathrm{𝟏𝟎}=(\mathrm{𝟏}+\mathrm{𝟑},0)+(\mathrm{𝟐},\pm 1)+(\mathrm{𝟏},\pm 2)`$ (4) The adjoint representation of $`E_8`$ decomposes under $`(SU(2)\times U(1))_{\mathrm{non}\mathrm{max}}\times E_6`$ as follows: $`\mathrm{𝟐𝟒𝟖}`$ $`=`$ $`(\mathrm{𝟏},0,\mathrm{𝟕𝟖})+(\mathrm{𝟏}+\mathrm{𝟑},0,\mathrm{𝟏})+(\mathrm{𝟐},\pm 3,\mathrm{𝟏})`$ (5) $`+`$ $`(\mathrm{𝟏},2,\mathrm{𝟐𝟕})+(\mathrm{𝟐},1,\mathrm{𝟐𝟕})+(\mathrm{𝟏},2,\overline{\mathrm{𝟐𝟕}})+(\mathrm{𝟐},1,\overline{\mathrm{𝟐𝟕}})`$ These decompositions are not altered if the compactification is performed on the space $`B=(S/R)\times (Z_2^{E_8}/Z_2^{E_6})`$, i.e. when $`R`$ is replaced by $`\stackrel{~}{R}R\times Z_2^{E_6}`$ and the corresponding fields are taken to be $`Z_2^{E_6}`$ singlets, with $`Z_2^{E_6}`$ defined as previously. Since $`C_{E_8}(R)=E_6(\times Z_2^R\times U(1))`$, with $`Z_2^R\times U(1)`$ the center of $`SU(2)\times U(1)`$ (the superscripts of the various $`Z_2`$ symmetries in this paper have each obviously different meaning), and $`C_{E_8^{}}(E_6^{})=SU(3)^{}`$, the centralizer $`C_G(\stackrel{~}{R})`$ is equal to $$H=E_6^D\times SU(3)^{}(\times Z_2^R\times U(1)),$$ (6) where $`E_6^D`$ is the diagonal subgroup of the $`E_6`$ subgroups of the two $`E_8`$’s. That $`H`$ given above is indeed the surviving gauge symmetry after compactification can be checked explicitly by enumerating the spin-1 degrees of freedom which are left invariant by the action of $`\stackrel{~}{R}`$. In the absence of the $`Z_2^{E_6}`$ symmetry, $`H`$ would have been given by $`E_6\times E_8^{}(\times Z_2^R\times U(1))`$. The role of $`Z_2^{E_6}`$ is to keep only the diagonal subgroup of $`E_6\times E_6^{}E_8\times E_8^{}`$ unbroken, eliminating all skew-symmetric contributions. Physically, it renders the compactification process more symmetric with respect to the two $`E_8`$’s. Furthermore, the SM-fermion quantum numbers under $`Z_2^R\times U(1)`$, the center of $`SU(2)\times U(1)`$, are equal to $`(1,2),(1,1,)`$ and $`(1,1)`$. The center survives after CSDR, and it is identified in the following with the family symmetry of the SM which differentiates between the three SM generations. This symmetry is taken to be global and the $`U(1)`$ coupling is accordingly switched-off in order to avoid problems with flavour-changing neutral currents at lower energies. Under the gauge structure $`H=E_6^D\times SU(3)^{}(\times Z_2^R\times U(1))`$ defined above, the CSDR rules give the following surviving 4-dimensional fields: $`\mathrm{spin}1/2:`$ $`(\mathrm{𝟐𝟕},\mathrm{𝟑})`$ $`\mathrm{Katoptrons}`$ $`(\mathrm{𝟕𝟖},\mathrm{𝟏})+(\mathrm{𝟏},\mathrm{𝟖})`$ $`\mathrm{Vector}\mathrm{fermions}`$ $`2\times (\mathrm{𝟐𝟕},\mathrm{𝟏})_a`$ $`\mathrm{Include}\mathrm{SM}\mathrm{fermions}`$ $`\mathrm{spin}0:`$ $`2\times (\mathrm{𝟐𝟕},\mathrm{𝟏})_a`$ $`\mathrm{Higgs}\mathrm{sector}`$ (7) where the fields transforming under $`E_6^D`$ as $`\mathrm{𝟐𝟕}`$ and $`\overline{\mathrm{𝟐𝟕}}`$ are identified by the Majorana condition, and the subscript $`a=1,2,3`$ serves as a generation index corresponding to $`Z_2^R\times U(1)`$ . In the above, we also indicate in which sector fields which are known to us from the SM and katoptron model are contained. The vector fermions transforming like $`(\mathrm{𝟏},\mathrm{𝟖})`$ and $`(\mathrm{𝟕𝟖},\mathrm{𝟏})`$ are not protected by any gauge symmetry, so according to the general argumentation on the survival hypothesis they acquire large gauge-invariant masses of the order of the compactification scale and disappear from the low-energy spectrum. The torsion of the non-symmetric space $`B`$ is taken to be such that katoptrons remain massless. The present coset space admits two different scales , something that could be useful in the subsequent breaking of $`H`$, as will be discussed later. Furthermore, it can be checked that this breaking leads to an anomaly-free 4-dimensional theory, since it satisfies the equation $$l(G)=60,$$ (8) where $`l(G)`$ is the sum of the indices of all the representations of $`R`$ appearing when the adj($`G`$) representation is decomposed under $`R\times H`$ . After having analyzed the geometrical rank reduction of $`G`$ down to $`H`$, one has to study the subsequent breaking of $`H`$ down to the SM. ## 3 Symmetry breaking to the Standard Model ### 3.1 Breaking by Wilson lines The simply-connected group $`S=Sp(4)`$ considered has a $`Z_2`$ symmetry as center (recall that $`Sp(4)/Z_2SO(5)`$), and this can be employed here to serve in a gauge-symmetry breaking mechanism by Wilson lines . The embedding of this abelian discrete symmetry in the $`SU(2)_L`$ subgroup of $`E_6^D`$, which we denote by $`Z_2Z_2^{SU(2)_L}`$, can be defined via the following homomorphism involving its non-trivial group-element $`g`$ : $$Z_2^{SU(2)_L}gU_g=\mathrm{𝟏}\mathrm{𝟏}\left(\begin{array}{ccc}1& & \\ & 1& \\ & & 1\end{array}\right)SU(3)_C\times SU(3)_R\times SU(3)_L$$ (9) where $`SU(3)^3`$ is a maximal subgroup of $`E_6^D`$, the second and third row of the (diagonal) $`SU(3)_L`$ factor correspond to $`SU(2)_L`$, and of course $`U_g^2=1`$. We then consider dimensional reduction over the coset space $$\stackrel{~}{B}=\left(S/(R\times Z_2^{SU(2)_L})\right)\times (Z_2^{E_8}/Z_2^{E_6}),$$ (10) The original gauge group $`G`$ is broken at the compactification scale, which is identified here with the gauge-coupling unification scale $`\mathrm{\Lambda }_{GUT}`$, down to $$H^{}=SU(6)\times SU(2)_L\times SU(3)^{}(\times Z_2^R\times U(1)),$$ (11) where $`SU(6)\times SU(2)_LE_6^D`$. The original $`\|\chi /2\|=2`$ copies of fermion and scalar fields are then further reduced to a single copy due to the action of $`Z_2^{SU(2)_L}`$. Obviously, one has to distinguish the topological role of this symmetry from the role of $`Z_2^R`$ which merely differentiates the fermion families via quantum numbers. The effect of $`Z_2^{SU(2)_L}`$ is apparently at the heart of the parity asymmetry of our world. One of its side-effects is to break the original supersymmetry at the compactification scale, since fermion fields lose some of their bosonic partners. It is reminded here that the present model does not need low-energy supersymmetry, since the hierarchy problem is solved by the gauge symmetry $`SU(3)^{}`$ . ### 3.2 Further breaking by a Higgs mechanism One of the scenarios presented in is subsequently realized. The Higgs fields transform under $`SU(6)\times SU(2)_L`$ like $`(\mathrm{𝟔},\mathrm{𝟐})+(\mathrm{𝟏𝟓},\mathrm{𝟏})`$. Only the $`(\mathrm{𝟏𝟓},\mathrm{𝟏})`$ Higgses which are invariant under $`Z_2^{SU(2)_L}`$ remain light, and one of their copies is taken to develop a non-zero vacuum expectation value at the compactification scale. This breaks spontaneously the gauge symmetry further down to $$H^{\prime \prime }=SU(4)_{PS}\times SU(2)_R\times SU(2)_L\times SU(3)^{}(\times Z_2^R\times U(1)),$$ (12) where $`SU(4)_{PS}`$ is the usual Pati-Salam symmetry. The $`\mathrm{𝟐𝟕}`$ representation of $`E_6`$ decomposes under $`SU(4)_{PS}\times SU(2)_R\times SU(2)_L`$ like $$\mathrm{𝟐𝟕}=(\overline{\mathrm{𝟒}},\mathrm{𝟏},\mathrm{𝟐})+(\mathrm{𝟏},\mathrm{𝟐},\mathrm{𝟐})+(\mathrm{𝟒},\mathrm{𝟐},\mathrm{𝟏})+(\overline{\mathrm{𝟒}},\mathrm{𝟏},\mathrm{𝟏})+(\mathrm{𝟏},\mathrm{𝟐},\mathrm{𝟏})+(\mathrm{𝟏},\mathrm{𝟏},\mathrm{𝟏})$$ (13) The remaining generations of Higgses transform under $`SU(4)_{PS}\times SU(2)_R`$ like $`(\mathrm{𝟒},\mathrm{𝟐})+(\overline{\mathrm{𝟒}},\mathrm{𝟏})+(\mathrm{𝟏},\mathrm{𝟐})+(\mathrm{𝟏},\mathrm{𝟏})`$. The $`(\mathrm{𝟒},\mathrm{𝟐})`$ Higgs field is subsequently taken to acquire a non-zero vacuum expectation value and break spontaneously the gauge symmetry $`SU(4)_{PS}\times SU(2)_R`$ further down to $`SU(3)_C\times U(1)_Y`$ at the Pati-Salam scale $`\mathrm{\Lambda }_{PS}`$, giving the final symmetry $$H^{\prime \prime \prime }=SU(3)_C\times SU(2)_L\times U(1)_Y\times SU(3)^{}(\times Z_2^R\times U(1)),$$ (14) which includes the familiar SM groups. The fact that the coset space considered in the last section admits two different scales could be at the origin of the relatively small hierarchy between the Pati-Salam symmetry breaking scale $`\mathrm{\Lambda }_{PS}`$ and the unification scale $`\mathrm{\Lambda }_{GUT}`$ . There is no surviving symmetry preventing scalar particles from obtaining large masses after these breakings. Spin-1/2 particles remain light only if they are chiral, the others gaining compactification-scale masses. One then recovers the gauge and matter content which is the starting point of , by taking the fermions in the representations of the original $`E_8`$ and $`E_8^{}`$ groups to be left-handed and right-handed respectively, as stated in the beginning. Therefore, one reproduces the three SM generations and a katoptron generation which has interchanged left-right $`SU(2)_L\times SU(2)_R`$ quantum numbers and transforms in addition in the fundamental representation of the gauge group $`SU(3)^{}`$. Under $`H^{\prime \prime \prime }`$, the fields transform like $`\mathrm{SM}\mathrm{fermions}`$ $`\mathrm{Katoptrons}`$ $`q_L:(\mathrm{𝟑},\mathbf{\hspace{0.33em}2},\mathrm{\hspace{0.33em}1}/3,\mathbf{\hspace{0.33em}1})_a`$ $`q_R^K:(\mathrm{𝟑},\mathbf{\hspace{0.33em}2},\mathrm{\hspace{0.33em}1}/3,\mathbf{\hspace{0.33em}3})`$ $`l_L:(\mathrm{𝟏},\mathbf{\hspace{0.33em}2},1,\mathbf{\hspace{0.33em}1})_a`$ $`l_R^K:(\mathrm{𝟏},\mathbf{\hspace{0.33em}2},1,\mathbf{\hspace{0.33em}3})`$ $`q_R^c:(\overline{\mathrm{𝟑}},\mathbf{\hspace{0.33em}1},{}_{+2/3}{}^{4/3},\mathbf{\hspace{0.33em}1})_a`$ $`q_L^{Kc}:(\overline{\mathrm{𝟑}},\mathbf{\hspace{0.33em}1},{}_{+2/3}{}^{4/3},\mathbf{\hspace{0.33em}3})`$ $`l_R^c:(\mathrm{𝟏},\mathbf{\hspace{0.33em}1},{}_{2}{}^{0},\mathbf{\hspace{0.33em}1})_a`$ $`l_L^{Kc}:(\mathrm{𝟏},\mathbf{\hspace{0.33em}1},{}_{2}{}^{0},\mathbf{\hspace{0.33em}3}),`$ where the superscript $`K`$ denotes katoptron fields, $`c`$ charge conjugation, the subscripts $`L`$ and $`R`$ left- and right-handed fields, and $`q`$ and $`l`$ quark and lepton fields respectively. The group $`SU(3)^{}`$ is asymptotically free and provides the mechanism responsible for the dynamical breaking of the electroweak symmetry at the right scale via katoptron condensates. ## 4 Discussion Starting with a higher-dimensional gauge field theory, we presented an effort to produce a picture consistent with current phenomenology and which in addition includes a dynamical Higgs sector. The need to obtain eventually the SM group structure at lower energies in 4 dimensions places severe constrains on the compactification manifolds considered. The gauge-symmetry-breaking sequence of can be reproduced by use of Wilson lines for example, if the group manifold $`S`$ has $`Z_2`$ as center, so $`S=Sp(4)`$ is left as a unique choice (for $`S`$ semisimple) leading to a 6-dimensional non-symmetric manifold $`B`$. Moreover, in order to make connection with the unification picture presented in , one has to note that the $`E_6`$ group with three generations seems to be favored over $`SO(10)`$ with 4 generations considered in that reference as a unification symmetry. In all other respects, the results and conclusions of remain unaltered, since the scenario with 4 generations was rejected there for other reasons. The likelihood of the scenario presented here should be tested not only for its theoretical consistency but also for its phenomenological relevance in forthcoming experiments . It would then constitute one interesting example trying to connect the abstract mathematical world of the Planck scale with the experimental physical reality of collider data. Acknowledgements We thank P. Forgacs, A. Kehagias, D. Lüst and D. Suematsu for interesting discussions and for reading the manuscript. G.T. thanks the National Technical University of Athens for their warm hospitality and support during the months when this work was completed. G.Z. is partially supported by the E.C. project ERBFMRXCT960090.
warning/0006/hep-th0006038.html	ar5iv	text	# 1 Introduction ## 1 Introduction In non-supersymmetric theories, the Nambu-Goldstone (NG) theorem tells us that there appear as many massless NG bosons as the number of broken generators, namely $`\mathrm{dim}(G/H)`$, when a global symmetry $`G`$ spontaneously breaks down to its subgroup $`H`$. The NG bosons parameterize a vacuum degeneracy which has one-to-one correspondence with the freedom of the embedding $`H`$ into $`G`$. The effective Lagrangian of massless bosons can be expanded by the number of space-time derivatives, and the leading term, with two derivatives, is described by nonlinear sigma models on target manifolds $`G/H`$, whose coordinates are parametrized by NG bosons. On this manifold, the unbroken symmetry $`H`$ is realized linearly, while the broken symmetry $`G`$ is realized nonlinearly by NG bosons . For non-supersymmetric cases, low-energy theorems tell us that low-energy scattering amplitudes of NG bosons are determined solely by the symmetries $`G`$ and $`H`$, and do not depend on details of the underlying theory (for a review, see Ref. ). The effective Lagrangian reproduces these low-energy theorems. In supersymmetric theories, there appear additional massless bosons called quasi-NG (QNG) bosons (and their fermionic superpartners).<sup>1</sup><sup>1</sup>1 Fermions, called the QNG fermions, would be interesting particles, when we regard quarks and leptons as QNG fermions. But we do not discuss QNG fermions in this paper. Leading terms of massless effective Lagrangian are described by $`N=1`$ supersymmetric nonlinear sigma models (for example see Ref. ). Target manifolds of $`N=1`$ nonlinear sigma models are Kähler manifolds : A manifold whose metric is given by a Kähler potential $`K(\phi ,\phi ^{})`$ $$g_{ij^{}}(\phi ,\phi ^{})=\frac{^2K(\phi ,\phi ^{})}{\phi ^i\phi ^j},$$ is called a Kähler manifold. $`\phi (x)`$ is a complex scalar component of a chiral superfield. NG and QNG bosons are coordinates of a complex coset manifold $`G^𝐂/\widehat{H}`$, where $`G^𝐂`$ is the complexification of $`G`$ and $`\widehat{H}`$ is the complex subgroup often larger than $`H^𝐂`$, the complexification of $`H`$. Kähler potentials of $`G^𝐂/\widehat{H}`$ have been constructed by Bando, Kuramoto, Maskawa and Uehara (BKMU) (for a review, see Ref. ). If $`\widehat{H}=H^𝐂`$, the number of QNG bosons is the same as that of NG bosons, and nonlinear realizations in these cases are called “maximal realizations” or “fully-doubled realizations”. On the other hand, if $`\widehat{H}(H^𝐂)`$ becomes larger, the number of QNG bosons decreases. If there is no QNG boson, realizations are called “pure realizations”, and studied extensively . Pure realizations cannot be obtained as the low-energy limit of underlying linear theories since there remains at least one QNG boson . If there is, however, gauge symmetry, it is possible to absorb pairs of a QNG and a NG bosons by the supersymmetric Higgs mechanism. Hence pure realizations are in some cases obtained as low-energy theories of gauged linear sigma models . To investigate low-energy theories with supersymmetry, it is important to understand geometric structures of supersymmetric sigma models. In the cases of pure realizations, the geometry of the target space is well understood, because Kähler potentials are uniquely determined by the metric of $`G/H`$. When there are QNG bosons, however, the coset space, $`G/H`$ where NG bosons reside, is a subspace of the target space. Since the metric in directions along QNG bosons is not determined by the geometry of its subspace $`G/H`$, the effective Lagrangian is not unique in this case and depends on an arbitrary function of many $`G`$-invariant variables . When there are many $`G`$-invariant variables, it is complicated to study geometric structures of target spaces in general. In this paper, we investigate the $`O(N)`$ model whose Kähler potential contains an arbitrary function of a single variable, but generalizations to other models are straightforward. This paper is organized as follows. We review our previous results in the rest of this section. The low-energy theorems at the symmetric points are explained. In Sect. 2, we study non-symmetric points where unbroken symmetry $`H`$ is reduced to a smaller group $`H^{}`$. In Sect. 3, to investigate the geometrical structure of the supersymmetric nonlinear sigma model with $`O(N)`$ symmetry, we show explicitly how the different compact homogeneous manifolds $`G/H`$ and $`G/H^{}`$ are embedded in the full target manifold $`G^𝐂/\widehat{H}`$ by using the method of Shore . We see how some QNG bosons change to NG bosons at the non-symmetric points. In Sect. 4, we derive the low-energy theorems of NG and QNG bosons at the general points of the target manifold when the Kähler potential is the simplest one. Sect. 5 is devoted to conclusion and discussion. In Appendix A, we explain the Kähler normal coordinate which is used to calculate the low-energy theorems. In Appendix B, some geometric quantities are calculated for the most general $`O(N)`$-invariant model. The general low-energy effective Lagrangian of massless bosons $`\varphi ^\alpha (x)`$ is a nonlinear sigma model whose target manifold has the metric $`g_{\alpha \beta }(\varphi )`$, $`={\displaystyle \frac{1}{2}}g_{\alpha \beta }(\varphi )_\mu \varphi ^\alpha ^\mu \varphi ^\beta .`$ (1.1) Low-energy scattering amplitudes are unchanged by a field redefinition, which is a general coordinate transformation in the target manifold. By expanding this in the Riemann normal coordinate $`\varphi ^i`$ up to the forth order, and regarding the fourth order terms as interaction terms $`_{\mathrm{int}}`$, low-energy two-body scattering amplitudes of the massless bosons $`\varphi ^i`$ (with momenta $`p_i`$) $`\varphi ^k(p_k),\varphi ^l(p_l)\|i_{int}\|\varphi ^i(p_i),\varphi ^j(p_j)`$ $`=i(2\pi )^4\delta ^{(4)}(p_k+p_lp_ip_j)(\varphi ^i(p_i),\varphi ^j(p_j)\varphi ^k(p_k),\varphi ^l(p_l))`$ (1.2) can be calculated by summing up all the tree graphs<sup>2</sup><sup>2</sup>2 To calculate the next-to leading order $`𝒪(p^4)`$, we need to sum up one-loop graphs of the leading term and tree graphs of the four derivative terms, with obeying Weinberg’s counting theorem., given by $`(\varphi ^i,\varphi ^j\varphi ^k,\varphi ^l)={\displaystyle \frac{1}{3f_\pi ^4}}[(su)R_{kijl}+(ut)R_{ijkl}+(ts)R_{kjli}].`$ (1.3) Here $`f_\pi `$ is the decay constant of the NG bosons (pions), $`R_{ijkl}`$ is the curvature tensor of the target manifold, and we have defined the Mandelstam’s variables by $`s\stackrel{\mathrm{def}}{=}(p_i+p_j)^2=+2p_ip_j=+2p_lp_k,`$ $`t\stackrel{\mathrm{def}}{=}(p_ip_k)^2=2p_ip_k=2p_lp_j,`$ $`u\stackrel{\mathrm{def}}{=}(p_ip_l)^2=2p_ip_l=2p_jp_l.`$ (1.4) We consider cases that a global symmetry $`G`$ spontaneously breaks down to its subgroup $`H`$. We express broken and unbroken generators by <sup>3</sup><sup>3</sup>3 The Lie algebras of the groups $`G`$ and $`H`$ are denoted by $`𝒢`$ and $``$, respectively. $`X_i𝒢,H_a,(T_A𝒢).`$ (1.5) In the cases of symmetric spaces $`G/H`$ (in which there is a symmetry $`X_iX_i,H_aH_a`$), the curvature tensor can be calculated by using structure constants of $`G`$, $`f_{AB}^{}{}_{}{}^{C}`$, to yield $`R_{ijkl}=f_\pi ^2f_{ij}^{}{}_{}{}^{a}f_{akl},`$ (1.6) and the low-energy theorems become $`(\varphi ^i,\varphi ^j\varphi ^k,\varphi ^l)={\displaystyle \frac{1}{3f_\pi ^2}}[(su)f_{ki}^{}{}_{}{}^{a}f_{ajl}+(ut)f_{kl}^{}{}_{}{}^{a}f_{aij}+(ts)f_{kj}^{}{}_{}{}^{a}f_{ali}].`$ (1.7) In $`N=1`$ supersymmetric theories, the low-energy effective Lagrangian of massless chiral superfields $`\mathrm{\Phi }^i(x,\theta ,\overline{\theta })=\phi ^i(x)+\sqrt{2}\theta \psi ^i(x)+\theta \theta F^i(x)`$ (where $`\phi ^i`$, $`\psi ^i`$ and $`F^i`$ are complex scalar fields, Weyl fermions, auxiliary scalar fields, respectively.) is a supersymmetric nonlinear sigma model , $``$ $`=`$ $`{\displaystyle d^2\theta d^2\overline{\theta }K(\mathrm{\Phi },\mathrm{\Phi }^{})}`$ (1.8) $`=`$ $`g_{ij^{}}(\phi ,\phi ^{})_\mu \phi ^i^\mu \phi ^j+ig_{ij^{}}\overline{\psi }^j\overline{\sigma }^\mu (_\mu \psi ^i+\mathrm{\Gamma }_{}^{i}{}_{lk}{}^{}_\mu \phi ^l\psi ^k)`$ $`+{\displaystyle \frac{1}{4}}R_{ij^{}kl^{}}\psi ^i\psi ^k\overline{\psi }^j\overline{\psi }^l.`$ Here the metric tensor is calculated by the Kähler potential as $`g_{ij^{}}(\phi ,\phi ^{})=_i_j^{}K(\phi ,\phi ^{}),`$ (1.9) and $`R_{ij^{}kl^{}}`$ and $`\mathrm{\Gamma }_{}^{i}{}_{lk}{}^{}`$ are the complex curvature and the connection, respectively. In Eq. (1.8), the auxiliary fields $`F^i`$ have been eliminated by using their equations of motion. The massless chiral NG superfields appear when the global symmetry $`G`$ spontaneously breaks down to its subgroup $`H`$ with preserving $`N=1`$ supersymmetry. We denote complex broken and unbroken generators as <sup>4</sup><sup>4</sup>4 Complex generators are complex linear combinations of Hermitian generators of $`𝒢`$. We use indices $`R,S,T`$ and $`L,M,N`$ for complex broken and unbroken generators, respectively. $`Z_R𝒢^𝐂\widehat{},K_M\widehat{}.`$ (1.10) Their commutation relations are $`[K_M,K_N]=if_{MN}^{}{}_{}{}^{L}K_L,[Z_R,K_M]=if_{RM}^{}{}_{}{}^{S}Z_S,[Z_R,Z_S]=if_{RS}^{}{}_{}{}^{M}K_M,`$ (1.11) where we have assumed the existence of an automorphism $`Z_RZ_R,K_MK_M.`$ (1.12) The target manifold (which is a $`G^𝐂`$-orbit of vacuum vector $`\stackrel{}{v}`$) is a complex coset manifold $`G^𝐂/\widehat{H}`$, and its representative is $`\xi (\mathrm{\Phi })=\mathrm{e}^{i\mathrm{\Phi }Z}G^𝐂/\widehat{H},\mathrm{\Phi }Z={\displaystyle \underset{i=1}{\overset{N_\mathrm{\Phi }}{}}}\mathrm{\Phi }^iZ_R\delta _i^R.`$ (1.13) Here $`\mathrm{\Phi }^i(x,\theta ,\overline{\theta })`$ are the NG chiral superfields and $`N_\mathrm{\Phi }`$ is a number of $`\mathrm{\Phi }^i`$. The left action of $`G`$ on the coset representative is $`\xi \stackrel{g}{}\xi ^{}=g\xi \widehat{h}^1(g,\xi ),gG,`$ (1.14) where $`\widehat{h}(g,\xi )\widehat{H}`$ is called an $`\widehat{H}`$-compensator. It is known that, when a vacuum vector $`\stackrel{}{v}`$ is in the real representation of $`G`$, or when $`G/H`$ is a symmetric space, only maximal realizations are possible . So we discuss maximal realizations, where there appear the same numbers of NG and QNG bosons (at a symmetric point defined below). The low-energy effective Kähler potential can be written as $`K(\mathrm{\Phi },\mathrm{\Phi }^{})=f(\stackrel{}{v}^{}\xi ^{}(\mathrm{\Phi }^{})\xi (\mathrm{\Phi })\stackrel{}{v}),`$ (1.15) where $`f`$ is an arbitrary function, which cannot be determined by symmetry.<sup>5</sup><sup>5</sup>5 If there are some $`G`$-invariants, a Kähler potential can be written as an arbitrary function of such invariants. This arbitrariness is a characteristic feature of non-pure realizations. Note that this Kähler potential is $`G`$-invariant by Eq. (1.14) but $`not`$ $`G^𝐂`$-invariant: A $`G`$-action is a general coordinate transformation preserving the metric (the Kähler potential), while a $`G^𝐂`$-action does not preserve the metric. This fact has an important consequence: The symmetry of the action is still the compact real group $`G`$, although the target space is a $`G^𝐂`$-orbit of vacuum vector $`\stackrel{}{v}`$. A holomorphic vielbein $`E_i^R`$ and a canonical $`\widehat{H}`$-connection $`W_i^M`$ can be read as coefficients of broken and unbroken elements of the Maurer-Cartan $`1`$-form $`{\displaystyle \frac{1}{i}}\xi (\phi )^1d\xi (\phi )=(E_i^R(\phi )Z_R+W_i^M(\phi )K_M)d\phi ^i.`$ (1.16) We define symmetric points by points with the largest unbroken symmetry. (As seen in the next section, there exist points with smaller unbroken symmetry.) We can take a coordinate system $`\phi `$ on $`G^𝐂/\widehat{H}`$ so that the origin of $`\phi `$ is a symmetric point. Then at a symmetric point $`\phi =0`$, the vielbein and the $`\widehat{H}`$ connection take the form of $`E_i^R\|_{\phi =0}=\delta _i^R,W_i^M\|_{\phi =0}=0,`$ (1.17) respectively, and differentiations of the vielbein with respect to coordinates at the point are $`_jE_i^R\|_{\phi =0}=0.`$ (1.18) From Eqs. (1.16) to (1.18), We can calculate the curvature tensor at the symmetric point, given by $`R_{ij^{}kl^{}}=f_1(\stackrel{}{v}^{}Z_S^{}Z_V^{}Z_UZ_R\stackrel{}{v})\delta _i^R(\delta _j^S)^{}\delta _k^U(\delta _l^V)^{}+g^2(\delta _{ik}\delta _{j^{}l^{}}+\delta _{ij^{}}\delta _{kl^{}}+\delta _{il^{}}\delta _{kj^{}}),`$ (1.19) where we have defined $`v^2\stackrel{\mathrm{def}}{=}\stackrel{}{v}^{}\stackrel{}{v}`$, $`f_1\stackrel{\mathrm{def}}{=}f^{}(v^2),f_2\stackrel{\mathrm{def}}{=}f^{\prime \prime }(v^2)`$ etc. and a constant $`g`$ by $`g^2\stackrel{\mathrm{def}}{=}f_2v^2.`$ (1.20) We express complex scalar fields by $`\phi ^i(x)=A^i(x)+iB^i(x)`$, where $`A^i(x)`$ and $`B^i(x)`$ are real scalar fields. In maximal realization cases, $`A^i(x)`$ and $`B^i(x)`$ are NG and QNG bosons, respectively. In the real basis of $`A^i`$ and $`B^i`$, the Kähler condition on the curvature tensor becomes $`\{\begin{array}{cc}R_{A^iA^jA^kA^l}=R_{B^iB^jB^kB^l}=R_{A^iA^jB^kB^l}=R_{B^iB^jA^kA^l},\hfill & \\ R_{B^iA^jB^kA^l}=R_{A^iB^jA^kB^l}=R_{B^iA^jA^kB^l}=R_{A^iB^jB^kA^l},\hfill & \\ R_{B^iA^jA^kA^l}=R_{A^iB^jA^kA^l}=R_{A^iB^jB^kB^l}=R_{B^iA^jB^kB^l},\hfill & \\ R_{A^iA^jB^kA^l}=R_{A^iA^jA^kB^l}=R_{B^iB^jA^kB^l}=R_{B^iB^jB^kA^l}.\hfill & \end{array}`$ (1.21) We can calculate real components of the curvature tensor (at symmetric points), which are directly related with low-energy scattering amplitudes, given by $`R_{A^iA^jA^kA^l}=f_\pi ^2f_{RS}^{}{}_{}{}^{M}f_{MUV}\delta _i^R\delta _j^S\delta _k^U\delta _l^V,`$ $`R_{B^iA^jB^kA^l}=f_1\stackrel{}{v}^{}(Z_S\{Z_V,Z_U\}Z_R+Z_R\{Z_V,Z_U\}Z_S)\stackrel{}{v}\delta _i^R\delta _j^S\delta _k^U\delta _l^V`$ $`4g^2(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{kj}+\delta _{ij}\delta _{kl}),`$ $`R_{B^iA^jA^kA^l}=R_{A^iA^jB^kA^l}=0,`$ (1.22) where we have defined $`f_\pi ^2\stackrel{\mathrm{def}}{=}2f_1v^2`$. We thus obtain low-energy ($`𝒪(p^2)`$) scattering amplitudes of the NG and QNG bosons, by substituting Eqs. (1.22) and (1.21) to Eq. (1.3). We conclude that, at a symmetric point, there exist low-energy theorems of amplitudes which include only the NG bosons, where higher derivatives of the arbitrary function cancel out, and they coincide with scattering amplitudes among NG bosons in non-supersymmetric theories on a symmetric space $`G/H`$ (see Eq. (1.7)). Amplitudes among only the QNG bosons coincide with those of the corresponding NG bosons by the Kähler conditions (1.21). Amplitudes for even number of the NG and QNG bosons depend on the second derivative of the arbitrary function. We would like to generalize these results to low-energy theorems at general points. At non-symmetric points, some of the QNG bosons turn to NG bosons, corresponding to the fewer unbroken symmetry. ## 2 Non-symmetric points and supersymmetric vacuum alignment In the last section we have discussed low-energy theorems of NG and QNG bosons at a symmetric point. In supersymmetric low-energy theories, there can exist points with smaller unbroken symmertry in the same vacuum manifold, as a result of the supersymmetric vacuum alignment. In this section we discuss how this phenomenon occurs. ### 2.1 Non-symmetric points In maximal realizations, the number of the QNG bosons is equal to that of QNG bosons at symmetric points. If we leave from the symmetric point by a $`G`$-action as $`\stackrel{}{v}^{}=g\stackrel{}{v}(gG)`$, they are also symmetric points and the unbroken symmetry remains unchanged : $`H^{}=gHg^1H`$. The Kähler potential, the metric and the curvature tensor do not change and the low-energy theorems do not change either. All of them are equivalent vacua. The full target manifold is, however, constructed by $`G^𝐂`$-actions on $`\stackrel{}{v}`$. If we move to another vacua by a $`G^𝐂`$-action, the unbroken symmetry $`H`$ varies depending on the choice of vacuum. A $`G^𝐂`$-action on the symmetric point $`\stackrel{}{v}`$ is $`\stackrel{}{v}^{}=g_0\stackrel{}{v},g_0G^𝐂.`$ (2.1) Complex unbroken subgroups at $`\stackrel{}{v}`$ and $`\stackrel{}{v}^{}`$, defined by $`\widehat{H}\stackrel{}{v}=\stackrel{}{v},\widehat{H}^{}\stackrel{}{v}^{}=\stackrel{}{v}^{},`$ (2.2) are related by $`\widehat{H}^{}=g_0\widehat{H}g_{0}^{}{}_{}{}^{1}\widehat{H}:K_{M}^{}{}_{}{}^{}=g_0K_Mg_{0}^{}{}_{}{}^{1}.`$ (2.3) In this sense, the complex unbroken generators are equivalent at any point on the manifold . Complex broken generators are also related by $`𝒢^𝐂\widehat{}^{}=g_0(𝒢^𝐂\widehat{})g_{0}^{}{}_{}{}^{1}:Z_{R}^{}{}_{}{}^{}=g_0Z_Rg_{0}^{}{}_{}{}^{1}.`$ (2.4) Since $`G^𝐂`$-orbits of $`\stackrel{}{v}`$ and $`\stackrel{}{v}^{}`$ are homeomorphic to each other $`G^𝐂/\widehat{H}G^𝐂/\widehat{H}^{},`$ (2.5) the transformation of $`G^𝐂`$ is just an automorphism on the target manifold. To be precise, (bosonic part of) the representatives of the complex cosets $`\xi =\mathrm{exp}(i\phi ^RZ_R)G^𝐂/\widehat{H},\xi ^{}=\mathrm{exp}(i\phi ^RZ_{R}^{}{}_{}{}^{})G^𝐂/\widehat{H}^{}`$ (2.6) are related by a right action of $`G^𝐂`$ through the relation $`\xi \stackrel{}{v}=\xi g_0^1g_0\stackrel{}{v}=\left[\xi g_0^1\widehat{h}^1(\xi ,g_0)\right]\stackrel{}{v}^{}=\xi ^{}\stackrel{}{v}^{},`$ (2.7) as $`\xi ^{}=\xi g_0^1\widehat{h}^1(\xi ,g_0),\widehat{h}^{}\widehat{H}^{}.`$ (2.8) This relation is sketched in Fig. 1. We can summarize these facts as follows: Suppose we choose a coordinate system whose origin is a symmetric point $`\stackrel{}{v}`$, and move to a non-symmetric point by an action of $`g_0`$ and define a new coordinate system whose origin is $`g_0\stackrel{}{v}`$. Then, unless $`g_0`$ belongs to the isometry $`G`$ of the metric, the origin of the new coordinate system $`\phi ^{}`$ is no longer a symmetric point. The right action (2.8) can be written explicitly as $`e^{i\phi ^{}Z^{}}=e^{i\phi Z}g_{0}^{}{}_{}{}^{1}e^{iu^{}(\xi ,g_0)K^{}},\widehat{h}^{}=e^{iu^{}(\xi ,g_0)K^{}}\widehat{H}^{},`$ (2.9) where $`u^{}`$ is a function of $`g_0`$ and $`\phi `$. It can be rewritten as $`e^{i\phi ^{}Z}=g_{0}^{}{}_{}{}^{1}e^{i\phi Z}e^{iu^{}(\xi ,g_0)K},e^{iu^{}(\xi ,g_0)K}\widehat{H},`$ (2.10) and if $`g_0`$ is restricted in $`G`$, it reduces to the ordinary left action (1.14). The right action does not change the Kähler potential from Eq. (2.7), $`\stackrel{}{v}^{}\xi ^{}\xi \stackrel{}{v}\stackrel{}{v}^{}\xi ^{}\xi ^{}\stackrel{}{v}^{}=\stackrel{}{v}^{}\xi ^{}\xi \stackrel{}{v}.`$ (2.11) We should again emphasize that this is just a coordinate transformation from the coordinate whose origin is a symmetric point $`\stackrel{}{v}`$ to one whose origin is a non-symmetric point $`\stackrel{}{v}^{}`$, but not a symmetry. ### 2.2 Supersymmetric vacuum alignment This subsection is devoted to an another interpretation, in terms of the group theory, about phenomena discussed in the last subsection, and then can be skipped. The compact subgroup $`G`$ is called a real form of $`G^𝐂`$. The operation of $`𝒢`$ on the complex algebra $`𝒢^𝐂`$ or its subalgebra picks up Hermitian generators. The real unbroken symmetry at the vacuum $`\stackrel{}{v}`$ is defined by $`H=\widehat{H}G,`$ (2.12) and at the vacuum $`\stackrel{}{v}^{}`$ by $`H^{}=\widehat{H}^{}GH,(H^{})(H),`$ (2.13) where $`()`$ denotes an equivalence class by the $`G`$ action: $`(H)=\{gHg^1\|gG\}`$. Hence, at the non-symmetric point, the unbroken symmetry group becomes smaller than at the symmetric point . This phenomenon is called the “supersymmetric vacuum alignment”. It comes from the different embedding of $`\widehat{H}`$ in $`G^𝐂`$ as in Fig. 2.<sup>6</sup><sup>6</sup>6 BKMU called the embedding corresponding to $`\stackrel{}{v}`$ and $`\stackrel{}{v}^{}`$ the “natural embedding” and the “twisted embedding”, respectively . However they discussed the natural embedding only. Kotcheff and Shore called $`\stackrel{}{v}`$ and $`\stackrel{}{v}^{}`$ the “symmetric embedding” and the “non-symmetric embedding”, respectively, since they discussed the case when $`G/H`$ is a symmetric space and $`G/H^{}`$ is a non-symmetric space . In the case of $`\widehat{H}^{}`$, $`K\stackrel{\mathrm{def}}{=}g_0Hg_0^1(H)`$ is not a subgroup of $`G`$, and the real unbroken symmetry is $`KG=H^{}(H)`$. At the non-symmetric point, the generators $`𝒦_{}^{}{}_{}{}^{𝐂}(𝒢=\varphi )`$ are not Hermitian generators. We call them “pseudo-Borel generators”.<sup>7</sup><sup>7</sup>7 A Borel algebra $``$ ($`\widehat{}`$) is defined as an algebra which satisfies, $`[,]`$. Since $`𝒦_{}^{}{}_{}{}^{𝐂}`$ does not satisfy this condition, it is not a Borel subalgebra. We show in the following sections that they correspond to NG bosons that appear at the non-symmetric point where the unbroken symmetry becomes smaller. The (real) $`G`$-orbits of $`\stackrel{}{v}`$ and $`\stackrel{}{v}^{}`$ are $`G/H`$ and $`G/H^{}(G/H)`$, respectively. They are compact submanifolds of $`G^𝐂/\widehat{H}`$ and parametrized by the NG bosons. Other directions of the total target space $`G^𝐂/\widehat{H}`$ are non-compact and correspond to the QNG bosons. ## 3 NG coset subspaces In this section, we discuss geometric structures of complex coset manifolds. To be specific we treat an simple example, the $`O(N)`$ model, but generalizations to other models are straightforward. ### 3.1 Typical example : $`O(N)`$ model In this subsection, we discuss the simplest example, the $`O(N)`$-model, where the vacuum $`\stackrel{}{v}`$ is in the vector representation of $`G=O(N)`$: $`\stackrel{}{v}V=𝐑^N`$ . We can complexify the $`O(N)`$ group by replacing the real vector with the complex vector: $`VV^𝐂=𝐂^N`$. The generators of the $`O(N)`$ group are $`(T_{ij})_{}^{k}{}_{l}{}^{}={\displaystyle \frac{1}{i}}(\delta _{i}^{}{}_{}{}^{k}\delta _{jl}\delta _{j}^{}{}_{}{}^{k}\delta _{il})=\left(\begin{array}{cccc}& & & \\ & & i& \\ & i& & \end{array}\right),`$ (3.1) where only $`(i,j)`$ and $`(j,i)`$ elements are non-zero. They satisfy commutation relations and normalization conditions $`[T_{ij},T_{kl}]=i(\delta _{jk}T_{il}\delta _{ik}T_{jl}\delta _{jl}T_{ik}+\delta _{il}T_{jk}),`$ $`\mathrm{tr}(T_{ij}T_{kl})=2(\delta _{ik}\delta _{jl}\delta _{il}\delta _{jk}).`$ (3.2) For later convenience, we define $`X_i\stackrel{\mathrm{def}}{=}T_{Ni},X_{i}^{}{}_{}{}^{}\stackrel{\mathrm{def}}{=}T_{N1,i^{}},X_{N1}\stackrel{\mathrm{def}}{=}T_{N,N1}(i,i^{}=1,\mathrm{},N2).`$ (3.3) Let us classify the real and complex generators at 1) a symmetric point and 2) a non-symmetric point. 1) Symmetric points. The vacuum vectors of symmetric points can be transformed by a $`G`$-action to $`\stackrel{}{v}=\left(\begin{array}{c}0\\ \mathrm{}\\ 0\\ v\end{array}\right).`$ (3.4) We can immediately find a) the real unbroken algebra and b) the complex unbroken algebra. 1-a) Real unbroken Lie algebra. Hermitian unbroken generators are $`(N1)\times (N1)`$ matrices which act the first $`N1`$ components of (3.4), and others are broken generators: $`𝒢=\left(\begin{array}{cccc}& & & \\ & & & \\ & & & \\ & & & \\ & & & \mathcal{0}\end{array}\right),`$ (3.9) where “$``$” denote Hermitian broken generators. The Hermitian broken generators are of the form $`X_i=T_{Ni}=\left(\begin{array}{cccc}& & & \mathrm{}\\ & \mathrm{}& & i\\ & & & \mathrm{}\\ & & & \\ \mathrm{}& i& \mathrm{}& 0\end{array}\right)𝒢(i=1,\mathrm{},N1),`$ (3.14) where dots denote zero components and only $`i`$-th elements are nonzero. Thus the symmetry breaking pattern at symmetric points is $`G=O(N)H=O(N1)`$. A real target manifold, parametrized by NG bosons for this breaking, is a compact homogeneous manifold $`O(N)/O(N1)S^{N1}`$. 1-b) Complex unbroken Lie algebra. To discuss the full target manifold, we must discuss by the comlexification of $`G`$. By the complexification, however, no new generator appears that leave the vacuum expectation value invariant at the symmetric point. Broken and unbroken generators are simply $`\{\begin{array}{cc}Z_R=X_i𝒢^𝐂\widehat{},\hfill & \\ K_M=H_a\widehat{}.\hfill & \end{array}`$ (3.15) All $`Z_R`$ are Hermitian generators. Chiral superfields, whose bosonic parts are coset coordinates, proportional to Hermitian generators are called the “mixed-type superfield”. A real part of a mixed type superfield is a NG boson while an imaginary part corresponds to a QNG boson. Since numbers of the NG and QNG bosons are both $`N1`$ in this case, this nonlinear realization is called the maximal realization. 2) Non-symmetric points. We discuss the symmetry breaking in general points. We move the vacuum expectation value $`\stackrel{}{v}`$ to $`\stackrel{}{v}^{}`$ by the following element of $`G^𝐂`$: $`g_0`$ $`=`$ $`\mathrm{exp}(i\theta X_{N1})`$ (3.22) $`=`$ $`\left(\begin{array}{ccc}1& 0& \\ & & \multicolumn{-1}{c}{}\\ & \mathrm{cos}\theta & \mathrm{sin}\theta \\ 0& \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)=\left(\begin{array}{ccc}1& 0& \\ & & \multicolumn{-1}{c}{}\\ & \mathrm{cosh}\stackrel{~}{\theta }& i\mathrm{sinh}\stackrel{~}{\theta }\\ 0& i\mathrm{sinh}\stackrel{~}{\theta }& \mathrm{cosh}\stackrel{~}{\theta }\end{array}\right)G^𝐂,`$ where $`\theta \stackrel{\mathrm{def}}{=}i\stackrel{~}{\theta }`$ is a pure imaginary angle. Although we have chosen the rotation by $`X_{N1}`$, rotations by other broken generators are equivalent, since they can be transformed by an action of $`G`$ to each other. Then the vacuum vector at the non-symmetric points can be written, without loss of generality, as $`\stackrel{}{v}^{}=g_0\stackrel{}{v}=\left(\begin{array}{c}0\\ \mathrm{}\\ 0\\ iv\mathrm{sinh}\stackrel{~}{\theta }\\ v\mathrm{cosh}\stackrel{~}{\theta }\end{array}\right)=\left(\begin{array}{c}0\\ \mathrm{}\\ 0\\ \alpha \\ \beta \end{array}\right),`$ (3.23) where we have defined $`\{\begin{array}{cc}\beta \stackrel{\mathrm{def}}{=}v\mathrm{cosh}\stackrel{~}{\theta }:\mathrm{real}\hfill & \\ \alpha \stackrel{\mathrm{def}}{=}iv\mathrm{sinh}\stackrel{~}{\theta }:\mathrm{pure}\mathrm{imaginary},\hfill & \end{array}`$ (3.24) and these satisfy $`\beta ^2+\alpha ^2=v^2`$. By these constants, $`g_0`$ and its inverse can be written as $`g_0=\left(\begin{array}{ccc}1& 0& \\ & & \\ & \beta /v& \alpha /v\\ 0& \alpha /v& \beta /v\end{array}\right),g_{0}^{}{}_{}{}^{1}=\left(\begin{array}{ccc}1& 0& \\ & & \\ & \beta /v& \alpha /v\\ 0& \alpha /v& \beta /v\end{array}\right).`$ (3.31) The magnitudes of vacuum expectation values are $`\stackrel{}{v}^{}\stackrel{}{v}^{}=\beta ^2\alpha ^2=\beta ^2+\stackrel{~}{\alpha }^2\stackrel{\mathrm{def}}{=}v_{}^{}{}_{}{}^{2},`$ (3.32) $`\stackrel{}{v}^2=\stackrel{}{v}^2=\beta ^2+\alpha ^2=\beta ^2\stackrel{~}{\alpha }^2=v^2.`$ (3.33) We can find from Eq. (3.23) that the non-compact directions are hyperbolic as in Fig. 3. The full target space is a spheroidal hyperboloid and the compact coset $`G/H`$ is embedded in the symmetric point. (The compact coset $`G/H^{}`$, at non-symmetric point, is discussed in the next subsection.) Let us discuss real and complex Lie algebras at non-symmetric points. 2-a) Real Lie algebra. At non-symmetric points, the whole generators can be divided into real unbroken algebra $`^{}`$ and real broken generators as $`𝒢=\left(\begin{array}{ccccc}& & & & \\ & ^{}& & & \\ & & & & \\ & & & & \\ & & & \mathcal{0}& \\ & & & & \mathcal{0}\end{array}\right),`$ (3.39) where “$``$” denote broken generators. The broken generators can be written explicitly as $`X_{i}^{}{}_{}{}^{}=\left(\begin{array}{ccccc}& & & \mathrm{}& \mathrm{}\\ & \mathrm{}& & i& 0\\ & & & \mathrm{}& \mathrm{}\\ & & \multicolumn{-1}{c}{}& & \\ \mathrm{}& i& \mathrm{}& 0& 0\\ \mathrm{}& 0& \mathrm{}& 0& 0\end{array}\right),X_i=\left(\begin{array}{ccccc}& & & \mathrm{}& \mathrm{}\\ & \mathrm{}& & 0& i\\ & & & \mathrm{}& \mathrm{}\\ & & \multicolumn{-1}{c}{}& & \\ \mathrm{}& 0& \mathrm{}& 0& 0\\ \mathrm{}& i& \mathrm{}& 0& 0\end{array}\right),`$ (3.50) $`X_{N1}=\left(\begin{array}{ccccc}& & & & \\ & \mathrm{}& & \mathrm{}& \mathrm{}\\ & & & & \\ & & \multicolumn{-1}{c}{}& & \\ & \mathrm{}& & 0& i\\ & \mathrm{}& & i& 0\end{array}\right)𝒢^{}(i=1,\mathrm{},N2),`$ (3.56) where dots denote zero components, and non-zero elements in the first two equations are $`i`$-th components. The symmetry breaking pattern at non-symmetric points turns out $`G=O(N)H^{}=O(N2)`$, which is smaller than symmetric points. A real target manifold, parametrized by NG bosons, is a compact homogeneous manifold $`G/H^{}=O(N)/O(N2)`$, which is larger than one of the symmetric point, $`G/H`$. (See Fig. 3.) Namely we have more NG bosons at non-symmetric points than at symmetric points. These newly emerged NG bosons must come from the QNG bosons, since the dimension of the full target manifold has to be unchanged. There is only one QNG boson because the number of the NG bosons is $`2N3`$ (and the total number is $`2N2`$). In the next subsection, we show how these different compact coset manifolds are embedded in the full manifold and how some of the QNG bosons change to the NG bosons at non-symmetric points. Before doing it, we investigate the complex symmetry at non-symmetric points, which give us the key point to understand such phenomena. 2-b) Complex Lie algebra. Complex broken and unbroken generators at non-symmetric points $`\stackrel{}{v}^{}`$ can be immediately calculated by using Eqs. (2.4) and (2.3), to yield $`Z_{R}^{}{}_{}{}^{}`$ $`=`$ $`g_0Z_Rg_0^1`$ (3.57) $`=`$ $`\{\begin{array}{cc}g_0X_ig_0^1=\frac{\alpha }{v}X_{i}^{}{}_{}{}^{}+\frac{\beta }{v}X_i\stackrel{\mathrm{def}}{=}Z_{I}^{}{}_{}{}^{}\hfill & \\ g_0X_{N1}g_0^1=X_{N1}\stackrel{\mathrm{def}}{=}Z_{N1}^{}\hfill & \end{array}𝒢^𝐂\widehat{}^{},`$ $`K_{M}^{}{}_{}{}^{}`$ $`=`$ $`g_0K_Mg_0^1`$ (3.58) $`=`$ $`\{\begin{array}{cc}g_0X_i^{}g_0^1=\frac{\beta }{v}X_{i}^{}{}_{}{}^{}\frac{\alpha }{v}X_i\stackrel{\mathrm{def}}{=}B_{I}^{}{}_{}{}^{}\hfill & \\ g_0H_{a}^{}{}_{}{}^{}g_0^1=H_{a}^{}{}_{}{}^{}^{}\hfill & \end{array}\widehat{}^{},`$ where $`Z_{I}^{}{}_{}{}^{}`$ and $`B_{I}^{}{}_{}{}^{}`$ $`(I^{}=1,\mathrm{},N2)`$ can be explicitly written as $`Z_{I}^{}{}_{}{}^{}=\left(\begin{array}{ccccc}& & & \mathrm{}& \mathrm{}\\ & \mathrm{}& & i\alpha /v& i\beta /v\\ & & & \mathrm{}& \mathrm{}\\ & & \multicolumn{-1}{c}{}& & \\ \mathrm{}& i\alpha /v& \mathrm{}& 0& 0\\ \mathrm{}& i\beta /v& \mathrm{}& 0& 0\end{array}\right),`$ (3.64) $`B_{I}^{}{}_{}{}^{}=\left(\begin{array}{ccccc}& & & \mathrm{}& \mathrm{}\\ & \mathrm{}& & i\beta /v& i\alpha /v\\ & & & \mathrm{}& \mathrm{}\\ & & \multicolumn{-1}{c}{}& & \\ \mathrm{}& i\beta /v& \mathrm{}& 0& 0\\ \mathrm{}& i\alpha /v& \mathrm{}& 0& 0\end{array}\right).`$ (3.70) We can classify these broken generators to pure-types or mixed-types as follows. First of all, the broken generator $`Z_{N1}^{}{}_{}{}^{}`$ corresponds to a mixed-type superfield, since $`Z_{N1}^{}{}_{}{}^{}`$ is a Hermitian generator. A real and a imaginary parts of a scalar component of a chiral superfield generated by $`Z_{N1}^{}{}_{}{}^{}`$ is a NG boson and a QNG boson, respectively. On the other hand, all other generators $`Z_{I}^{}{}_{}{}^{}`$ generate pure-type chiral superfields, where both scalar components are NG bosons, since they are non-Hermitian generators. We can count numbers of the NG and QNG bosons as $`2N3`$ and $`1`$, respectively, without using the fact that the total number of the NG and QNG bosons does not change. ### 3.2 Embedding of NG cosets $`G/H`$ and $`G/H^{}`$ In this subsection, we show how the different cosets 1) $`G/H`$ and 2) $`G/H^{}`$ are embedded into symmetric and non-symmetric points by using the Shore’s procedure . We can obtain coset representatives of NG bosons $`G/H`$ ($`G/H^{}`$) by putting all QNG bosons zero in the complex representative $`\xi `$ of the full complex coset $`G^𝐂/\widehat{H}`$, at symmetric (non-symmetric) points. In the cases when there are pure-type superfields, we need a local $`\widehat{H}`$-transformation from the right. 1) Embedding of $`G/H`$ at symmetric points. Since we do not need a local $`\widehat{H}`$-transformation, we can obtain the representative of $`G/H`$ by simply putting all QNG bosons zero : $`\xi \|_{\mathrm{QNG}=B^i=0}=e^{i\varphi X}G/H,`$ (3.71) where fields $`\varphi =\{A^i\}(i=1,\mathrm{},N1)`$ are NG bosons at symmetric points. 2) Embedding of $`G/H^{}`$ at non-symmetric points. The representative of the complex coset at non-symmetric points is $`\xi (\phi )=e^{i\phi Z^{}}=\mathrm{exp}i\left[\phi ^i\left({\displaystyle \frac{\beta }{v}}X_i+{\displaystyle \frac{\alpha }{v}}X_{i}^{}{}_{}{}^{}\right)+\phi ^{N1}X_{N1}\right]G^𝐂/\widehat{H}^{},`$ (3.72) where we have used a character $`\phi `$ as a coordinate. Since there are pure-type broken generators $`Z_{I}^{}{}_{}{}^{}`$, we need a local $`\widehat{H}`$-transformation $`\zeta ^{}(\phi ,\phi ^{})=\mathrm{exp}(id(\phi ,\phi ^{})B^{})\widehat{H}^{}`$ (3.73) from the right: $`\xi (\phi )`$ $`\widehat{\xi }(\widehat{A},\widehat{B})`$ $`=`$ $`\xi (\phi )\zeta ^{}(\phi ,\phi ^{})`$ (3.74) $`=`$ $`\mathrm{exp}i\left[\widehat{\phi }^i\left({\displaystyle \frac{\beta }{v}}X_i+{\displaystyle \frac{\alpha }{v}}X_{i}^{}{}_{}{}^{}\right)+\widehat{d}^i(\widehat{\phi },\widehat{\phi }^{})\left({\displaystyle \frac{\alpha }{v}}X_i+{\displaystyle \frac{\beta }{v}}X_{i}^{}{}_{}{}^{}\right)+\widehat{\phi }^{N1}X_{N1}\right]`$ $`=`$ $`\mathrm{exp}i\left[a^iX_i+b^iX_{i}^{}{}_{}{}^{}+(\widehat{A}^{N1}+i\widehat{B}^{N1})X_{N1}\right],`$ where $`\widehat{\phi }^i=\widehat{A}^i+i\widehat{B}^i`$ are transformed fields whose relation to $`\phi `$ is obtained below, Eq. (3.80), $`\widehat{d}^i`$ is a function of $`\widehat{\phi }`$ and $`\widehat{\phi }^{}`$ whose relation to $`d^i`$ is also obtained below, Eq. (3.79), and $`a^i`$ and $`b^i`$ are scalar fields. We can chose the function $`\widehat{d}`$ (or $`d`$) such that scalar fields $`a`$ and $`b`$ become real: $`\widehat{d}^i(\widehat{\phi },\widehat{\phi }^{})={\displaystyle \frac{i}{2\stackrel{~}{\alpha }\beta }}(v^2\widehat{\phi }^iv^2\widehat{\phi }^i),a^i=\widehat{A}^i{\displaystyle \frac{v}{\beta }},b^i=\widehat{B}^i{\displaystyle \frac{v}{\stackrel{~}{\alpha }}}.`$ (3.75) Since real scalar fields $`\widehat{A}^i`$ and $`\widehat{B}^i`$ are proportional to Hermitian generators $`X_i`$ and $`X_i^{}`$, respectively in the exponential, they parameterize compact directions of the target manifold. This is why both $`\widehat{A}^i`$ and $`\widehat{B}^i`$ are NG bosons, and $`\widehat{\mathrm{\Phi }}^i`$ can be considered pure-type superfields. On the other hand, since $`\widehat{A}^{N1}`$ and $`\widehat{B}^{N1}`$ are proportional to Hermitian and anti-Hermitian generators on the exponential, they parameterize compact and non-compact directions of the target space, respectively. Hence $`\widehat{A}^{N1}`$ and $`\widehat{B}^{N1}`$ are NG and QNG bosons, and then $`\widehat{\mathrm{\Phi }}^{N1}`$ is a mixed-type superfield. Then the superfields $`\mathrm{\Phi }^i`$ and $`\mathrm{\Phi }^{N1}`$, before the $`\widehat{H}`$-transformation, turn out to be pure-type and mixed-type superfields, respectively, since they coincide with $`\widehat{\mathrm{\Phi }}^i`$ and $`\widehat{\mathrm{\Phi }}^{N1}`$ at the linear order, as shown below, Eq. (3.80). We can obtain the real representative of the coset $`G/H^{}`$ by putting all QNG bosons zero: $`\widehat{\xi }\|_{\mathrm{QNG}=\widehat{B}^{N1}=0}=e^{i\varphi ^{}X}G/H^{},`$ (3.76) where fields $`\varphi ^{}=\{a^i,b^i,\widehat{A}^{N1}\}`$ are NG bosons at non-symmetric points. To understand why the compact coset manifold $`G/H^{}`$ at non-symmetric points is larger than $`G/H`$ at symmetric points, see Figs. 4 and 5. At symmetric points, there are $`N1`$ non-compact directions, while they change to one non-compact direction and $`N2`$ compact directions parametrized by newly emerged NG bosons.<sup>8</sup><sup>8</sup>8 We give a comment on the NG submanifold at non-symmetric points, $`G/H^{}`$. It can be considered as a $`H/H^{}S^{N2}`$ fiber bundle over a base manifold, $`G/HS^{N1}`$. By bringing $`\stackrel{}{v}^{}`$ to $`\stackrel{}{v}`$, the fiber shrinks but the base remains at finite size (radius $`v`$). Next we obtain the relation of the fields $`\phi `$ and $`\widehat{\phi }`$ (or $`\mathrm{\Phi }`$ and $`\widehat{\mathrm{\Phi }}`$). In general, the first equation of Eq. (3.74) can be written explicitly as $`e^{i\phi Z}e^{id(\phi ,\phi ^{})K}=e^{i(\widehat{\phi }Z+\widehat{d}(\widehat{\phi },\widehat{\phi }^{})K)},\zeta ^{}(\phi ,\phi ^{})=e^{id(\phi ,\phi ^{})K}\widehat{H}^{},`$ (3.77) where $`d`$ is a function of $`\phi `$ and $`\phi ^{}`$. By using the Baker-Campbell-Hausdorff formula on the left-side, we obtain relations $`\widehat{\phi }^R=\phi ^R+{\displaystyle \frac{1}{2}}f_{MS}^{}{}_{}{}^{R}\phi ^Sd^M+\mathrm{},`$ (3.78) $`\widehat{d}^M=d^M+{\displaystyle \frac{1}{12}}f_{NR}^{}{}_{}{}^{T}f_{TS}^{}{}_{}{}^{M}\phi ^R\phi ^Sd^N+\mathrm{}.`$ (3.79) Note that we have used only the fact that $`G^𝐂/\widehat{H}`$ is a symmetric space, and the result is model-independent. By these equations, two coordinates $`\widehat{\phi }`$ and $`\phi `$ are related by $`\widehat{\phi }=\phi +O(\phi ^2,\phi \phi ^{}),`$ (3.80) and coincide to each other at the first order.<sup>9</sup><sup>9</sup>9 Note that this transformation is not holomorphic and so the complex structures of the two coordinates $`\phi `$ and $`\widehat{\phi }`$ are different. For our purpose to calculate low-energy theorems, the difference can be neglected because the curvature tensors in both coordinates coincide at $`\phi =0`$. This is why $`\mathrm{\Phi }`$ and $`\widehat{\mathrm{\Phi }}`$ coincide to each other at the linear level and their identifications to pure- or mixed-type superfields coincide. ## 4 Low-energy theorems at general points In this section, we discuss low-energy theorems at non-symmetric points. We elaborate on the $`O(N)`$ model as an example, but generalizations to other models are straightforward. ### 4.1 Some formula for the $`O(N)`$ model Before discussing low-energy theorems, we give comments for a “linear” description of the model. The invariant Lagrangian can be written as $`={\displaystyle }d^4\theta \stackrel{}{\varphi }^{}\stackrel{}{\varphi }+({\displaystyle }d^2\theta W(\varphi )+(\mathrm{conj}.)),`$ (4.1) where $`\stackrel{}{\varphi }(x,\theta ,\overline{\theta })`$ consists of chiral superfields belonging to a linear representation of $`G`$. $`W(\varphi )`$ is a $`G`$-invariant superpotential, and it is actually $`G^𝐂`$-invariant due to its holomorphy. For the $`O(N)`$ model, its candidate is $`W(\varphi )=g\varphi _0(\stackrel{}{\varphi }^2a^2),`$ (4.2) where $`\varphi _0`$ is an additional $`G`$-singlet field and $`g`$ is a coupling constant. $`\varphi _0`$ tends to a non-dynamical auxiliary field in the heavy mass limit ($`g\mathrm{}`$). In this limit we can eliminate $`\varphi _0`$ by its equation of motion, which gives an F-term constraint among dynamical fields $`\varphi ^i`$. For the $`O(N)`$ model, the constraint is $`\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}a^2=0`$. The Kähler potential may suffer from a quantum correction, with preserving the global symmetry $`G`$. We thus obtain a nonlinear Kähler potential for NG chiral superfields, $`K(\mathrm{\Phi },\mathrm{\Phi }^{})=f(\stackrel{}{\varphi }^{}\stackrel{}{\varphi })\|_\mathrm{F}=f(\stackrel{}{v}^{}\xi ^{}\xi \stackrel{}{v}),`$ (4.3) where F denotes an F-term constraint. We have used a relation between linear superfields and NG superfields, $`\stackrel{}{\varphi }\|_\mathrm{F}=\xi \stackrel{}{v}`$. This recovers Eq. (1.15). If we restrict the problem to the $`O(N)`$ model, geometric quantities can be calculated by solving the constraint $`\stackrel{}{\varphi }^2=a^2`$ explicitly as $`\varphi ^N=\sqrt{a^2_{i=1}^{N1}(\varphi ^i)^2}`$. However we discuss in the coset formalism which can be generalized to other models straightforwardly. To obtain geometric quantities in the coset formalism, we need expectation values of broken generators, sandwiched by the vacuum vector $`\stackrel{}{v}^{}`$. Let us calculate them first. We use indices $`R,S,\mathrm{}=1,\mathrm{},N1`$ and $`I,J,\mathrm{}=1,\mathrm{},N2`$. We omit primes except for $`\stackrel{}{v}^{}`$. By noting that $`Z_{N1}^{}{}_{}{}^{}=Z_{N1}`$ and $`Z_{I}^{}{}_{}{}^{}Z_I`$, we calculate products of one complex generator on vacuum expectation values, given by $`Z_I\stackrel{}{v}^{}=\left(\begin{array}{c}0\\ \mathrm{}\\ iv\\ \mathrm{}\\ 0\end{array}\right),Z_{N1}\stackrel{}{v}^{}=\left(\begin{array}{c}0\\ \mathrm{}\\ \mathrm{}\\ i\beta \\ i\alpha \end{array}\right),Z_{I}^{}{}_{}{}^{}\stackrel{}{v}^{}={\displaystyle \frac{v_{}^{}{}_{}{}^{2}}{v^2}}\left(\begin{array}{c}0\\ \mathrm{}\\ iv\\ \mathrm{}\\ 0\end{array}\right),`$ (4.4) where the $`I`$-th elements are nonzero in the first and the third equations. Products of two complex generators on vacuum expectation values are also given by $`Z_RZ_S\stackrel{}{v}^{}=\delta _{RS}\stackrel{}{v}^{},Z_{I}^{}{}_{}{}^{}Z_J\stackrel{}{v}^{}=\delta _{I^{}J}\stackrel{}{v}^{},`$ $`Z_{I}^{}{}_{}{}^{}Z_{N1}\stackrel{}{v}^{}={\displaystyle \frac{c^2}{v_{}^{}{}_{}{}^{2}}}Z_{I}^{}{}_{}{}^{}\stackrel{}{v}^{},Z_{N1}Z_{I}^{}{}_{}{}^{}\stackrel{}{v}^{}=0,`$ (4.5) where we have defined $`c^2\stackrel{\mathrm{def}}{=}2\sqrt{v^{\mathrm{\hspace{0.17em}4}}v^4}=2\stackrel{~}{\alpha }\beta `$. We define convenient notations $`<R\mathrm{}>\stackrel{\mathrm{def}}{=}\stackrel{}{v}^{}Z_R\mathrm{}\stackrel{}{v}^{},R^{}\stackrel{\mathrm{def}}{=}Z_{R}^{}{}_{}{}^{}.`$ (4.6) Then expectation values of one to four generators can be calculated, to yield $`<I>=0,<N1>=c^2,`$ $`<RS>=v_{}^{}{}_{}{}^{2}\delta _{RS},<I^{}J>=v^2\delta _{I^{}J},`$ $`<IJK>=0,<I^{}JK>=0,`$ $`<N1,IJ>=c^2\delta _{IJ},<IJ,N1>=0,`$ $`<I,N1>=0,<,N1,I>=0,`$ $`<I,N1,N1>=0,<I^{},N1,N1>=0,`$ $`<N1,I,N1>=0,<N1,N1,N1>=c^2,`$ $`<R^{}S^{}UV>=v_{}^{}{}_{}{}^{2}\delta _{R^{}S^{}}\delta _{UV},`$ (4.7) where “$``$” denotes an arbitrary generator. These quantities are needed for the calculation of the curvature tensor. They can be generalized to other models straightforwardly. ### 4.2 Geometric quantities and low-energy theorems We can calculate geometric quantities of the $`O(N)`$ model by using the formulas obtained in the last subsection. In this section, we consider the most simple Kähler potential, $`K=f(x)=x`$. The general case is discussed in Appendix B. First of all the metric is given by (we omit prime on $`Z_R`$) $`g_{ij^{}}=_i_j^{}K=G_{RS^{}}E_i^R(E_j^S)^{},G_{RS^{}}=\stackrel{}{v}^{}Z_{S}^{}{}_{}{}^{}\xi ^{}\xi Z_R\stackrel{}{v}^{},`$ (4.8) where $`G_{RS^{}}`$ is called an auxiliary metric. The auxiliary metric at the point $`\phi =0`$ becomes $`G_{RS^{}}\|_{\phi =0}=<S^{}R>=\left(\begin{array}{cc}v^2\delta _{IJ^{}}& 0\\ 0& v_{}^{}{}_{}{}^{2}\end{array}\right).`$ (4.9) A vielbein and a $`\widehat{H}`$-connection at the point $`\phi =0`$ are given by $`E_i^R\|_{\phi =0}=\delta _i^R,W_i^M\|_{\phi =0}=0,`$ (4.10) respectively, and the differentiation of the vielbein with respect to the coordinate can be calculated, to yield $`_jE_i^R\|_{\phi =0}=0.`$ (4.11) Let us calculate a complex curvature $`R_{ij^{}kl^{}}=_i_j^{}_k_l^{}Kg^{mn^{}}(_j^{}_m_l^{}K)(_i_k_n^{}K),`$ (4.12) which is crucial to low-energy theorems. The complex curvature on the origin $`\phi =0`$ of $`G^𝐂/\widehat{H}`$ can be calculated by Eq. (4.7), to yield $`R_{ij^{}kl^{}}\|_{\phi =0}`$ $`=`$ $`[<S^{}V^{}UR>G^{XY^{}}\|_{\phi =0}<X^{}SV>^{}<Y^{}RU>]\delta _i^R(\delta _j^S)^{}\delta _k^U(\delta _l^V)^{}`$ (4.13) $`=`$ $`{\displaystyle \frac{v^4}{v_{}^{}{}_{}{}^{2}}}\delta _{ik}\delta _{j^{}l^{}}.`$ The first line is for general symmetric manifolds $`G^𝐂/\widehat{H}`$ with one vacuum expectation value, and the second line is for the $`O(N)`$ model. If we rescale fields so that the metric (4.9) becomes the Knonecker’s delta, components of the curvature tensor become $`R_{ij^{}kl^{}}\|_{\phi =0}=\{\begin{array}{cc}1/v_{}^{}{}_{}{}^{2}\delta _{ik}\delta _{j^{}l^{}}\text{ when }i,j,k,l=1,\mathrm{},N2,\hfill & \\ v^2/v_{}^{}{}_{}{}^{4}\delta _{ik}\delta _{j^{}l^{}}\text{ when only two indices are }N1\text{-th},\hfill & \\ v^4/v_{}^{}{}_{}{}^{6}\delta _{ik}\delta _{j^{}l^{}}\text{ when }i,j,k,l=N1.\hfill & \end{array}`$ (4.14) From these equations, we can calculate real components of the curvature in the rescaled coordinate given by $`(i,j,k,l=1,\mathrm{},N2)`$ $`R_{A^iA^jA^kA^l}=2{\displaystyle \frac{1}{v_{}^{}{}_{}{}^{2}}}(\delta _{ik}\delta _{jl}\delta _{il}\delta _{jk}),`$ $`R_{B^iA^jB^kA^l}=2{\displaystyle \frac{1}{v_{}^{}{}_{}{}^{2}}}(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}),`$ $`R_{B^iA^jA^kA^l}=R_{A^iA^jB^kA^l}=0,`$ $`R_{A^{N1}A^jA^{N1}A^l}=2{\displaystyle \frac{v^2}{v_{}^{}{}_{}{}^{4}}}\delta _{jl},`$ $`R_{A^{N1}A^{N1}A^kA^l}=0,`$ $`R_{B^{N1}A^jB^{N1}A^l}=2{\displaystyle \frac{v^2}{v_{}^{}{}_{}{}^{4}}}\delta _{jl},`$ $`R_{B^{N1}A^{N1}B^kA^l}=0,`$ $`R_{B^{N1}A^{N1}B^{N1}A^{N1}}=4{\displaystyle \frac{v^4}{v_{}^{}{}_{}{}^{6}}},`$ (4.15) where all quantities are evaluated at $`\phi =0`$. Let us discuss low-energy theorems. As discussed in Appendix A, components of the curvature tensor in an arbitrary coordinate and in normal coordinates coincide to each other in this order. Hence, by substituting these equations to Eq. (1.7), we can obtain low-energy theorems for two-body scattering amplitudes among NG and QNG bosons at general points of target spaces. At the symmetric point, $`v^{}=v`$, all coefficients become $`\frac{1}{v^2}`$. There, all $`A^i`$ and $`A^{N1}`$ fields are NG bosons of the symmetry breaking, $`O(N)`$ to $`O(N1)`$, and we can verify that their scattering amplitudes satisfy low-energy theorems from a equation, $`f_{ij}^{}{}_{}{}^{a}f_{akl}=(\delta _{ik}\delta _{jl}\delta _{il}\delta _{jk})`$. Their decay constant is $`f_\pi =v`$. All $`B`$ fields correspond to QNG bosons and their low-energy theorems coincide with those of NG partners as discussed in Sec. 1 and Ref . At the non-symmetric point, there appear new features of low-energy theorems for NG and QNG bosons. The unbroken symmetry $`O(N1)`$ at symmetric points further breaks down to $`O(N2)`$, and $`B^i(i=1,\mathrm{},N2)`$ change to NG bosons for the second breaking of $`O(N1)`$ to $`O(N2)`$. From the first equation of (4.15) and $`R_{B^iB^jB^kB^l}=R_{A^iA^jA^kA^l}`$, we can verify that low-energy theorems for $`B^i(i=1,\mathrm{},N2)`$ coincide with those of NG bosons of the second symmetry breaking, $`O(N1)`$ to $`O(N2)`$. Low-energy theorems among $`A^i`$ and $`A^N`$ (at symmetric points) for the first breaking of $`O(N)`$ to $`O(N1)`$ are distorted there, since the field $`A^{N1}`$ becomes to play a special role. Before closing this section, we give a comment on a relation between NG bosons at non-symmetric points in a supersymmetric theory and NG bosons in a non-supersymmetric theory. In a non-supersymmetric theory with spontaneously broken $`O(N)`$ symmetry to a subgroup $`O(N2)`$, two sets of linear fields $`\stackrel{}{\varphi }_1`$ and $`\stackrel{}{\varphi }_2`$, belonging to the vector representation, should have vacuum expectation values. The broken generators are $`𝒢=\left(\begin{array}{ccccc}& & & & \\ & 0& & f_\pi ^{(1)}& f_\pi ^{(2)}\\ & & & & \\ & & & & \\ & f_\pi ^{(1)}& & & f_\pi ^{(3)}\\ & & & & \\ & f_\pi ^{(2)}& & f_\pi ^{(3)}& \end{array}\right),`$ (4.21) where we have expressed three $`H`$-irreducible sectors of broken generators by decay constants $`f_\pi ^{(i)}(i=1,2,3)`$ of NG bosons corresponding to these generators.<sup>10</sup><sup>10</sup>10 Since there are two vacuum expectation values $`\stackrel{}{v}_1=<\stackrel{}{\varphi }_1>,\stackrel{}{v}_2=<\stackrel{}{\varphi }_2>𝐍`$ to break $`O(N)`$ to $`O(N2)`$, these three sectors correspond to three $`G`$-invariants, $`\stackrel{}{v}_{1}^{}{}_{}{}^{\mathrm{\hspace{0.17em}2}},\stackrel{}{v}_{2}^{}{}_{}{}^{\mathrm{\hspace{0.17em}2}},\stackrel{}{v}_1\stackrel{}{v}_2`$. These three free parameters in the non-supersymmetric theory are reduced to two parameters, $`v`$ and $`v^{}`$, to be embedded into a bosonic part of a supersymmetric theory. This is because there exist $`N2`$ pure-type multiplets and they relate two decay constants, $`f_\pi ^{(1)}`$ and $`f_\pi ^{(2)}`$, in Eq. (4.21). This was known at least in pure-realization cases, where there exist only pure-type multiplets . ## 5 Conclusion and discussion If a global symmetry spontaneously breaks in supersymmetric theories, there appear NG and QNG bosons and their fermions superpartners. The low-energy effective Lagrangian for these fields can be constructed as supersymmetric nonlinear sigma models. If symmetry breaking occur by a superpotential of a fundamental of an effective field theories, there must appear at least one QNG boson. Hence the target manifold inevitably becomes non-compact. As a result, supersymmrtic vacuum alignment occur; NG and QNG boson can change with the total number preserved. This has been understood by different embedding of $`\widehat{H}`$ into $`G^𝐂`$. Low-energy theorems of two-body scattering amplitudes for these bosons was known at symmetric points. In this paper, we have calculated low-energy theorems for NG and QNG bosons at general points. We have found new features of low-energy theorems. In a theory with the supersymmetric vacuum alignment, symmetry breaking occurs twice (or more times for other models). The low-energy theorems for the first breaking at symmetric points have been distorted at non-symmetric points; on the other hand, low-energy theorems for the second breaking coincide with non-supersymmetric cases. This is because one (or some for other models) NG boson must sit in mixed-type multiplet, and play a special role in the sense that its partner is QNG boson. Although we have illustrated the low-energy theorems at non-symmetric points in the $`O(N)`$-model with the simplest (linear) Kähler potential, generalizations to more complicated models are straightforward. (The calculation in the most general Kähler potential of the $`O(N)`$-model is discussed in Appendix B.) Consider the case that there are $`n`$ $`G^𝐂`$-invariants and $`m`$ $`G`$-invariants. (In the $`O(N)`$-model, they are $`m=n=1`$, since there is one $`G^𝐂`$-invariant, $`\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}`$, and one $`G`$-invariant, $`\|\stackrel{}{\varphi }\|^2`$.) The low-energy effective Kähler potential can be written as an arbitrary function of $`m`$ $`G`$-invariants . We can count parameters included in low-energy theorems of two-body scattering amplitudes. Since curvature tensor includes from one to four derivatives of the arbitrary function, there are certain numbers of the parameters concerned with the arbitrary function. Therefore the low-energy theorems include these parameters. As seen in this paper, in the case of $`O(N)`$ model, there were six parameters $`v,v^{},f_1,f_2,f_3,f_4`$ (see Appendix B). Although we have investigated two-body scattering amplitudes, a generalization to many-body scattering amplitudes can be calculated by using the Kähler normal coordinate to the desired order . An interaction Lagrangian can be written by the curvature tensor, covariant derivatives of the curvature tensor etc. If we calculate $`n`$-body scattering amplitudes, it contains from one to $`n`$ derivatives of the arbitrary function. We have investigated only low-energy theorems, namely low-energy scattering amplitudes at the leading order $`𝒪(p^2)`$. It is interesting, for a development of supersymmetric chiral perturbation theories, to investigate higher derivative terms such as next-to leading terms $`𝒪(p^4)`$ . At such order, we need a supersymmetric Wess-Zumino-Witten term , which correctly reproduces anomalies of the global symmetry (we did not need it at the lowest order). ## Acknowledgements We are grateful to Kazutoshi Ohta and Nobuyoshi Ohta for the arguments in the early stage of this work. The work of M. N. is supported in part by JSPS Research Fellowships. ## Appendix A Kähler normal coordinate expansion In this appendix, we show that the Kähler normal coordinate can be always obtained from general coordinates by a holomorphic coordinate transformation preserving the curvature tensor up to a constant order. The systematic method to obtain the Kähler normal coordinate to an arbitrary order is discussed in Ref. . Let $`\{z^i,z^i\}`$ are the general coordinate. We expand the Kähler potential by the Taylor expansion as $`K(z,z^{})`$ $`=`$ $`K\|_0+F(z)+F^{}(z^{})`$ (A.1) $`+g_{ij^{}}\|_0z^iz^j+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{i^{}jk}\|_0z^iz^jz^k+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{ij^{}k^{}}\|_0z^iz^jz^k`$ $`+{\displaystyle \frac{1}{4}}(R_{ij^{}kl^{}}+g_{mn^{}}\mathrm{\Gamma }_{}^{m}{}_{ik}{}^{}\mathrm{\Gamma }_{}^{n^{}}{}_{j^{}l^{}}{}^{})\|_0z^iz^kz^jz^l`$ $`+{\displaystyle \frac{1}{6}}_k\mathrm{\Gamma }_{l^{}ij}\|_0z^iz^jz^kz^l+{\displaystyle \frac{1}{6}}_k^{}\mathrm{\Gamma }_{li^{}j^{}}\|_0z^iz^jz^kz^l+O(z^5),`$ where $`F(z)=_iK\|z^i+{\displaystyle \frac{1}{2}}_i_jK\|z^iz^j+\mathrm{}`$ (A.2) is holomorphic and can be eliminated by a Kähler transformation. Here $`\mathrm{\Gamma }_{}^{i}{}_{jk}{}^{}`$ is the connection and $`R_{ij^{}kl^{}}`$ is the curvature tensor of the Kähler manifold . There are many non-covariant coefficients except for $`g_{ij^{}}`$ and $`R_{i^{}jk^{}l}`$. To eliminate them note that Eq. (A.1) can be written as $`K(z,z^{})`$ $`=K\|_0+F(z)+F^{}(z^{})`$ $`+g_{mn^{}}(z^m+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{}^{i}{}_{jk}{}^{}\|z^jz^k+{\displaystyle \frac{1}{6}}_k\mathrm{\Gamma }_{}^{m}{}_{ij}{}^{}\|z^iz^jz^k)(z^n+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{}^{i}{}_{jk}{}^{}\|z^jz^k+{\displaystyle \frac{1}{6}}_k\mathrm{\Gamma }_{}^{n}{}_{ij}{}^{}\|z^iz^jz^k)^{}`$ $`+{\displaystyle \frac{1}{4}}R_{i^{}jk^{}l}\|z^iz^kz^jz^l+O(z^5).`$ (A.3) By a holomorphic coordinate transformation $`\omega ^i=z^i+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{}^{i}{}_{jk}{}^{}\|z^jz^k+{\displaystyle \frac{1}{6}}_l\mathrm{\Gamma }_{}^{i}{}_{jk}{}^{}\|z^jz^kz^l+O(z^4),`$ (A.4) it can be rewritten as $`K(\omega ,\omega ^{})=K\|+\stackrel{~}{F}(\omega )+\stackrel{~}{F}^{}(\omega ^{})+g_{ij^{}}\|\omega ^i\omega ^j+{\displaystyle \frac{1}{4}}R_{i^{}jk^{}l}\|\omega ^i\omega ^k\omega ^j\omega ^l+O(\omega ^5),`$ (A.5) where $`\stackrel{~}{F}(\omega )\stackrel{\mathrm{def}}{=}F(z(\omega ))`$. The new coordinate $`\omega `$ is the Kähler normal coordinate to the forth order. In this coordinate, all coefficients are covariant quantities. By performing these transformations in the superfield level, we obtain a Kähler normal coordinate expansion of the Lagrangian given by $``$ $`=`$ $`g_{ij^{}}\|_{\phi =0}_\mu \phi ^i^\mu \phi ^j+ig_{ij^{}}\|_{\phi =0}\overline{\psi }^j\overline{\sigma }^\mu _\mu \psi ^i`$ (A.6) $`+`$ $`R_{ij^{}kl^{}}\|_{\phi =0}\phi ^k\phi ^l_\mu \phi ^i^\mu \phi ^j+{\displaystyle \frac{1}{4}}R_{ij^{}kl^{}}\|_{\phi =0}\psi ^i\psi ^k\overline{\psi }^j\overline{\psi }^l`$ $`+`$ $`iR_{ij^{}kl^{}}\|_{\phi =0}\phi ^j_\mu \phi ^i(\overline{\psi }^l\overline{\sigma }^\mu \psi ^k).`$ which has been used to calculate the low-energy theorems. First two terms are equation terms for bosons and fermions, and others can be considered interaction terms. (Although we concentrate on bosonic amplitudes in this paper, we can also obtain low-energy theorems including fermion partners of NG and QNG bosons by using this expansion.) Next let us discuss a relation between the curvature tensor in an arbitrary coordinate and in the Kähler normal coordinate. Since the Jacobian of (A.4) $`J_{}^{i}{}_{j}{}^{}=\omega ^i/z^j=\delta _j^i+O(z)`$ (A.7) is unit matrix up to constant order, components of the curvature tensor in the new coordinate $`\omega `$ is $`R_{}^{}{}_{m^{}no^{}p}{}^{}=R_{i^{}jk^{}l}(J_{}^{i}{}_{m}{}^{})^{}J_{}^{j}{}_{n}{}^{}(J_{}^{k}{}_{o}{}^{})^{}J_{}^{l}{}_{p}{}^{}=R_{i^{}jk^{}l}+O(\omega ).`$ (A.8) Therefore components of the curvature tensor is invariant up to constant order. We can use an arbitrary coordinate to calculate the curvature tensor in low-energy theorems (1.7), although low-energy theorems theirselves have been obtained in normal coordinates. ## Appendix B General Kähler potential In Sec. 4, we have calculated geometric quantities in the case of the simplest Kähler potential $`K=f(x)=x`$. In this appendix we calculate them in an arbitrary case $`K=f(x)`$. We can calculate the geometric quantities of the $`O(N)`$ model by using the formulas obtained in Sec. 4.1. First the metric is (we omit primes except for $`\stackrel{}{v}^{}`$) $`g_{ij^{}}=_i_j^{}K=G_{RS}E_i^R(E_j^S)^{},`$ $`G_{RS^{}}=f^{}(z)(\stackrel{}{v}^{}Z_{S}^{}{}_{}{}^{}\xi ^{}\xi Z_R\stackrel{}{v}^{})+f^{\prime \prime }(z)(\stackrel{}{v}^{}Z_{S}^{}{}_{}{}^{}\xi ^{}\xi \stackrel{}{v}^{})(\stackrel{}{v}^{}\xi ^{}\xi Z_R\stackrel{}{v}^{}),`$ $`z\stackrel{\mathrm{def}}{=}\stackrel{}{v}^{}\xi ^{}\xi \stackrel{}{v}^{}.`$ (B.1) At the point $`\phi =0`$, we define derivatives of the arbitrary function by $`f_1\stackrel{\mathrm{def}}{=}f^{}(\stackrel{}{v}^{}\stackrel{}{v}^{})=f^{}(v_{}^{}{}_{}{}^{2}),f_2\stackrel{\mathrm{def}}{=}f^{\prime \prime }(\stackrel{}{v}^{}\stackrel{}{v}^{})=f^{\prime \prime }(v_{}^{}{}_{}{}^{2}),\mathrm{}.`$ (B.2) The auxiliary metric at the point $`\phi =0`$ is $`G_{RS^{}}\|_{\phi =0}=f_1<S^{}R>+f_2<S^{}><R>,`$ (B.3) where we have used the notations in Eq. (4.6). This becomes for the $`O(N)`$ model $`G_{RS^{}}\|_{\phi =0}=\left(\begin{array}{cc}f_1v^2\delta _{IJ^{}}& 0\\ 0& f_1v_{}^{}{}_{}{}^{2}+f_2c^4\end{array}\right),`$ (B.4) from Eq. (4.7). The vielbein and the $`\widehat{H}`$-connection at the point $`\phi =0`$ are $`E_i^R\|_{\phi =0}=\delta _i^R,W_i^M\|_{\phi =0}=0,`$ (B.5) respectively, and differentiations of the vielbein with respect to coordinates are $`_jE_i^R\|_{\phi =0}=0.`$ (B.6) The curvature tensor (4.12) on the point $`\phi =0`$ of an arbitrary symmetric $`G^𝐂/\widehat{H}`$ is given by $`R_{RS^{}UV^{}}\|_{\phi =0}`$ $`=f_1<S^{}V^{}UR>+f_2(<V^{}U><S^{}R>+<S^{}U><V^{}R>+<S^{}V^{}><UR>`$ $`+<V^{}><S^{}UR>+<U><S^{}V^{}R>+<S^{}V^{}U><R>+<S^{}><V^{}UR>)`$ $`+f_3(<V^{}><U><S^{}R>+<V^{}><S^{}U><R>+<V^{}><S^{}><UR>`$ $`+<V^{}U><S^{}><R>+<U><S^{}V^{}><R>+<U><S^{}><V^{}R>)+f_4<V^{}><U><S^{}><R>`$ $`G^{XY^{}}\|_{\phi =0}`$ $`\times (f_1<X^{}SV>+f_2(<S><X^{}V>+<X^{}S><V>+<X^{}><SV>)+f_3<S><X^{}><V>)^{}`$ $`\times (f_1<Y^{}RU>+f_2(<R><Y^{}U>+<Y^{}R><U>+<Y^{}><RU>)+f_3<R><Y^{}><U>),`$ (B.7) where have defined $`R_{ij^{}kl^{}}=R_{RS^{}UV^{}}\delta _i^R(\delta _j^S)^{}\delta _k^U(\delta _l^V)^{}.`$ (B.8) In the case of the $`O(N)`$ model, it can be calculated from Eq. (4.7), to yield $`R_{IJ^{}KL^{}}=v^4\left[{\displaystyle \frac{f_{1}^{}{}_{}{}^{2}+f_1f_2v_{}^{}{}_{}{}^{2}}{f_1v_{}^{}{}_{}{}^{2}+f_2c^4}}\right]\delta _{IK}\delta _{J^{}L^{}}+f_2v^4(\delta _{IJ^{}}\delta _{KL^{}}+\delta _{IL^{}}\delta _{KJ^{}}),`$ $`R_{N1,J^{},KL^{}}=R_{I,N1,KL^{}}=0,`$ $`R_{N1,N1^{},KL^{}}=v^2\left[{\displaystyle \frac{f_1f_2v_{}^{}{}_{}{}^{2}+(f_1f_3f_{2}^{}{}_{}{}^{2})c^4}{f_1}}\right]\delta _{KL^{}},`$ $`R_{N1,J^{},N1,L^{}}=v^4\left[{\displaystyle \frac{f_{1}^{}{}_{}{}^{2}+f_1f_2v_{}^{}{}_{}{}^{2}(2f_{2}^{}{}_{}{}^{2}+f_1f_3)c^4}{f_1v_{}^{}{}_{}{}^{2}+f_2c^4}}\right]\delta _{J^{}L^{}},`$ $`R_{N1,N1^{},N1,L^{}}=R_{N1,N1^{},K,N1^{}}=0,`$ $`R_{N1,N1^{},N1,N1^{}}`$ $`={\displaystyle \frac{1}{f_1v_{}^{}{}_{}{}^{2}+f_2c^4}}[f_{1}^{}{}_{}{}^{2}v^4+f_1f_2v_{}^{}{}_{}{}^{2}(2v_{}^{}{}_{}{}^{4}+v^4)2f_{2}^{}{}_{}{}^{2}(v_{}^{}{}_{}{}^{8}+v^4v_{}^{}{}_{}{}^{4}2v^8)`$ $`+2f_1f_3c^4(2v_{}^{}{}_{}{}^{4}+v^4)+f_1f_4c^8v_{}^{}{}_{}{}^{2}+(f_2f_4f_{3}^{}{}_{}{}^{2})c^{12}],`$ (B.9) where $`c^4=v_{}^{}{}_{}{}^{4}v^4`$. We can show that these results coincide with direct calculations after solving the constraint $`\stackrel{}{\varphi }^2=a^2`$ as $`\varphi ^N=\sqrt{a^2_{i=1}^{N1}(\varphi ^i)^2}`$. However the coset formalism presented in this paper can be generalized to an arbitrary model straightforwardly. Before the calculation of real components of the curvature tensor, we give comments. 1. At the symmetric point, $`v_{}^{}{}_{}{}^{2}=v^2,c^2=\sqrt{v^{\mathrm{\hspace{0.17em}4}}v^4}=0`$, the curvatures can be written as $`R_{ij^{}kl^{}}={\displaystyle \frac{1}{2}}f_{\pi }^{}{}_{}{}^{2}\delta _{ik}\delta _{j^{}l^{}}+g^2(\delta _{ik}\delta _{j^{}l^{}}+\delta _{ij^{}}\delta _{kl^{}}+\delta _{il^{}}\delta _{kj^{}}),`$ (B.10) $`f_{\pi }^{}{}_{}{}^{2}=2f_1v^2,g^2=f_2v^4.`$ (B.11) We recovered previous results in Ref. and Eq. (1.19) for the $`O(N)`$ model. 2. For the case of the linear Kähler potential, $`f(x)=x:f_1=1,f_2=f_3=\mathrm{}=0,`$ (B.12) they reduce to results in Eq. (4.13). We can calculate real components of the curvature tensor which are directly concerned with low-energy theorems, given by $`(i,j,k,l=1,\mathrm{},N2)`$ $`R_{A^iA^jA^kA^l}=2v^4\left[{\displaystyle \frac{f_{1}^{}{}_{}{}^{2}f_{2}^{}{}_{}{}^{2}c^4}{f_1v_{}^{}{}_{}{}^{2}+f_2c^4}}\right](\delta _{ik}\delta _{jl}\delta _{il}\delta _{jk}),`$ $`R_{B^iA^jB^kA^l}=2v^4\left[{\displaystyle \frac{f_{1}^{}{}_{}{}^{2}+2f_1f_2v_{}^{}{}_{}{}^{2}+f_{2}^{}{}_{}{}^{2}c^4}{f_1v_{}^{}{}_{}{}^{2}+f_2c^4}}\right](\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk})4f_2v^4(\delta _{ij}\delta _{kl}),`$ $`R_{B^iA^jA^kA^l}=R_{A^iA^jB^kA^l}=0,`$ $`R_{A^{N1}A^jA^{N1}A^l}`$ $`={\displaystyle \frac{2v^2}{f_1(f_1v_{}^{}{}_{}{}^{2}+f_2c^4)}}[f_{1}^{}{}_{}{}^{3}v^2+f_{1}^{}{}_{}{}^{2}f_2v_{}^{}{}_{}{}^{2}(v^2v_{}^{}{}_{}{}^{2})2f_1f_{2}^{}{}_{}{}^{2}v^2c^4+f_{2}^{}{}_{}{}^{3}c^8,`$ $`f_{1}^{}{}_{}{}^{2}f_3(v^2+v_{}^{}{}_{}{}^{2})c^4f_1f_2f_3c^8]\delta _{jl},`$ $`R_{A^{N1}A^{N1}A^kA^l}=0,`$ $`R_{B^{N1}A^jB^{N1}A^l}`$ $`={\displaystyle \frac{2v^2}{f_1(f_1v_{}^{}{}_{}{}^{2}+f_2c^4)}}[f_{1}^{}{}_{}{}^{3}v^2f_{1}^{}{}_{}{}^{2}f_2v_{}^{}{}_{}{}^{2}(v^2+v_{}^{}{}_{}{}^{2})+2f_1f_{2}^{}{}_{}{}^{2}v^2c^4+f_{2}^{}{}_{}{}^{3}c^8`$ $`+f_{1}^{}{}_{}{}^{2}f_3(v^2v_{}^{}{}_{}{}^{2})c^4f_1f_2f_3c^8]\delta _{jl},`$ $`R_{B^{N1}A^{N1}B^kA^l}=4v^2\left[{\displaystyle \frac{f_1f_2v_{}^{}{}_{}{}^{2}+(f_1f_3f_{2}^{}{}_{}{}^{2})c^4}{f_1}}\right]\delta _{kl},`$ $`R_{B^{N1}A^{N1}B^{N1}A^{N1}}`$ $`={\displaystyle \frac{4}{f_1v_{}^{}{}_{}{}^{2}+f_2c^4}}[f_{1}^{}{}_{}{}^{2}v^4+f_1f_2v_{}^{}{}_{}{}^{2}(2v_{}^{}{}_{}{}^{4}+v^4)2f_{2}^{}{}_{}{}^{2}(v_{}^{}{}_{}{}^{8}+v^4v_{}^{}{}_{}{}^{4}2v^8)`$ $`+2f_1f_3c^4(2v_{}^{}{}_{}{}^{4}+v^4)+f_1f_4c^8v_{}^{}{}_{}{}^{2}+(f_2f_4f_{3}^{}{}_{}{}^{2})c^{12}].`$ (B.13) At the symmetric point, these reduce to $`(i,j,k,l=1,\mathrm{},N1)`$ $`R_{A^iA^jA^kA^l}=f_{\pi }^{}{}_{}{}^{2}(\delta _{ik}\delta _{jl}\delta _{il}\delta _{jk}),`$ $`R_{B^iA^jB^kA^l}=f_{\pi }^{}{}_{}{}^{2}(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk})4g^2(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}+\delta _{ij}\delta _{kl}),`$ $`R_{B^iA^jA^kA^l}=R_{A^iA^jB^kA^l}=0,`$ (B.14) which again coincide with those obtained in Ref. and Eq. (1.22) for the $`O(N)`$ model.
warning/0006/cond-mat0006385.html	ar5iv	text	# Importance of Correlation Effects on Magnetic Anisotropy in Fe and Ni \[ ## Abstract We calculate magnetic anisotropy energy of Fe and Ni by taking into account the effects of strong electronic correlations, spin-orbit coupling, and non-collinearity of intra-atomic magnetization. The LDA+U method is used and its equivalence to dynamical mean–field theory in the static limit is emphasized. Both experimental magnitude of MAE and direction of magnetization are predicted correctly near $`U=4eV`$ for Ni and $`U=3.5eV`$ for Fe. Correlations modify one–electron spectra which are now in better agreement with experiments. PACS numbers: 71.15.Mb 71.15.Rf 71.27.+a 75.30.Gw 75.40.Mg \] The calculation of the magneto-crystalline anisotropy energy (MAE) of magnetic materials containing transition-metal elements from first principles electronic structure calculations is a long-standing problem. The MAE is defined as the difference of total energies with the orientations of magnetization pointing in different, e.g., (001) and (111), crystalline axis. The difference is not zero because of spin-orbit effect, which couples the magnetization to the lattice, and determines the direction of magnetization, called the easy axis. Being a ground state property, the MAE should be accessible in principle via density functional theory (DFT) . Despite the primary difficulty related to the smallness of MAE ($`1\mu eV/`$atom), great efforts to compute the quantity with advanced total energy methods based on local density approximation (LDA) combined with the development of faster computers, have seen success in predicting its correct orders of magnitudes . However, the correct easy axis of Ni has not been predicted by this method and the fundamental problem of understanding MAE is still open. A great amount of work has been done to understand what is the difficulty in predicting the correct axis for Ni. Parameters within the LDA calculation have been varied to capture physical effects which might not be correctly described. These include (i) scaling spin-orbit coupling in order to enlarge its effect on the MAE , (ii) calculating torque to avoid comparing large numbers of energy , (iii) studying the effects of the second Hunds rule in the orbital polarization theory , (iv) analyzing possible changes in the position of the Fermi level by changing the number of valence electrons , (v) using the state tracking method , and (iv) real space approach . In this paper we take a new view that the correlation effects within the $`d`$ shell are important for the magnetic anisotropy of $`3d`$ transition metals like Ni. These effects are not captured by the LDA but are described by Hubbard–like interactions presented in these systems and need to be treated by first principles methods. Another effect which has not been investigated in the context of magnetic anisotropy calculations is the non-collinear nature of intra-atomic magnetization . It is expected to be important when spin-orbit coupling and correlation effects come into play together. In this article we show that when we include these new ingredients into the calculation we solve the long-standing problem of predicting the correct easy axis of Ni. We believe that the physics of transition metal compounds is intermediate between atomic limit where the localized $`d`$ electrons are treated in the real space and fully itinerant limit when the electrons are described by band theory in k space. A many–body method incorporating these two important limits is the dynamical mean–field theory (DMFT) . The DMFT approach has been extensively used to study model Hamiltonian of correlated electron systems in the weak, strong and intermediate coupling regimes. It has been very successful in describing the physics of realistic systems, like the transition metal oxides and, therefore, is expected to treat properly the materials with $`d`$ or $`f`$ electrons. Electron-electron correlation matrix $`U_{\gamma _1\gamma _2\gamma _3\gamma _4}=m_1m_3\left\|v_C\right\|m_2m_4\delta _{s_1s_2}\delta _{s_3s_4}`$ for $`d`$ orbitals is the quantity which takes strong correlations into account. This matrix can be expressed via Slater integrals $`F^{(i)}`$, $`i=0,2,4,6`$ in the standard manner. The inclusion of this interaction generates self–energy $`\mathrm{\Sigma }_{\gamma _1\gamma _2}(i\omega _n,\stackrel{}{k})`$ on top of the one–electron spectra. Within DMFT it is approximated by momentum independent self–energy $`\mathrm{\Sigma }_{\gamma _1\gamma _2}(i\omega _n)`$. A central quantity of the dynamical mean–field theory is the one–electron on–site Green function $`G_{\gamma _1\gamma _2}(i\omega _n)=`$ $`{\displaystyle \underset{\stackrel{}{k}}{}}`$ $`[(i\omega _n+\mu )O_{\gamma _1\gamma _2}(\stackrel{}{k})H_{\gamma _1\gamma _2}^0(\stackrel{}{k})`$ (1) $`+`$ $`v_{dc}\mathrm{\Sigma }_{\gamma _1\gamma _2}(i\omega _n)]^1.`$ (2) where $`H_{\gamma _1\gamma _2}^0(\stackrel{}{k})`$ is the one-electron Hamiltonian standardly treatable within the LDA. Since the latter already includes the electron-electron interactions in some averaged way, we subtract the double counting term $`v_{dc}`$ . The use of realistic localized orbital representation such as linear muffin–tin orbitals leads us to include overlap matrix $`O_{\gamma _1\gamma _2}(\stackrel{}{k})`$ into the calculation. The DMFT reduces the problem to solving effective impurity model where the correlated $`d`$ orbitals are treated as an impurity level hybridized with the bath of conduction electrons. The role of hybridization is played by the so–called bath Green function defined as follows: $$[𝒢_0^1]_{\gamma _1\gamma _2}(i\omega _n)=G_{\gamma _1\gamma _2}{}_{}{}^{1}(i\omega _n)+\mathrm{\Sigma }_{\gamma _1\gamma _2}(i\omega _n).$$ (3) Solving this impurity model gives access to the self–energy $`\mathrm{\Sigma }_{\gamma _1\gamma _2}(i\omega _n)`$ for the correlated electrons. The one–electron Green function (2) is now modified with new $`\mathrm{\Sigma }_{\gamma _1\gamma _2}(i\omega _n)`$, which generates a new bath Green function. Therefore, the whole problem requires self–consistency. In this paper we confine ourselves to zero temperature and make an additional assumption on solving the impurity model using the Hartree–Fock approximation. In this approximation the self–energy reduces to $$\mathrm{\Sigma }_{\gamma _1\gamma _2}=\underset{\gamma _3\gamma _4}{}(U_{\gamma _1\gamma _2\gamma _3\gamma _4}U_{\gamma _1\gamma _2\gamma _4\gamma _3})\overline{n}_{\gamma _3\gamma _4}$$ (4) where $`\overline{n}_{\gamma _1\gamma _2}`$ is the average occupation matrix for the correlated orbitals. The off-diagonal elements of the occupancy matrix are not zero when spin-orbit coupling is included . The latter can be implemented following the prescription of Andersen or more recent one by Pederson . In the Hartree–Fock limit the self–energy is frequency independent and real. The self–consistency condition of DMFT can be expressed in terms of the average occupation matrix: Having started from some $`\overline{n}_{\gamma _1\gamma _2}`$ we find $`\mathrm{\Sigma }_{\gamma _1\gamma _2}`$ according to (4). Fortunately, the computation of the on–site Green function (2) needs not to be performed. Since the self–energy is real, the new occupancies can be calculated from the eigenvectors of the one–electron Hamiltonians with $`\mathrm{\Sigma }_{\gamma _1\gamma _2}v_{dc}`$ added to its $`dd`$ block. The latter can be viewed as an orbital–dependent potential which has been introduced by the LDA+U method . The LDA$`+`$U method has been very successful compared with experiments at zero temperature in ordered compounds. By establishing its equivalence to the static limit of the DMFT we see clearly that dynamical mean–field theory is a way of improving upon it, which is crucial for finite temperature properties. In this work we study the effect of the Slater parameter $`F_0`$ that is the Hubbard on site interaction U on the magnetic anisotropy energy. The latter is calculated by taking the difference of two total energies with different directions of magnetization (MAE=$`E(111)E(001)`$). The total energies are obtained via fully self consistent solutions. Since the total energy calculation requires high precision, full potential LMTO method has been employed. For the $`\stackrel{}{k}`$ space integration, we follow the analysis given by Trygg and co–workers and use the special point method with a Gaussian broadening of $`15mRy`$. The validity of this procedure has been tested in their work . For convergence of the total energies within desired accuracy, about $`15000k`$-points are needed. We used $`28000k`$-points to reduce possible numerical noise. Our calculations include non-spherical terms of the charge density and potential both within the atomic spheres and in the interstitial region . All low-lying semi-core states are treated together with the valence states in a common Hamiltonian matrix in order to avoid unnecessary uncertainties. These calculations are spin polarized and assume the existence of long-range magnetic order. Spin-orbit coupling is implemented according to the suggestions by Andersen . We also treat magnetization as a general vector field, which realizes non-collinear intra-atomic nature of this quantity. Such general magnetization scheme has been recently discussed in . We now discuss our calculated MAE. We first test our method in case of LDA ($`U=0`$). To compare with previous calculations, we turn off the non-collinearity of magnetization which makes it collinear with the quantization axis. The calculation gives correct orders of magnitude for both fcc Ni and bcc Fe but with the wrong easy axis for Ni, which is the same result as the previous result . Turning on the non-collinearity results in a a larger value of the absolute value of the MAE ($`2.9\mu eV`$) for Ni but the easy axis predicted to be (001) which is still wrong. The magnitude of the experimental MAE of Ni is $`2.8\mu eV`$ aligned along $`(111)`$ direction . We now describe the effect of strong correlations, which is crucial in predicting the correct axis of Ni (see Fig. 1). As U increases, the MAE of Ni smoothly increases until $`U`$ reaches $`2.5eV`$ and then smoothly decreases up to the value $`3.8\mu eV`$. Around $`U=3.9eV`$, the MAE decreases abruptly to negative value. Around $`U=4.0eV`$, the experimental order of magnitude and the correct easy axis (111) are restored. The change from the wrong easy axis to the correct easy axis occurs over the range of $`\delta U0.2eV`$, which is the order of spin-orbit coupling constant ($`0.1eV`$). For Fe, the MAE decreases on increasing $`U`$ to negative values, where the magnetization takes the wrong axis. From $`U=2.7eV`$, it increases back to the correct direction of easy axis (positive MAE). Around $`U=3.5eV`$, it restores the correct easy axis and the experimental value of MAE is reproduced. It is remarkable that the values of $`U`$ necessary to reproduce the correct magnetic anisotropy energy are very close to the values which are needed to describe photoemission spectra of these materials . This shows an internal consistency of our approach and emphasizes the importance of correlations. We find direct correlation between the dependency of the MAE as a function of $`U`$ and the difference of magnetic moments ($`\mathrm{\Delta }M=(M(111)M(001)`$) behaving similarly. For Ni, the difference of magnetic moments is nearly $`U`$ independent up to $`U=3eV`$. For large $`U`$, it smoothly decreases from the positive value to the negative one. It also decreases rapidly around $`U=3.9eV`$ in accord with MAE. For Fe, the difference is positive at $`U=0`$. It decreases slightly to the negative values and then increases to the positive value over the range of $`U<2.7eV`$ where MAE decreases. For larger $`U`$’s, where MAE is coming back to positive value, its slope is significantly larger than that at smaller $`U`$’s. This concurrent change of MAE and the difference of magnetic moments suggests why some previous attempts based on force theorem failed in predicting the correct easy axes. Force theorem replaces the difference of the total energies by the difference of one–electron energies. In this approach, the contribution from the slight difference in magnetic moments does not appear and, therefore, is not counted in properly. Unfortunately, we could not find any experimental data of magnetic moments to the desired precision ($`10^4\mu _B`$) to compare with. We also have problems in reaching the convergence of the total energy with desired accuracy for large values of $`U`$ in both Fe and Ni. We now present implications of our results on the calculated electronic structure for the case of Ni. One important feature which emerges from the calculation is the absence of the $`X_2`$ pocket (see Fig. 2). This has been predicted by LDA but has not been found experimentally . The band corresponding to the pocket is pushed down well below the Fermi level. This is expected since correlation effects are more important for slower electrons and the velocity near the pocket is rather small. It turns out that the whole band is submerged under the Fermi level. There has been some suspicions that the incorrect position of the $`X_2`$ band within LDA was responsible for the incorrect prediction of the easy axis within this theory. Daalderop and coworkers removed the $`X_2`$ pocket by increasing the number of valence electrons and found the correct easy axis. We therefore conclude that the absence of the pocket is one of the central elements in determining the magnetic anisotropy, and there is no need for any ad-hoc adjustment within a theory which takes into account the correlations. We now describe the effects originated from (near) degenerate states close to the Fermi surface. These have been of primary interest in past analytic studies . We will call such states degenerate Fermi surface crossing (DFSC) states. The contribution to MAE by non-DFSC states comes from the fourth order perturbation. Hence it is of the order of $`\lambda ^4`$. The energy splitting between DFSC states due to spin-orbit coupling is of the order of $`\lambda `$ because the contribution comes from the first order perturbation. Using linear approximation of the dispersion relation $`ϵ(\stackrel{}{k}\lambda )`$, the relevant volume in $`k`$-space was found of the order $`\lambda ^3`$. Thus, these DFSC states make contribution of the order of $`\lambda ^4`$. However, there may be accidentally DFSC states appearing along a line on the Fermi surface, rather than at a point. We have found this case in our LDA calculation for Ni. Therefore the contribution of DFSC states is as important as the bulk non-DFSC states though the degeneracies occur only in small portion of the Brillouin zone. The importance of the DFSC states leads us to comparative analysis of the LDA and LDA+U band structures near the Fermi level. In LDA, five bands are crossing the Fermi level at nearly the same points along the $`\mathrm{\Gamma }X`$ direction. Two of the five bands are degenerate for the residual symmetry and the other three bands accidentally cross the Fermi surface at nearly the same points. There are two $`sp`$ bands with spin up and spin down, respectively. The other three bands are dominated by $`d`$ orbitals. In LDA$`+`$U, one of the $`d`$ bands is pushed down below the Fermi surface. The other four bands are divided into two degenerate pieces at the Fermi level (see Fig. 2): Two symmetry related degenerate $`d_{}`$ bands and two near degenerate $`sp_{}`$ and $`sp_{}`$ bands. In LDA, we found that two bands are accidentally near degenerate along the line on the Fermi surface within the plane $`\mathrm{\Gamma }XL`$. One band is dominated by $`d_{}`$ orbitals. The other is dominated by $`s_{}`$ orbitals near $`X`$ and by $`d_{}`$ orbitals off $`X`$. In LDA+U, these accidental DFSC states disappear(see Fig. 2). Instead, there is new two-fold DFSC states along $`\mathrm{\Gamma }L`$ direction, both of which are dominated by $`d_{}`$ orbitals. As we have seen, the on-site repulsion $`U`$ reduces the number of DFSC states along $`\mathrm{\Gamma }X`$ direction while increasing that of DFSC states along $`\mathrm{\Gamma }L`$ direction. Based on the tight–binding model, Mori and coworkers have shown that DFSC states along $`\mathrm{\Gamma }L`$ direction result in the magnetization aligned along $`(1,1,1)`$ direction and DFSC states along $`\mathrm{\Gamma }X`$ direction result in the magnetization aligned along $`(0,0,1)`$ direction. Since the strong correlation does precisely this, we conclude that disappearance of DFSC states along $`\mathrm{\Gamma }X`$ direction and their appearance along $`\mathrm{\Gamma }L`$ direction is another important element that determines the easy axis of Ni. Unlike LDA, we have found two extra very tiny $`L`$ pockets in LDA$`+`$U (see Fig. 2). Both of them are dominated by $`sp`$ orbital with opposite spins. Being small, these extra pockets may be artifacts of LDA$`+`$U. To conclude, we have demonstrated that it is possible to perform highly precise calculation of the total energy in order to obtain both the correct easy axes and the magnitudes of MAE for Fe and Ni. This has been accomplished by including the strong correlation effect via the Hubbard-like on-site repulsion $`U`$ and incorporating the non–collinear magnetization. In both Fe and Ni, the on-site $`U`$ takes physically acceptable values consistent with the values known from atomic physics. The calculations performed are state of the art in what can currently be achieved for realistic treatments of correlated solids. Further studies should be devoted to improving the quality of the solution of the impurity model within DMFT and extending the calculation to finite temperatures. This research was supported by the ONR grant No. 4-2650. GK would like to thank K. Hathaway for discussing the origin of magnetic anisotropy and G. Lonzarich for discussing dHvA data. We thank R. Chitra for stimulating discussion. We thank V. Oudovenko for setting up the computer cluster used to perform these calculations. We have also used the supercomputer at the Center for Advanced Information Processing, Rutgers. IY thanks K. H. Ahn for discussions.
warning/0006/nucl-th0006083.html	ar5iv	text	# MKPH-T-00-10 The GDH sum rule for the deuteron ## 1 Introduction The Gerasimov-Drell-Hearn (GDH) sum rule connects the anomalous magnetic moment of a particle with the energy weighted integral - henceforth denoted by $`I^{GDH}`$ \- from threshold up to infinity over the difference of the total cross sections for the absorption of circularly polarized photons on a target with spin parallel ($`\sigma ^P(k)`$) and antiparallel ($`\sigma ^A(k)`$) to the spin of the photon. In detail it reads for a particle of mass $`M_t`$, charge $`eQ`$, anomalous magnetic moment $`\kappa `$ and spin $`I`$ $$I^{GDH}=4\pi ^2\kappa ^2\frac{e^2}{M_t^2}I=_0^{\mathrm{}}\frac{dk}{k}\left(\sigma ^P(k)\sigma ^A(k)\right),$$ (1) where the anomalous magnetic moment is defined by the total magnetic moment operator of the particle $`\stackrel{}{M}=(Q+\kappa )\frac{e}{M_t}\stackrel{}{S}`$, where $`\stackrel{}{S}`$ denotes the spin operator of the target. This sum rule gives a very interesting relation between a magnetic ground state property of a particle and an integral property of its whole excitation spectrum. It shows that the existence of a nonvanishing anomalous magnetic moment is directly tied to an internal dynamic structure of the particle. Furthermore, it tells us, because the lhs of (1) is positive, that the integrated, energy-weighted total absorption of a circularly polarized photon on a particle with its spin parallel to the photon spin is larger than the one on a target with its spin antiparallel, if the anomalous magnetic moment does not vanish. The GDH sum rule has first been derived by Gerasimov and, shortly afterwards, independently by Drell and Hearn and also, less well known, by Hosada and Yamamoto. The last authors have used current algebra relations while the others based the derivation on the low energy theorem for the Compton scattering amplitude of a particle and the assumption of an unsubtracted dispersion relation for the difference of the elastic forward scattering amplitudes for circularly polarized photons and a completely polarized target with spin parallel and antiparallel to the photon spin. First, I will briefly discuss a general class of photoabsorption sum rules of which the GDH sume rule is a special case. Then I will address the question whether the GDH sume rule of the neutron can be determined from the GDH sum rule of the deuteron in the absence of free neutron targets. Subsequently, I will present results on an evaluation of the GDH sum rule for the deuteron by explicit integration of the GDH integral up to a photon energy of 550 MeV including the photodisintegration channel as well as coherent and incoherent single pion photoproduction channels. I will close with some conclusions and an outlook. ## 2 A general class of photoabsorption sum rules The GDH sum rule belongs to a larger class of photoabsorption sum rules related to the various contributions to the total photoabsorption cross section for the general case of beam and target polarization $`\sigma _{\mathrm{tot}}(k,\rho ^\gamma ,\rho ^t)={\displaystyle \frac{1}{2}}{\displaystyle \underset{J=0}{\overset{2I}{}}}P_J^t`$ $`[`$ $`(1+()^J)\sigma _J^{11}(k)`$ (2) $`+`$ $`(1()^J)P_c^\gamma \sigma _J^{11}(k)P_J(\mathrm{cos}\theta _t)`$ $`+`$ $`(1+()^J)P_l^\gamma \sigma _J^{11}(k)d_{20}^J(\theta _t)\mathrm{cos}(2\varphi _t)],`$ where $`P_l^\gamma `$ and $`P_c^\gamma `$ denote the degree of linear and circular photon polarization, respectively. Furthermore, the target polarization parameters $`P_J^t`$ with respect to an orientation direction, characterized by the angles $`\theta _t`$ and $`\varphi _t`$, are defined by the target polarization density matrix $`\rho _{MM^{}}^t={\displaystyle \frac{()^{IM^{}}}{\widehat{I}}}{\displaystyle \underset{J,m}{}}`$ $`\widehat{J}`$ $`\left(\begin{array}{ccc}I& I& J\\ M^{}& M& m\end{array}\right)P_J^te^{im\varphi _t}d_{m0}^J(\theta _t).`$ (3) The separate contributions $`\sigma _J^{\lambda ^{}\lambda }`$ are related to the forward Compton scattering amplitude via the optical theorem $`\sigma _J^{\lambda ^{}\lambda }(k)={\displaystyle \frac{4\pi }{k}}\mathrm{}mT_{\lambda ^{}\lambda }^J(k),`$ (4) where $`T_{\lambda ^{}\lambda }^J`$ is defined by $`T_{\lambda ^{}M^{},\lambda M}(k)=()^{IM}\widehat{I}{\displaystyle \underset{J=0}{\overset{2I}{}}}\widehat{J}\left(\begin{array}{ccc}I& J& I\\ M^{}& \lambda \lambda ^{}& M\end{array}\right)T_{\lambda ^{}\lambda }^J(k).`$ (5) It can be expressed in terms of generalized polarizabilities $`T_{\lambda ^{}\lambda }^J(k)={\displaystyle \frac{\widehat{J}}{\widehat{I}}}{\displaystyle \underset{L^{}L}{}}()^{L^{}+L}\left(\begin{array}{ccc}L& L^{}& J\\ \lambda & \lambda ^{}& \lambda ^{}\lambda \end{array}\right)P_J^{L^{}L\lambda ^{}\lambda }(k)`$ (6) with $`P_J^{L^{}L\lambda ^{}\lambda }(k)={\displaystyle \underset{\nu ^{},\nu =0,1}{}}\lambda ^\nu ^{}\lambda ^\nu P_J(M^\nu ^{}L^{},M^\nu L;k),`$ (7) where $`M^0=E`$ (electric) and $`M^1=M`$ (magnetic) multipole. The $`T_{\lambda ^{}\lambda }^J`$ are also related to the expansion of the scattering amplitude in terms of a complete set of operators $`\tau ^{[J]}`$ with $`J=0,1,\mathrm{},2I`$ in the ground state spin space with reduced matrix elements $`I\tau ^{[J]}I=\widehat{I}\widehat{J}`$ $`T_{\lambda ^{}M^{},\lambda M}(k)={\displaystyle \underset{J=0}{\overset{2I}{}}}()^{\lambda ^{}+\lambda }IM^{}\|\tau _{\lambda \lambda ^{}}^{[J]}\|IM\mathrm{\Omega }_{\lambda ^{}\lambda }^{[J]}(k).`$ (8) Comparison with (5) leads to the simple relation $`T_{\lambda ^{}\lambda }^J(k)=\mathrm{\Omega }_{\lambda ^{}\lambda }^{[J]}(k).`$ (9) Specifically one has $`\sigma _J^{11}`$ $`=`$ $`{\displaystyle \frac{\widehat{J}}{\widehat{I}}}{\displaystyle \underset{M}{}}()^{IM}\left(\begin{array}{ccc}I& J& I\\ M& 0& M\end{array}\right)\sigma _{1M},`$ (10) where $`\sigma _{1M}`$ denotes the total cross section for the absorption of a photon with helicity $`\lambda =1`$ by a target with definite spin projection M on the photon momentum. Corresponding expressions hold for $`\sigma _J^{11}`$ with respect to the absorption of linearly polarized photons. In detail one finds for $`J=0,1,2`$ $`\sigma _0^{11}`$ $`=`$ $`{\displaystyle \frac{1}{\widehat{I}^2}}{\displaystyle \underset{M}{}}\sigma _{1M},`$ (11) $`\sigma _1^{11}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{\widehat{I}^2\sqrt{I(I+1)}}}{\displaystyle \underset{M}{}}M\sigma _{1M},`$ (12) $`\sigma _2^{11}`$ $`=`$ $`{\displaystyle \frac{\sqrt{5}}{\widehat{I}^2\sqrt{I(I+1)}}}{\displaystyle \underset{M}{}}{\displaystyle \frac{(3M^2I(I+1))}{\sqrt{(2I1)(2I+3)}}}\sigma _{1M},`$ (13) and for the spin asymmetry ($`\sigma ^{P/A}=\sigma _{1,\pm I}`$) $`\sigma ^P\sigma ^A`$ $`=`$ $`\widehat{I}{\displaystyle \underset{J}{}}\widehat{J}(1()^J)\left(\begin{array}{ccc}I& I& J\\ I& I& 0\end{array}\right)\sigma _J^{11}`$ (14) $`=`$ $`{\displaystyle \frac{2\sqrt{3I}}{\sqrt{I+1}}}\sigma _1^{11}+\mathrm{}`$ Now, crossing symmetry implies $`(T_{\lambda ^{}\lambda }^J(k))^{}=()^JT_{\lambda ^{}\lambda }^J(k).`$ (15) For $`J=\text{even}`$ one takes a once-subtracted dispersion relation for $`T_{\lambda ^{}\lambda }^J(k)`$ $`\mathrm{}e\left(T_{\lambda ^{}\lambda }^J(k)T_{\lambda ^{}\lambda }^J(0)\right)`$ $`=`$ $`{\displaystyle \frac{2k^2}{\pi }}𝒫{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk^{}}{k^{}}}{\displaystyle \frac{\mathrm{}mT_{\lambda ^{}\lambda }^J(k^{})}{k^2k^2}}`$ (16) $`=`$ $`{\displaystyle \frac{k^2}{2\pi ^2}}𝒫{\displaystyle _0^{\mathrm{}}}𝑑k^{}{\displaystyle \frac{\sigma _J^{\lambda ^{}\lambda }(k^{})}{k^2k^2}},`$ while for $`J=\text{odd}`$ an unsubtracted dispersion relation applies $`\mathrm{}eT_{\lambda ^{}\lambda }^J(k)`$ $`=`$ $`{\displaystyle \frac{2k}{\pi }}𝒫{\displaystyle _0^{\mathrm{}}}𝑑k^{}{\displaystyle \frac{\mathrm{}mT_{\lambda ^{}\lambda }^J(k^{})}{k^2k^2}}`$ (17) $`=`$ $`{\displaystyle \frac{k}{2\pi ^2}}𝒫{\displaystyle _0^{\mathrm{}}}𝑑k^{}k^{}{\displaystyle \frac{\sigma _J^{\lambda ^{}\lambda }(k^{})}{k^2k^2}}.`$ A power series expansion according to (15) $`\mathrm{}eT_{\lambda ^{}\lambda }^J(k)=\{\begin{array}{cc}_{\nu =0}^{\mathrm{}}t_\nu ^{\lambda ^{}\lambda ,J}k^\nu \hfill & \text{ for }J\text{ even,}\hfill \\ _{\nu =0}^{\mathrm{}}t_\nu ^{\lambda ^{}\lambda ,J}k^{\nu +1}\hfill & \text{ for }J\text{ odd,}\hfill \end{array}`$ (20) yields a class of sum rules $`t_\nu ^{\lambda ^{}\lambda ,J}=\{\begin{array}{cc}\frac{1}{2\pi ^2}_0^{\mathrm{}}𝑑k^{}\frac{\sigma _J^{\lambda ^{}\lambda }(k^{})}{k^{2\nu }}\hfill & \text{ for }J\text{ even and }\nu =1,2,\mathrm{},\hfill \\ & \\ \frac{1}{2\pi ^2}_0^{\mathrm{}}𝑑k^{}\frac{\sigma _J^{\lambda ^{}\lambda }(k^{})}{k^{2\nu +1}}\hfill & \text{ for }J\text{ odd and }\nu =0,1,\mathrm{},\hfill \end{array}`$ (24) one of which is the GDH, namely for $`J=`$ odd and $`\nu =0`$. Because from the low-energy expansion of the Compton amplitude $`T_{\lambda M,\lambda M}(k)=e^2{\displaystyle \frac{Q^2}{M_t}}+\lambda \kappa ^2{\displaystyle \frac{e^2}{M_t^2}}S_z_{IM}k+𝒪(k^2),`$ (25) one finds specifically $`T_{\lambda ^{}\lambda }^J(k)`$ $`=`$ $`\{\begin{array}{cc}\delta _{\lambda ^{}\lambda }\delta _{J0}e^2\frac{Q^2}{M_t}+𝒪(k^2)\hfill & \text{ for }J\text{ even,}\hfill \\ k\left(\delta _{\lambda ^{}\lambda }\delta _{J1}\lambda \kappa ^2\frac{e^2}{M_t^2}\frac{\sqrt{I(I+1)}}{\sqrt{3}}+𝒪(k^2)\right)\hfill & \text{ for }J\text{ odd.}\hfill \end{array}`$ (28) The latter yields the GDH sum rule in the form $`4\pi ^2{\displaystyle \frac{\kappa ^2e^2}{M_t^2}}I`$ $`=`$ $`{\displaystyle \frac{2\sqrt{3I}}{\sqrt{I+1}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk^{}}{k^{}}}\sigma _1^{11}(k^{}),`$ (29) from which (1) follows because of (14) and the fact that the higher order terms $`\sigma _J^{11}`$ for $`J>1`$ do not contribute to the left hand side. ## 3 Is it possible to get the GDH sum rule for the neutron from the one of the deuteron? With respect to the GDH sum rule for the neutron, it has been suggested to measure in the absence of neutron targets its spin asymmetry using a polarized deuteron target. It would rest on two assumptions: (i) A vector polarized deuteron constitutes effectively a polarized neutron target. (ii) The contribution of the meson production to the spin asymmetry of the deuteron is dominated by the quasifree process, and one can neglect binding and final state interaction effects arising from the presence of the spectator nucleon, so that the deuteron spin asymmetry is an incoherent sum of proton and neutron contributions. However, I would like to point out a few “caveats” which make it very unlikely that one can determine the spin asymmetry of the neutron in this way: (i) First of all, the neutron is not completely polarized in a completely vector polarized deuteron target, i.e., $`P(n)=11.5p_D`$, and its polarization degree is slightly model dependend due to the appearance of the deuteron $`D`$-wave probability $`p_D`$ which is not observable. (ii) A model calculation of the impulse approximation (IA) for the incoherent pion photoproduction on the deuteron by R. Schmidt et al. shows already for the unpolarized cross section that the complete IA is not the incoherent sum of proton and neutron contributions. In addition, final state interaction and other two-body effects will very likely add further complications thus spoiling this simple idea. (iii) Since polarization observables show in general a stronger sensitivity to small dynamical and coherence effects, the spin asymmetry might be even more sensitive to the mentioned disturbing effects. (iv) The contribution from coherent $`\pi ^0`$ production on the deuteron is non-negligible which certainly is not an incoherent sum of $`\pi ^0`$ production on proton and neutron. Thus all these factors will prohibit a simple determination of the neutron spin asymmetry by subtracting from the meson production part of the deuteron’s spin asymmetry the proton one. That does not mean, that one does not learn anything about neutron properties. On the contrary, pion photoproduction on the deuteron will provide a very important test for the understanding of the production process on the neutron. But this can be achieved only in the context of a reliable theoretical model which takes into acount all important two-body effects. ## 4 The GDH sum rule for the deuteron In the case of the deuteron, one finds a very interesting cancellation of large contributions. The deuteron has isospin zero, excluding most of the contribution of the large nucleon anomalous magnetic moments to its magnetic moment, and thus one finds a very small anomalous magnetic moment, namely $`\kappa _d=.143`$ resulting in a GDH prediction of $`I_d^{GDH}=0.65\mu `$b, which is more than two orders of magnitude smaller than the nucleon values. Considering the possible absorption processes, one notes first that the incoherent pion production on the deuteron is dominated by the quasifree production on the nucleons bound in the deuteron. This gives a rough estimate for its contribution to the GDH value, namely the sum of the proton and neutron GDH values of 438 $`\mu `$b. Another contribution arises from the coherent $`\pi ^0`$ production channel. On the other hand, for the additional photodisintegration channel which is the only photoabsorption process below the pion production threshold, one finds at very low energies near threshold a sizeable negative contribution which arises from the $`M1`$-transition to the virtual $`{}_{}{}^{1}S_{0}^{}`$ state, because this state can only be reached if the spins of photon and deuteron are antiparallel, and is forbidden for the parallel situation as has been pointed out, for example in Ref.. We have evaluated explicitly the finite GDH sum rule for the deuteron by integrating up to a photon energy of 550 MeV. Three contributions have been included: (i) the photodisintegration channel $`\gamma dnp`$, (ii) the coherent pion production $`\gamma d\pi ^0d`$, and (iii) the incoherent pion production $`\gamma d\pi NN`$. The upper integration limit of 550 MeV has been chosen because on the one hand one finds sufficient convergence for the photodisintegration channel, while on the other hand only single pion photoproduction has been considered, thus limiting the applicability of the present theoretical treatment to energies not too far above the two pion production threshold as long as significant contributions from multipion production cannot be expected. I will now discuss the three contributions separately and refer for details to Ref.. The photodisintegration channel is evaluated within the nonrelativistic framework as is described in detail in Ref. but with inclusion of the most important relativistic contributions. Explicitly, all electric and magnetic multipoles up to the order $`L=4`$ are considered which means inclusion of the final state interaction in all partial waves up to $`j=5`$. For the calculation of the initial deuteron and the final n-p scattering wave functions we use the realistic Bonn potential (r-space version). In the current operator we distinguish the one-body currents with Siegert operators (N), explicit meson exchange contributions (MEC) beyond the Siegert operators, essentially from $`\pi `$\- and $`\rho `$-exchange, contributions from isobar configurations of the wave functions (IC), calculated in the impulse approximation, and leading order relativistic contributions (RC) of which the spin-orbit current is by far the most dominant part. The results are summarized in Fig. 1, where the spin asymmetry and the GDH integral is shown. The GDH values are listed in Tab. 1. One readily notes the huge negative contribution from the $`{}_{}{}^{1}S_{0}^{}`$-state at low energies (see the upper left panel of Fig. 1). Here, the effects from MEC are relatively strong, resulting in an enhancement of the negative value by about 15 percent. Isobar effects are significant in the region of the $`\mathrm{\Delta }`$-resonance, as expected. They give a positive contribution, but considerably smaller in absolute size than MEC. The largest positive contribution stems from RC in the energy region up to about 100 MeV (see the upper right panel of Fig. 1) reducing the GDH value in absolute size by more than 30 percent. This strong influence from relativistic effects is not surprising in view of the fact, that the correct form of the term linear in the photon momentum of the low energy expansion of (25) is only obtained if leading order relativistic contributions are included. The total sum rule value from the photodisintegration channel then is $`I_{\gamma dnp}^{GDH}(550\text{MeV})=413\mu `$b. Its absolute value almost equals within less than ten percent the sum of the free proton and neutron values. This may not be accidental since the large value is directly linked to the nucleon anomalous magnetic moment as is demonstrated by the fact that one finds indeed a very small but positive value $`I_{\gamma dnp}^{GDH}(550\text{MeV})=7.3\mu `$b if the nucleon anomalous magnetic moment is switched off in the e.m. one-body current operator. The theoretical model used to calculate the contribution of the coherent pion production channel is described in detail in Ref.. The reaction is clearly dominated by the magnetic dipole excitation of the $`\mathrm{\Delta }`$ resonance from which one obtains a strong positive $`I_{\gamma dd\pi ^0}^{GDH}`$ contribution because the $`\mathrm{\Delta }`$-excitation is favoured if photon and nucleon spins are parallel. The model takes into account pion rescattering by solving a system of coupled equations for the N$`\mathrm{\Delta }`$, NN$`\pi `$ and NN channels. The inclusion of the rescattering effects is important and leads in general to a significant reduction of the unpolarized cross section in reasonable agreement with the differential cross section data available in the $`\mathrm{\Delta }`$ region. Fig. 2 shows the result of our calculation. One sees the strong positive contribution from the $`\mathrm{\Delta }`$-excitation giving a value $`I_{\gamma dd\pi ^0}^{GDH}(550\text{MeV})=63\mu `$b. The comparison with the unpolarized cross section, also plotted in Fig. 2, demonstrates the dominance of $`\sigma ^P`$ over $`\sigma ^A`$. Furthermore, one notes quite satisfactory convergence. The calculation of the $`\gamma d\pi NN`$ contributions to the GDH integral is based on the spectator nucleon approach discussed in Ref.. In this framework, the reaction proceeds via pion production on one nucleon while the other nucleon acts merely as a spectator. Thus, the $`\gamma d\pi NN`$ operator is given as the sum of the elementary $`\gamma N\pi N`$ operators of the two nucleons. For this elementary operator, we have taken the standard pseudovector Born terms and the contribution of the $`\mathrm{\Delta }`$ resonance, and a satisfactory description of pion photoproduction on the nucleon is achieved in the $`\mathrm{\Delta }`$-resonance region. Although the spectator model does not include any final state interaction except for the resonant $`M_{1+}^{3/2}`$ multipole, it gives quite a good description of available data on the total cross section demonstrating the dominance of the quasifree production process, for which the spectator model should work quite well. The results are collected in Fig. 3. The upper part shows the individual contributions from the different charge states of the pion and their total sum to the cross section difference for pion photoproduction on both the deuteron and for comparison on the nucleon. One notes qualitatively a similar behaviour although the maxima and minima are smaller and also slightly shifted towards higher energies for the deuteron. In the lower part of Fig. 3 the corresponding GDH integrals are shown. A large positive contribution comes from $`\pi ^0`$-production whereas the charged pions give a negative but - in absolute size - smaller contribution to the GDH value. Up to an energy of 550 MeV one finds for the total contribution of the incoherent pion production channels a value $`I_{\gamma dNN\pi }^{GDH}(550\text{MeV})=167\mu `$b which is remarkably close to the sum of the neutron and proton values for the given elementary model $`I_n^{GDH}(550\text{MeV})+I_p^{GDH}(550\text{MeV})=163\mu `$b. However, as is evident from Fig. 3, convergence is certainly not reached at this energy. Thus it is not clear whether this result is accidental, and in addition, one has to see what the effect of neglected rescattering and other two-body contributions will be. Furthermore, the elementary pion production operator used in this work had been constructed primarily to give a realistic description of the $`\mathrm{\Delta }`$ resonance region. In fact, it underestimates the GDH inegral up to 550 MeV by about a factor two compared to a corresponding evaluation based on a multipole analysis of experimental pion photoproduction data. The contributions from all three channels and their sum are listed in Tab. 2. A very interesting and important result is the large negative contribution from the photodisintegration channel near and not too far above the break-up threshold with surprisingly large relativistic effects below 100 MeV. Hopefully, this low energy feature of the spin asymmetry can be checked experimentally in the near future. For the total GDH value from explicit integration up to 550 MeV, we find a negative value $`I_d^{GDH}(550\text{MeV})=183\mu `$b. However, as I have mentioned above, some uncertainty lies in the contribution of the incoherent pion production channel because of shortcomings of the model of the elementary production amplitude above the $`\mathrm{\Delta }`$ resonance. If one uses instead of the model value $`I_{\gamma dNN\pi }^{GDH}(550\text{MeV})=167\mu `$b the sum of the GDH values of neutron and proton by integrating the cross section difference obtained from a multipole analysis of experimental data (fit SM95 from the VPI-analysis), giving $`I_n^{GDH}(550\text{MeV})+I_p^{GDH}(550\text{MeV})=331\mu `$b, one finds for the deuteron $`I_d^{GDH}(550\text{MeV})=19\mu `$b, which I consider a more realistic estimate. Since this value is still negative, a positive contribution of about the same size should come from contributions at higher energies in order to fulfil the small GDH sum rule for the deuteron, provided that the sum rule is valid. These contributions should come from the incoherent single pion production above 550 MeV, because for this channel convergence had not been reached in contrast to the other two channels, and in addition, from multipion production. ## 5 Conclusions and Outlook In order to summarize let me draw a few important conclusions: The spin asymmetry of the deuteron is a very interesting observable of its own value because of a strong anticorrelation of photodisintegration and pion production. It is also very sensitive to relativistic effects at quite low energies which have never been tested in detail in this observable. It is very doubtful, if not impossible, that one can extract in a simple manner the neutron spin asymmetry from the spin asymmetry of the deuteron. However, the spin asymmetry of the deuteron will provide more detailed tests of our understanding of pion photoproduction in a nucleus, and in particular on the neutron. Future theoretical and experimental studies should be devoted to the following topics: (a) Theory: Improvement of the elementary pion photoproduction amplitude above the two pion threshold and inclusion of multiple pion production. For photodisintegration the inclusion of a retarded $`\mathrm{\Delta }N`$ interaction, of higher nucleon resonances, and of relativistic contributions is needed. Furthermore, other sum rules, alluded to in Sect. 2, should be studied. (b) Experiment: A careful measurement of the spin asymmetry for the proton over a larger energy range. With respect to the deuteron, a measurement of the spin asymmetry is needed, separately for the photodisintegration channel, in particular close to the break-up threshold, as well as for the meson production channel.
warning/0006/hep-ph0006084.html	ar5iv	text	# 1 Introduction ## 1 Introduction In order to seek for a clue to the unified understanding of quarks and leptons, many attempts to give a unified description of the quark and lepton mass matrices have been proposed. The universal seesaw mass matrix model is one of the promising attempts to view the unified description, where the mass matrices $`M_f`$ for the conventional quarks and leptons $`f_i`$ ($`f=u,d,\nu ,e`$; $`i=1,2,3`$) are given by $$(\overline{f}_L\overline{F}_L)\left(\begin{array}{cc}0& m_L\\ m_R& M_F\end{array}\right)\left(\begin{array}{c}f_R\\ F_R\end{array}\right),$$ $`(1.1)`$ and $`m_L`$ and $`m_R`$ are universal for all fermion sectors $`f`$. For $`O(M_F)O(m_R)O(m_L)`$, the mass matrix (1.1) leads to the well-known seesaw expression $$M_fm_LM_F^1m_R.$$ $`(1.2)`$ As a specific version of such universal seesaw model, Fusaoka and one of the authors (Y.K.) have proposed a so-called “democratic” seesaw model : The heavy fermion matrices $`M_F`$ have a simple structure \[(unit matrix)+(democratic matrix)\], i.e., $$M_F=m_0\lambda _f(\mathrm{𝟏}+3b_fX),$$ $`(1.3)`$ $$\mathrm{𝟏}=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),X=\frac{1}{3}\left(\begin{array}{ccc}1& 1& 1\\ 1& 1& 1\\ 1& 1& 1\end{array}\right),$$ $`(1.4)`$ on the basis on which the matrices $`m_L`$ and $`m_R`$ are diagonal: $$m_L=\frac{1}{\kappa }m_R=m_0Z=m_0\left(\begin{array}{ccc}z_1& 0& 0\\ 0& z_2& 0\\ 0& 0& z_3\end{array}\right),$$ $`(1.5)`$ where the parameters $`z_1`$, $`z_2`$ and $`z_3`$ are normalized as $`z_1^2+z_2^2+z_3^2=1`$, and $`m_0`$ is of the order of the electroweak symmetry breaking scale, i.e., $`m_010^2`$ GeV. Since the parameter $`b_f`$ in the charged lepton sector is taken as $`b_e=0`$, the parameters $`z_i`$ are fixed as $$\frac{z_1}{\sqrt{m_e}}=\frac{z_2}{\sqrt{m_\mu }}=\frac{z_3}{\sqrt{m_\tau }}=\frac{1}{\sqrt{m_\tau +m_\mu +m_e}}.$$ $`(1.6)`$ For the up-type quark sector, the parameter $`b_f`$ is taken as $`b_u=1/3`$, which leads to det$`M_U=0`$, and the seesaw mechanism does not work for one of the three families, and hence we obtain the mass $`m_tm_0/\sqrt{3}`$ without the seesaw suppression factor $`\kappa /\lambda _u`$ (we identify it as the top quark mass). Furthermore, we also obtain a relation $`m_u/m_c3m_e/m_\mu `$, which is in good agreement with the observed values. Moreover, when we take $`b_d1`$ ($`b_d=e^{i\beta _d}`$ with $`\beta _d=18^{}`$) for the down-type quark sector, we can obtain reasonable quark mass ratios and the Cabbibo-Kobayashi-Maskawa (CKM) matrix. The neutrino mass matrix in the universal seesaw mass matrix model is given as follows: $$\left(\begin{array}{cccc}\overline{\nu }_L& \overline{\nu }_R^c& \overline{N}_L& \overline{N}_R^c\end{array}\right)\left(\begin{array}{cccc}0& 0& 0& m_L\\ 0& 0& m_R^T& 0\\ 0& m_R& M_{NL}& M_D\\ m_L^T& 0& M_D^T& M_{NR}\end{array}\right)\left(\begin{array}{c}\nu _L^c\\ \nu _R\\ N_L^c\\ N_R\end{array}\right),$$ $`(1.7)`$ where $`\psi _L^c(\psi _L)^c=C\psi _L^T`$. \[We consider a SO(10)$`{}_{L}{}^{}\times `$SO(10)<sub>R</sub> model , where fermions $`(f_L+F_R^c)`$ and $`(f_R+F_L^c)`$ are assigned to (16,1) and (1,16) under SO(10)$`{}_{L}{}^{}\times `$SO(10)<sub>R</sub>, respectively. Hereafter, we will denote the Majorana mass matrices $`M_{NL}`$ and $`M_{NR}`$ of the neutral heavy leptons $`N_L`$ and $`N_R`$ as $`M_R=M_{NL}`$ and $`M_L=M_{NR}`$, respectively.\] For $`O(m_L)O(m_R)O(M_D),O(M_L),O(M_R)`$, we obtain the following $`6\times 6`$ seesaw mass matrix for $`(\nu _L^c,\nu _R)`$ $$M^{(6\times 6)}\left(\begin{array}{cc}0& m_L\\ m_R^T& 0\end{array}\right)\left(\begin{array}{cc}M_R& M_D\\ M_D^T& M_L\end{array}\right)^1\left(\begin{array}{cc}0& m_R\\ m_L^T& 0\end{array}\right),$$ $`(1.8)`$ which leads to the 3$`\times `$3 seesaw matrices for $`\nu _L`$ and $`\nu _R`$ $$M(\nu _L)m_LM_L^1m_L^T,$$ $`(1.9)`$ $$M(\nu _R)m_RM_R^1m_R^T.$$ $`(1.10)`$ The scenario corresponding to $`O(m_LM_L^1m_L^T)O(m_RM_R^1m_R^T)`$ has already been investigated by one of the authors (Y.K.) . He has concluded that although either the atmospheric or solar neutrino data can be explained by the mixings $`\nu _\mu \nu _\tau `$ or $`\nu _e\nu _\mu `$, however, simultaneous explanation of the both data cannot be obtained in this model. In the present paper, we consider another possibility $`O(m_LM_L^1m_L^T)O(m_RM_R^1m_R^T)`$. In this case, mixings between $`\nu _{iL}`$ and $`\nu _{iR}`$ are induced. The solar neutrino data are understood from a small mixing between $`\nu _{eL}`$ and $`\nu _{eR}`$. The atmospheric and the LSND neutrino data are explained by the mixings $`\nu _{\mu L}\nu _{\tau L}`$ and $`\nu _{eL}\nu _{\mu L}`$, respectively. The vantage point of the democratic seesaw model is that parameters $`z_i`$ in the mass matrices $`m_L`$ and $`m_R`$ are given in terms of the charged lepton masses and thereby the mass spectrum and mixings of $`\nu _{iL}`$ and $`\nu _{iR}`$ can also be predicted in terms of the charged lepton masses. ## 2 Parameter $`b_\nu `$ In the present paper, for simplicity, we assume that all the neutral heavy fermion mass matrices $`M_D`$, $`M_L`$ and $`M_R`$ have the same flavor structure $$\frac{1}{\lambda _D}M_D=\frac{1}{\lambda _L}M_L=\frac{1}{\lambda _R}M_R=m_0(\mathrm{𝟏}+3b_\nu X),$$ $`(2.1)`$ and we will investigate only the case $`b_\nu =1/2`$. The excuse for considering only the case $`b_\nu =1/2`$ is as follows. The choices of $`b_f`$ ($`b_e=0,b_u=1/3,b_d1`$) have given the successful description of the quark masses and mixings in terms of the charged lepton masses. When, instead of the expression (1.3), we denote $`M_F`$ as $$M_F=m_0\lambda _f\sqrt{1+2b_f+3b_f^2}(\mathrm{cos}\varphi _fE\mathrm{sin}\varphi _fS),$$ $`(2.2)`$ $$E=\frac{1}{\sqrt{3}}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),S=\frac{1}{\sqrt{6}}\left(\begin{array}{ccc}0& 1& 1\\ 1& 0& 1\\ 1& 1& 0\end{array}\right),$$ $`(2.3)`$ where $`E`$ and $`S`$ have been normalized as $`\mathrm{Tr}E^2=\mathrm{Tr}S^2=1`$ and $`\mathrm{tan}\varphi _f=\sqrt{2}b_f/(1+b_f)`$, the cases $`b_e=0`$, $`b_u=1/3`$ and $`b_d=1`$ correspond to $`(\mathrm{cos}\varphi _f,\mathrm{sin}\varphi _f)=(1,0)`$, $`(\sqrt{2/3},\sqrt{1/3})`$ and $`(0,1)`$, respectively. Considering an empirical relation $`\varphi _d=\pi /2\varphi _e`$ for $`(\mathrm{cos}\varphi _e,\mathrm{sin}\varphi _e)=(1,0)`$ and $`(\mathrm{cos}\varphi _d,\mathrm{sin}\varphi _d)=(0,1)`$, we consider that the value of $`b_\nu `$ is also given by $`\varphi _\nu =\pi /2\varphi _u`$ for $`(\mathrm{cos}\varphi _u,\mathrm{sin}\varphi _u)=(\sqrt{2/3},\sqrt{1/3})`$, i.e., we assume $$(\mathrm{cos}\varphi _\nu ,\mathrm{sin}\varphi _\nu )=(\sqrt{1/3},\sqrt{2/3}),$$ $`(2.4)`$ which corresponds to the case $`b_\nu =1/2`$. Besides, from the phenomenological point of view, the case $`b_\nu =1/2`$ is also interesting. The inverse matrix of the $`M_L`$ with $`b_\nu =1/2`$ $$M_L=m_0\lambda _L(\mathrm{𝟏}\frac{1}{2}3X)=\frac{1}{2}m_0\lambda _L\left(\begin{array}{ccc}1& 1& 1\\ 1& 1& 1\\ 1& 1& 1\end{array}\right),$$ $`(2.5)`$ is given by $$M_L^1=\frac{1}{m_0\lambda _L}\left(\begin{array}{ccc}0& 1& 1\\ 1& 0& 1\\ 1& 1& 0\end{array}\right),$$ $`(2.6)`$ so that the seesaw matrix $`M_\nu m_LM_L^1m_L^T`$ is expressed as $$M_\nu m_0\frac{1}{\lambda _L}\left(\begin{array}{ccc}0& z_1z_2& z_1z_3\\ z_1z_2& 0& z_2z_3\\ z_1z_3& z_2z_3& 0\end{array}\right).$$ $`(2.7)`$ The form (2.7) is just a Zee-type mass matrix , which has recently been revived as a promising neutrino mass matrix form. ## 3 Mass spectrum and mixing For the specific form (2.1) with $`b_\nu =1/2`$, the 6$`\times `$6 seesaw matrix $`M^{(6\times 6)}`$ given by Eq. (1.8) becomes $$M^{(6\times 6)}m_0\left(\begin{array}{cc}0& Z\\ \kappa Z& 0\end{array}\right)\left(\begin{array}{cc}\lambda _RY& \lambda _DY\\ \lambda _DY& \lambda _LY\end{array}\right)^1\left(\begin{array}{cc}0& \kappa Z\\ Z& 0\end{array}\right)$$ $$=m_0\frac{1}{\lambda _R\lambda _L\lambda _D^2}\left(\begin{array}{cc}\lambda _RZY^1Z& \kappa \lambda _DZY^1Z\\ \kappa \lambda _DZY^1Z& \kappa ^2\lambda _LZY^1Z\end{array}\right),$$ $`(3.1)`$ where $$Y=\mathrm{𝟏}+3b_\nu X,Y^1=\mathrm{𝟏}+3a_\nu X,$$ $`(3.2)`$ $$a_\nu =b_\nu /(1+3b_\nu ).$$ $`(3.3)`$ Therefore, the matrix $`M^{(6\times 6)}`$ is diagonalized by the 6$`\times `$6 unitary matrix $`U^{(6\times 6)}`$ $$U^{(6\times 6)}=\left(\begin{array}{cc}\mathrm{cos}\theta U& \mathrm{sin}\theta U\\ \mathrm{sin}\theta U& \mathrm{cos}\theta U\end{array}\right),$$ $`(3.4)`$ as $$U^{(6\times 6)}M^{(6\times 6)}U^{(6\times 6)}=\mathrm{diag}(m_{\nu _{1L}},m_{\nu _{2L}},m_{\nu _{3L}},m_{\nu _{1R}},m_{\nu _{2R}},m_{\nu _{3R}})$$ $$=m_0\mathrm{diag}(\xi _L\rho _1,\xi _L\rho _2,\xi _L\rho _3,\xi _R\rho _1,\xi _R\rho _2,\xi _R\rho _3),$$ $`(3.5)`$ where $$U^{}ZY^1ZU=\mathrm{diag}(\rho _1,\rho _2,\rho _3),$$ $`(3.6)`$ $$\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{cc}\lambda _R& \kappa \lambda _D\\ \kappa \lambda _D& \kappa ^2\lambda _L\end{array}\right)\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)=\left(\begin{array}{cc}\lambda _L^{}& 0\\ 0& \lambda _R^{}\end{array}\right),$$ $`(3.7)`$ $$\xi _L=\frac{\lambda _L^{}}{\lambda _R\lambda _L\lambda _D^2},\xi _R=\frac{\lambda _R^{}}{\lambda _R\lambda _L\lambda _D^2},$$ $`(3.8)`$ $$\left(\begin{array}{c}\lambda _L^{}\\ \lambda _R^{}\end{array}\right)=\frac{1}{2}(\lambda _R+\kappa ^2\lambda _L)\frac{1}{2}(\lambda _R\kappa ^2\lambda _L)\sqrt{1+\mathrm{tan}^22\theta }.$$ $`(3.9)`$ The mixing angle $`\theta `$ between $`\nu _{iL}`$ and $`\nu _{iR}`$ is given by $$\mathrm{tan}2\theta =\frac{2\kappa \lambda _D}{\lambda _R\kappa ^2\lambda _L}.$$ $`(3.10)`$ The light neutrino masses $`m(\nu _{iL})`$ and $`m(\nu _{iR})`$ are given by $$m(\nu _{iL})=m_0\xi _L\rho _i,m(\nu _{iR})=m_0\xi _R\rho _i.$$ $`(3.11)`$ For the case of $`b_\nu =1/2`$, the eigenvalues $`\rho _i`$ of the matrix $`ZY^1Z`$ are given by $$\rho _12z_1^2,\rho _2\left(z_2+\frac{z_1^2}{2z_2}z_1^2\right),\rho _3z_2+\frac{z_1^2}{2z_2}+z_1^2,$$ $`(3.12)`$ so that $$\rho _3^2\rho _2^24z_2z_1^2,\rho _2^2\rho _1^2z_2^2.$$ $`(3.13)`$ The $`3\times 3`$ mixing matrix $`U`$ for the case $`b_\nu =1/2`$ is given by $$U\left(\begin{array}{ccc}1& \frac{1}{\sqrt{2}}\frac{z_1}{z_2}(1z_2)& \frac{1}{\sqrt{2}}\frac{z_1}{z_2}(1+z_2)\\ \frac{z_1}{z_2}& \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\\ z_1& \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\end{array}\right).$$ $`(3.14)`$ ## 4 Explanations of the neutrino data The atmospheric and solar neutrino data are explained by the mixings $`\nu _{\mu L}\nu _{\tau L}`$ and $`\nu _{eL}\nu _{eR}`$, respectively. As seen in the mixing matrix (3.14), the neutrinos $`\nu _{\mu L}`$ and $`\nu _{\tau L}`$ are maximally mixed. On the other hand, the mixing between $`\nu _{eL}`$ and $`\nu _{eR}`$ is given by Eq. (3.10). Since the solar neutrino data disfavor sterile neutrino with a large mixing angle, we take the small mixing angle solution in the Mikheyev-Smirnov-Wolfenstein (MSW) mechanism , $$\mathrm{\Delta }m_{solar}^24.0\times 10^6\mathrm{eV}^2,\mathrm{sin}^22\theta _{solar}6.9\times 10^3.$$ $`(4.1)`$ Here, the values in Eq. (4.1) have been quoted from the recent analysis for $`\nu _e\nu _s`$ by Bahcall, Krastev and Smirnov . The value $`\mathrm{sin}^22\theta _{solar}7\times 10^3`$ can be fitted by adjusting the parameters $`\lambda _L`$, $`\lambda _R/\kappa ^2`$ and $`\lambda _D/\kappa `$ in Eq. (3.10). As seen from Eqs. (3.5) and (3.13), the ratio of $`\mathrm{\Delta }m_{solar}^2=(m_{\nu _{1R}})^2(m_{\nu _{1L}})^2`$ to $`\mathrm{\Delta }m_{atm}^2=(m_{\nu _{3L}})^2(m_{\nu _{2L}})^2`$ is given by $$\frac{\mathrm{\Delta }m_{solar}^2}{\mathrm{\Delta }m_{atm}^2}\frac{\lambda _R^2\lambda _L^2}{\lambda _L^2}\frac{4z_1^4}{4z_2z_1^2}(R^21)\frac{m_e}{\sqrt{m_\mu m_\tau }}=(R^21)\times 1.15\times 10^3,$$ $`(4.2)`$ where $$R=\frac{\lambda _R^{}}{\lambda _L^{}}=\frac{\xi _R}{\xi _L}=\frac{m(\nu _{iR})}{m(\nu _{iL})}.$$ $`(4.3)`$ The recent best fit value $`\mathrm{\Delta }m_{atm}^2=3.2\times 10^3`$ eV<sup>2</sup> gives the ratio $$\frac{\mathrm{\Delta }m_{solar}^2}{\mathrm{\Delta }m_{atm}^2}\frac{4.0\times 10^6\mathrm{eV}^2}{3.2\times 10^3\mathrm{eV}^2}1.3\times 10^3.$$ $`(4.4)`$ By comparing Eqs. (4.2) and (4.4), we obtain $`R1.4`$. Note that the observed value (4.4) is in good agreement with the value $`m_e/\sqrt{m_\mu m_\tau }`$, so that we are tempted to consider a model with $`R0`$. However, the sign of $`\mathrm{\Delta }m_{solar}^2`$ in the small mixing angle MSW solution must be positive, so that we cannot consider the case $`R0`$. In the present model, $`R`$ is only a phenomenological parameter with the constraint $`R>1`$. The LSND data is explained by the mixing $`\nu _{eL}\nu _{eR}`$. The mass-squared difference $`\mathrm{\Delta }m_{LSND}^2=m_{\nu _{2L}}^2m_{\nu _{1L}}^2`$ and the $`\nu _{eL}\nu _{eR}`$ mixing angle are given by the ratio $$\frac{\mathrm{\Delta }m_{LSND}^2}{\mathrm{\Delta }m_{atm}^2}\frac{z_2}{4z_1}\frac{1}{4}\sqrt{\frac{m_\mu }{m_e}}=2.2\times 10^2,$$ $`(4.5)`$ and $$\mathrm{sin}^22\theta _{LSND}4U_{e1}^2U_{\mu 1}^24\left(\frac{z_1}{z_2}\right)^24\frac{m_e}{m_\mu }0.019,$$ $`(4.6)`$ respectively. The best fit value $`\mathrm{\Delta }m_{atm}^23.2\times 10^3`$ $`\mathrm{eV}^2`$ give a prediction $`\mathrm{\Delta }m_{LSND}^20.70`$ $`\mathrm{eV}^2`$. However, the region $`\mathrm{\Delta }m_{LSND}^20.34`$ $`\mathrm{eV}^2`$ in the LSND favored region at $`\mathrm{sin}^22\theta =0.02`$ has been excluded by the recent KARMEN2 experiment . Therefore, only when we take the value $`\mathrm{\Delta }m_{LSND}^20.33`$ $`\mathrm{eV}^2`$, we can obtain the prediction $`\mathrm{\Delta }m_{atm}^21.5\times 10^3`$ $`\mathrm{eV}^2`$ which is barely inside the 90% C.L. allowed region ($`1.5\times 10^3\mathrm{eV}^2\mathrm{\Delta }m_{atm}^25\times 10^3\mathrm{eV}^2`$) in the recent Super-Kamiokande atmospheric neutrino data . Hereafter, we will adopt this pinpoint solution: $$\mathrm{\Delta }m_{LSND}^20.33\mathrm{eV}^2,\mathrm{\Delta }m_{atm}^21.5\times 10^3\mathrm{eV}^2.$$ $`(4.7)`$ Then, the parameter $`R`$ is fixed as $$R1.8,$$ $`(4.8)`$ from Eq. (4.2), and the neutrino masses are predicted as follows: $$m(\nu _{3L})m(\nu _{2L})0.57\mathrm{eV},m(\nu _{1L})1.3\times 10^3\mathrm{eV},$$ $`(4.9)`$ $$m(\nu _{3R})m(\nu _{2R})1.05\mathrm{eV},m(\nu _{1R})2.5\times 10^3\mathrm{eV},$$ $`(4.10)`$ where we have used the relation $`m(\nu _{2L})\sqrt{\mathrm{\Delta }m_{LSND}^2}`$. In the present scenario, there are three light sterile neutrinos $`\nu _{iR}(i=1,2,3)`$. However, those neutrinos do not spoil the big bang nucleosynthesis (BBN) scenario, which puts the following constraint for a mixing between the active neutrino $`\nu _\alpha (\alpha =e,\mu ,\tau )`$ and a sterile neutrino $`\nu _s`$, $$(\mathrm{sin}^22\theta _{\alpha s})^2\mathrm{\Delta }m_{\alpha s}^2<3.6\times 10^4\mathrm{eV}^2.$$ $`(4.11)`$ The value of $`(\mathrm{sin}^22\theta )^2\mathrm{\Delta }m^2`$ in our model is less than $`10^4\mathrm{eV}^2`$, because the mixing angle $`\theta `$ in the present model is sufficiently small, i.e., $`(\mathrm{sin}^22\theta )^2=(6.9\times 10^3)^2=4.8\times 10^5`$. However, we have another severe constraint on the neutrino masses from the cosmic structure formation in a low-matter-density universe $$N_\nu m_\nu <1.8\mathrm{eV}(1.5\mathrm{eV}),$$ $`(4.12)`$ for flat universe (for open universes), where $`N_\nu `$ is the number of almost degenerate neutrinos with the highest mass. The present model gives $`N_\nu m_\nu 3.2`$ $`\mathrm{eV}`$, so that the model dose not satisfies the constraint (4.12). We will go optimistically for this problem. The mixing between $`\nu _{eL}`$ and $`\nu _{\tau L}`$ is given by $$U_{e3}\frac{1}{\sqrt{2}}\frac{z_1}{z_2}(1+z_2)\sqrt{\frac{m_e}{2m_\mu }}\left(1+\sqrt{\frac{m_\mu }{m_\tau }}\right)0.061,$$ $`(4.13)`$ which safely satisfies the constraint $`\|U_{e3}\|(0.220.14)`$ obtained from the CHOOZ reactor neutrino experiment . ## 5 Conclusion and discussion In conclusion, we have investigated a neutrino mass matrix in the framework of the “democratic” universal seesaw model. Although the model has three light sterile neutrinos $`\nu _{iR}`$ $`(i=1,2,3)`$, they do not spoil the BBN scenario, because the mixing angle $`\theta `$ between the active and sterile neutrinos is taken as $`\mathrm{sin}^22\theta 7\times 10^3`$. The atmospheric, solar and LSND neutrino data are explained by the mixings $`\nu _{\mu L}\nu _{\tau L}`$, $`\nu _{eL}\nu _{eR}`$ and $`\nu _{eL}\nu _{\mu L}`$, respectively. The model with the parameter $`b_\nu =1/2`$ gives the predictions in terms of the charged lepton masses, $$\frac{\mathrm{\Delta }m_{solar}^2}{\mathrm{\Delta }m_{atm}^2}(R^21)\frac{m_e}{\sqrt{m_\mu m_\tau }},\frac{\mathrm{\Delta }m_{LSND}^2}{\mathrm{\Delta }m_{atm}^2}\frac{1}{4}\sqrt{\frac{m_\mu }{m_e}},$$ $`(5.1)`$ $$\mathrm{sin}^22\theta _{atm}1,\mathrm{sin}^22\theta _{LSND}4\frac{m_e}{m_\mu },$$ $`(5.2)`$ where $`R=m(\nu _{iR})/m(\nu _{iL})`$. In the present model, the prediction $`\mathrm{\Delta }m_{solar}^2/\mathrm{\Delta }m_{atm}^2`$ includes a free parameter $`R`$. Only a parameter independent prediction is $`\mathrm{\Delta }m_{LSND}^2/\mathrm{\Delta }m_{atm}^2`$ together with $`\mathrm{sin}^22\theta _{LSND}4m_e/m_\mu `$. Since the most part of the allowed region of the $`\nu _e`$-$`\nu _\mu `$ oscillation in the LSND data is ruled out by the KARMEN2 data (but a narrow region still remains), the predictability of the present model is somewhat faded from the point of view of the neutrino phenomenology. However, the motivation of the present paper is not to give the explanation of the LSND data, but to seek for a possible unification model of the quark and lepton mass matrices. The presence of the light right-handed neutrinos $`\nu _{iR}`$ will offer fruitful new physics to the near future neutrino experiments. In the present scenario, the following intermediate energy scales have been considered: The neutral leptons $`N_L`$ and $`N_R`$ acquire large Majorana masses $`M_R`$ and $`M_L`$ at $`\mu =\mathrm{\Lambda }_{NL}=m_0\lambda _R`$ and $`\mu =\mathrm{\Lambda }_{NR}=m_0\lambda _L`$, respectively. The fermions $`N`$ and $`F`$ $`(F=U,D,E)`$ acquire large Dirac masses $`M_D`$ and $`M_F`$ at $`\mu =\mathrm{\Lambda }_D=m_0\lambda _D`$ and $`\mu =\mathrm{\Lambda }_F=m_0\lambda _F`$, respectively. The gauge symmetries SU(2)<sub>R</sub> and SU(2)<sub>L</sub> are broken at $`\mu =\mathrm{\Lambda }_R=m_0\kappa `$ and $`\mu =\mathrm{\Lambda }_L=m_0`$, respectively. For $`\mathrm{tan}^22\theta 1`$, form Eq. (3.9), we obtain the approximate relations $$\lambda _L^{}\kappa ^2\lambda _L,\lambda _R^{}\lambda _R,$$ $`(5.3)`$ so that $$R=\frac{\lambda _R^{}}{\lambda _L^{}}\frac{\lambda _R}{\kappa ^2\lambda _L}.$$ $`(5.4)`$ The numerical result $`R=O(1)`$ means $`\lambda _R/\lambda _L\kappa ^2`$, i.e., $$\frac{\mathrm{\Lambda }_{NL}}{\mathrm{\Lambda }_{NR}}\left(\frac{\mathrm{\Lambda }_R}{\mathrm{\Lambda }_L}\right)^2.$$ $`(5.5)`$ Since $$m(\nu _{2L})=\xi _L\rho _2m_0\frac{\kappa ^2\lambda _L}{\lambda _R\lambda _L\lambda _D^2}\rho _2m_0\frac{1}{\lambda _R/\kappa ^2}\sqrt{\frac{m_\mu }{m_\tau }}m_0,$$ $`(5.6)`$ we estimate $$\frac{\lambda _R}{\kappa ^2}\sqrt{\frac{m_\mu }{m_\tau }}\frac{m_0}{m(\nu _{2L})}10^{11},$$ $`(5.7)`$ where we have used $`m_010^2`$ GeV, so that we obtain $`\mathrm{\Lambda }_{NL}\kappa ^2\times 10^{13}`$ GeV. If we consider that $`\mathrm{\Lambda }_{NL}`$ must be smaller than the Planck mass $`M_P10^{19}`$ GeV, we obtain the constraint $$\kappa \mathrm{\Lambda }_R/\lambda _L<10^3.$$ $`(5.8)`$ Since the case $`\kappa 1`$ is experimentally ruled out, we conclude that $$O(10^3)\mathrm{GeV}<\mathrm{\Lambda }_R<O(10^5)\mathrm{GeV}.$$ $`(5.9)`$ From (3.10), we estimate $$\frac{\lambda _D}{\kappa }\frac{1}{2}\left(1\frac{1}{R}\right)\frac{\lambda _R}{\kappa ^2}\mathrm{tan}2\theta 10^9.$$ $`(5.10)`$ On the other hand, we have known that $$\frac{\mathrm{\Lambda }_R}{\mathrm{\Lambda }_F}=\frac{\kappa }{\lambda _F}10^2$$ $`(5.11)`$ from the study of the quark mass spectrum . Therefore, we cannot take an idea that the Dirac masses $`M_D`$ and $`M_F`$ ($`FN`$) are generated at the same energy scale $`\mu =\mathrm{\Lambda }_D=\mathrm{\Lambda }_F`$. Note that in the conventional universal seesaw model, the neutrino masses are of the order of $`\mathrm{\Lambda }_L^2/\mathrm{\Lambda }_{NR}=m_0/\lambda _L`$, because of $`M(\nu _L)m_LM_L^1m_L^T`$, so that we consider $`\lambda _L10^9`$. In contrast with the conventional model, in the present model, the value of $`\lambda _L`$ is $`\lambda _L\lambda _R/\kappa ^210^{11}`$. Therefore, for example, the conclusion on the intermediate energy scales based on the SO(10)$`{}_{L}{}^{}\times `$SO(10)<sub>R</sub> model in Ref. is not applicable to the present model, because in Ref. the solutions have been investigated under the condition $`\lambda _L10^9`$. It is a future task to seek for a unification model which satisfies these constraints on the intermediate energy scales, (5.5) and (5.7)-(5.11). Acknowledgments One of the authors (Y.K.) would like to thank Professor O. Yasuda for his helpful comments on the cosmological constraints on the neutrino masses and informing the references and . He also thanks Professor M. Tanimoto and Professor A. Yu. Smirnov for pointing out a mistake (the sign of the $`\mathrm{\Delta }m_{solar}^2`$) in the first version of the paper. A.G. is supported by the Japan Society for Promotion of Science (JSPS), Postdoctoral Fellowship for Foreign Researches in Japan (Grant No. 99222).
warning/0006/cond-mat0006372.html	ar5iv	text	# Symmetry and designability for lattice protein models ## I Introduction The folded structures of proteins are often highly ordered. They are comprised of secondary structures, and often have striking regularities in their tertiary organization . What is the origin of symmetry in natural proteins? We approach this question by exploring the symmetries in simple lattice models of protein folding. Lattice models for proteins have been a rich source of information on protein structure. Yue and Dill observed certain protein-like secondary structures and tertiary symmetries in HP lattice model proteins that have low degeneracies, i.e., a small number of low energy states. More recently, Li et al. noticed that the most “designable” structures, namely those with a large number of sequences folding into them, also often have global symmetries. Since the most designable structures also have other protein-like properties – they have sharp thermal folding transitions and are fast folders , the connection to symmetry is intriguing. In these earlier studies, no quantitative measure was used to define symmetry. Here, we explore in detail the connection between designability and global symmetry, based on a quantitative, but simple, measure of symmetry. Within the hydrophobic model , we quantify the relation between designability and symmetry for 6x6 compact lattice proteins. This article is organized as follows: Sec. II reviews the hydrophobic model and the designabilities of structures. In Sec. III, we relate symmetry to designability and identify the importance of the surface-core pattern to the particular emerging symmetries. To understand the origin of enhanced symmetry, in Sec. IV we explore, first, the role of surface-to-core transitions and, second, the extent to which symmetric folds result from the repeated use of designable substructures. For comparison, Sec. V addresses symmetry in a model not based on hydrophobicity. Sec. VI is the summary and conclusion. ## II Hydrophobic Model In this section, we review the hydrophobic model and the designabilities of structures. For more details, the reader is referred to Li et al. . The hydrophobic model is a combination of the HP model and the solvation model citeEisenberg86. In an HP model, the twenty different amino acids of proteins are replaced by two monomer types, Hydrophobic or Polar, according to their affinities for water. Each protein is therefore a sequence of H’s and P’s. In a lattice HP model, the amino acids are restricted to fall only on the sites of a regular lattice, typically a square lattice in two dimensions or a cubic lattice in three dimensions. The allowed conformations are self-avoiding, and hence cannot visit a single lattice site more than once. Here we use a variant that we call the hydrophobic model, in which only the maximally compact structures are considered as possible ground states. This simplification still allows us to capture the essence of the HP model, but with two advantages: a substantial reduction in computational cost, and a conceptually useful method to represent sequences and structures within the same kind of abstract spatial representation, described below. In the hydrophobic model, the energy of a compactly folded protein is taken to be simply minus the number of H monomers in the “core” (cf. Fig. 1). Therefore, in the hydrophobic model, the energy of an HP-protein sequence folded into a particular compact structure depends only on the structure’s ordering of surface and core sites. Thus, a structure can be represented by a string $`𝐬`$ of 0’s and 1’s: sites in the core region are represented by 1’s and sites on the surface are represented by 0’s, as illustrated in Fig. 1. Sequences are also represented by strings of 0’s (P) and 1’s (H), $`𝐡=(h_1,h_2,\mathrm{},h_N)`$, where $`h_i`$ denotes the hydrophobicity of the monomer at position $`i`$ of the sequence. The energy of a sequence folded into a particular structure is therefore given by $`H={\displaystyle \underset{i=1}{\overset{N}{}}}s_ih_i`$ where $`s_i`$ is the structure string. An equivalent way of writing the energy is, $`H={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}(s_ih_i)^2{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}s_i^2{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}h_i^2.`$ The number of core sites is the same for all structures of the same size, thus $`\frac{1}{2}s_i^2`$ is a constant and can be dropped. Similarly, the last term, $`\frac{1}{2}h_i^2`$, is a constant for each sequence and so does not influence which structure is the ground state for that sequence. Therefore, the only relevant term is the first term, which measures the Hamming distance between the structure string and the sequence string in an $`N`$-dimensional Euclidean space. A sequence with string $`𝐡`$ will have a particular structure with string $`𝐬`$ as its unique ground state if and only if $`𝐡`$ is closer to $`𝐬`$ than to any $`𝐬^{}`$ corresponding to another structure. The designability of a structure can therefore be obtained from the following geometric construction: Draw bisector planes between $`𝐬`$ and all of its neighboring structures in the $`N`$-dimensional space. The volume enclosed by these planes is called the Voronoi polytope around $`𝐬`$. The designability of a structure is the number of sequences lying entirely within the Voronoi polytope around that structure. This is schematically represented in Fig. 2. Each vertex represents a sequence. Those vertices corresponding to structures are circled. Intuitively, the designability of a structure is closely related to how far away its nearest neighbors are. The further away its neighbors are, the more designable it is. The histogram of the number of structures versus designability for the 6x6 hydrophobic model is shown in Fig. 3. The distribution has a long tail of highly designable structures compared to a Poisson distribution with the same mean. If sequences were randomly assigned to structures, the resulting distribution of designabilities would be Poisson. It is clear from Fig. 3 that the structures in the tail have anomalously high designabilities. That is, they are unique ground states of many more than their share of sequences. The many sequences that have a particular highly designable native structure are related to each other by point mutations . For the model we consider, a point mutation is simply the replacement of a hydrophobic monomer “1” by a polar monomer “0”, or vice versa. Often, many monomers can be independently mutated without destabilizing the native state . Therefore, the folding of these sequences is relatively insensitive to mutations. One can think of highly designable structures as those which remain most stable under sequence mutations. ## III Symmetry and Designability In Li et al. , it was noted that highly designable structures tend to be highly symmetric, with global mirror symmetries as well as regular local motifs. In this study, we explore the connection between designability and symmetry in detail. We focus on 6x6 2D square-lattice proteins. To measure the symmetry of a structure, we look at how well that structure is preserved under rigid global transformations. Specifically, the transformations are the mirror reflections about the x/y axes, the mirror reflections about the two diagonal directions, and $`90^o`$ and $`180^o`$ rotations. The symmetry scores for a given structure are the number of overlapping bonds between that structure and each of its transformed versions. The maximum possible symmetry score for a 6x6 compact structure is $`35`$. ### A Hydrophobic Model with Centered Core We begin by studying the trends of symmetry versus designability for the hydrophobic model. The symmetry scores, averaged over designability bins, are plotted versus the designability in Fig. 4. It is observed that, on average, the x/y-mirror symmetry (the larger of the x-mirror symmetry score and the y-mirror symmetry score) increases with designability (Fig. 4(a)). A similar trend is observed for $`180^o`$ rotation symmetry, which is consistent with the x/y-mirror symmetry result since a $`180^o`$ rotation is simply an x-mirror operation followed by a y-mirror operation. For the other symmetry operations, $`90^o`$ rotation and diagonal mirrors, the trend is reversed – higher designability implies lower symmetry scores for these symmetries. Thus, for the hydrophobic model, there is indeed a connection between designability and symmetry as previously noted . However, different symmetries behave differently with increasing designability. In this case, the x/y-mirror symmetry and $`180^o`$ rotation symmetry are enhanced for highly designable structures. ### B Hydrophobic model with shifted core For the hydrophobic model, the surface-core pattern has the symmetry of a square (Fig. 1). What if this is not the case? Does higher designability always lead to higher x/y-mirror symmetry scores, even when the surface-core pattern is disrupted? To address this question, we study a shifted-core version of the hydrophobic model. The core sites have been shifted to the lower corner (Fig. 5). In this shifted-core model, the energy of a compactly folded protein is taken to be simply minus the number of H monomers in the new off-center “core” . The histogram of designability for the shifted-core model is shown in Fig. 6. Again, the region of high designability is characterized by a long tail, qualitatively similar to that of the ordinary 6x6 model. With the “core” shifted to one corner, the surface-core pattern no longer has mirror symmetries about the x and y axes. In fact, only one diagonal mirror symmetry is left. In Fig. 7, we plot the averaged symmetry scores versus designability for the shifted-core model. When x/y-mirror symmetry is plotted versus designability there is no correlation. Instead, only the diagonal-mirror symmetry which is present in the surface-core pattern increases significantly with designability (Fig. 7(a)). These results indicate that the surface-core pattern is important in determining which symmetries are favored in highly designable structures. In both cases considered, the preferred symmetries follow the surface-core pattern. ## IV Origin of Symmetry Why is there an enhancement of symmetry for highly designable structures? Also, why are some symmetries enhanced and not others? In this section, we examine two possible origins of symmetry. First, perhaps global symmetries arise in designable structures because of a high number of surface-to-core transitions. Or second, perhaps designable structures have global symmetries because they arise from repeated highly designable substructures. ### A Surface-to-core transitions A possible candidate for the link between designability and symmetry is a local property of structures – the number of surface-to-core transitions. A surface-to-core transition occurs when monomer $`i`$ of the chain is in the core and monomer $`i+1`$ is on the surface, or vice versa. Previously , it was observed that highly designable structures have an excess of surface-to-core transitions. The connection can be understood as follows: (1) A structure with a large number of surface-to-core transitions is difficult to rearrange without exchanging many surface and core sites. Such structures are therefore likely to be far from their neighbors in the space of strings, and thus have a chance for high designability (cf. Fig. 2). (2) In turn, a structure with a large number of surface-to-core transitions has a geometrical regularity which may naturally lead to global symmetry. Moreover, the geometrical regularities, and hence the enhanced global symmetry, should reflect the symmetry of the surface-core pattern, consistent with our results using the shifted-core model. We tested whether surface-to-core transitions form the link between designability and global symmetry. We find that, qualitatively, both correlations (1) and (2) are present, however, quantitatively, they fail to account for the observed enhancement of global symmetry. To quantify the first correlation, the number of surface-to-core transitions averaged over structures of a given range of designabilities is plotted against designability in Fig. 8(a) for the original hydrophobic model. High designability clearly implies an enhanced number of surface-to-core transitions. To demonstrate the second correlation, the x/y symmetry score averaged over structures with a given number of surface-to-core transitions is plotted against the number of surface-to-core transitions in Fig. 8(b). Symmetry does increase with the number of surface-to-core transitions when the number of transitions is large . Is the chain of correlations from designability to surface-to-core transitions to global symmetry strong enough to explain the observed enhancement of global symmetry? In Fig. 8(c), the x/y-symmetry score is plotted against designability assuming the connection between them is only through the correlation of each with the number of surface-to-core transitions. Specifically, for a given designability, the corresponding average number of surface-to-core transitions is obtained from panel (a), then the corresponding average x/y-symmetry score for that number of surface-to-core transitions is obtained from panel (b). The predicted x/y-symmetry score thus obtained is then plotted against designability. The actual x/y-symmetry score versus designability from Fig. 4(a) is also plotted. We see that surface-to-core transitions account for only a fraction of the observed connection between symmetry and designability. ### B Designable substructures A second possible explanation for why designable folds are so symmetric is (1) they arise from designable substructures, and (2) symmetries are a natural consequence of assembling anything from identical substructures. The most designable structure for the 6x6 hydrophobic model is shown in the left half of Fig. 9 . We take the right half of the 6x6 surface-core pattern (a 3x6 rectangle) and calculate the designabilities of all possible structures for this 3x6 hydrophobic model. The most designable 3x6 structure is shown in the right half of Fig. 9. Comparing the two structures, we see that the most designable 3x6 structure is very similar to one half of the most designable 6x6 structure of the original model. We conclude that this 6x6 structure is highly designable because it is composed of two highly designable substructures. The role of symmetry in this case can then be understood as duplicating a winning solution. It is not yet clear how to quantify this concept of designable substructures. Any scheme that involves breaking and reforming bonds, as would be necessary to relate the structures in Fig. 9, seems arbitrary and unsatisfactory. Nevertheless, a connection between designable components and global symmetry seems to us likely, and may have implications for understanding global symmetries in real proteins. ## V Symmetry beyond the hydrophobic model As a final question, we ask if the connection between designability and symmetry is particular to models based on hydrophobicity, or whether it occurs more generally. In the hydrophobic model each structure is characterized by its ordering of surface and core sites. As an alternative, we consider a model in which each structure is characterized by its complete contact matrix, as described below. Those structures with large contact-matrix distance from their neighbors are considered to be highly designable. Within the contact-matrix model, the designable structures do not show significantly enhanced symmetry. Each structure has a contact matrix for its monomers. The elements of the contact matrix between monomers are 1 if they are next to each other in the structure, but not adjacent on the chain, and 0 otherwise. A compact structure is uniquely defined by its contact matrix, up to rigid rotations and inversion. Thus, contact matrices and structures are related by a one-to-one mapping. The distance between any two structures can be measured by the overlap between their contact matrices. The more overlap, the more bonds they have in common. Hence, contact-matrix distance measures the similarity of structures without particular emphasis on surface and core sites. Just as in the hydrophobic model, where structures with few neighboring strings emerged as highly designable, structures with few neighbors in contact-matrix distance would emerge as highly designable for more general models of amino-acid interaction . For the set of compact 6x6 structures, we take the number of neighbors within a contact-matrix distance of 16 as a measure designability, with low number of neighbors implying high designability. The histogram of number of structures versus number of near neighbors is shown in Fig. 10(a). Shown in Fig. 10(b) is a plot of the averaged top symmetry scores versus the number of neighboring structures within the cutoff distance of $`16`$. The top symmetry score for a structure is the highest score for all possible rotations and reflections. The horizontal line indicates the average top symmetry score of $`27.7`$. The region of few neighbors, and hence high designability, is at the left of the figure and has only a very slightly enhanced symmetry score with respect to the average. We conclude that enhanced global symmetry of designable structures does not emerge generally from models with arbitrary interactions among amino acids. Rather, the enhancement of global symmetries is particular to models in which the interaction between amino acids is dominated by hydrophobicity . It appears that the correlation between the designability and symmetry of a native protein is a consequence of the key role played by hydrophobic solvation, and the approximate radial symmetries that result from it. ## VI Conclusion In this work, we have examined the connection between the designability and symmetry of protein structures within the hydrophobic model of Li et al. . The designable structures, namely those which are unique ground states of many more than their share of sequences, had been previously identified to have enhanced global symmetry, as well as other protein-like attributes such as thermodynamic stability and stability against mutations . To quantify the relation between symmetry and designability we focused on the set of two-dimensional compact structures which fill the sites of a 6x6 square lattice. We found that the designable structures have strongly enhanced symmetry for x/y reflection and $`180^o`$ rotation. For a related model in which the “core” is shifted to one corner, the only enhanced symmetry was a diagonal reflection. This indicates that the enhanced symmetry of the designable structures reflects a symmetry of the surface-core pattern. To explore the origin of symmetry, we examined the relation between designability, number of surface-to-core transitions, and global symmetry. We conclude that an increase in surface-to-core transitions among designable structures can only account for a fraction of the observed enhancement of global symmetry. Our working hypothesis is that the global symmetry of designable structures results from the repetition of designable substructures. Finally, from a comparison model based on contact-matrix distances we conclude that the relation between designability and symmetry originates from surface-core symmetries, which in turn, result mainly from hydrophobic interactions.
warning/0006/hep-ph0006015.html	ar5iv	text	# Technique for Direct eV-Scale Measurements of the Mu and Tau Neutrino Masses Using Supernova Neutrinos \[ ## Abstract Early black hole formation in a core-collapse supernova will abruptly truncate the neutrino fluxes. The sharp cutoff can be used to make model-independent time-of-flight neutrino mass tests. Assuming a neutrino luminosity of $`10^{52}`$ erg/s per flavor at cutoff and a distance of 10 kpc, SuperKamiokande can detect an electron neutrino mass as small as 1.8 eV, and the proposed OMNIS detector can detect mu and tau neutrino masses as small as 6 eV. This Letter presents the first technique with direct sensitivity to eV-scale mu and tau neutrino masses. \] Introduction: Despite decades of experimental effort, the values of the neutrino masses remain elusive. While the laboratory bound on the electron neutrino mass is about 3 eV , the laboratory bounds on the mu and tau neutrino masses are much weaker: 170 keV and 18 MeV , respectively. Only recently have neutrino oscillation experiments found strong evidence for nonzero differences of squared neutrino masses. Once discovered, the values of the neutrino masses may provide important clues to physics beyond the Standard Model. In some scenarios, e.g., with the see-saw mechanism , the mu and tau neutrino masses are expected to be much larger than the electron neutrino mass. If they are at the eV scale or greater, the neutrino masses could also be important cosmologically as a component of the long-sought dark matter. It is therefore crucial to devise direct tests of the mu and tau neutrino masses with sensitivity reaching the eV scale. While neutrino mass tests based on cosmological considerations may reach the eV scale, they are indirect (no neutrinos are detected) and depend upon the other cosmological parameters being independently known . The best known possibility for directly measuring the mu and tau neutrino masses is by time-of-flight measurements of supernova neutrinos, comparing the arrival time of the mu and tau neutrinos to that of the electron neutrinos. However, this is complicated by the long intrinsic duration ($`10`$ s) of the neutrino signal and the fact that its detailed characteristics are model-dependent. Beacom and Vogel have shown that a technique based on the average arrival times $`t`$ is model-independent and is sensitive to delays as small as $`0.1`$. This would allow detection of mu or tau neutrino masses down to 45 eV in SuperKamiokande (SK) and 30 eV in the Sudbury Neutrino Observatory (SNO). If the mu and tau neutrino masses (strictly speaking, those of the relevant mass eigenstates) are nearly degenerate, as suggested by the atmospheric neutrino results , then the sensitivity would improve by about $`\sqrt{2}`$. Unfortunately, it seems difficult to improve the results with this technique, since the mass sensitivity scales with the detector mass $`M_D`$ as $`1/M_D^{1/4}`$ . To reach the few-eV scale would require detectors $`10^4`$ times larger, which seems impossible. In this Letter, we discuss a new time-of-flight technique for measuring neutrino masses that can reach the eV scale. This technique is applicable if the proto-neutron star forms a black hole early enough to abruptly terminate the neutrino signal. We state only our most important results; the details will be discussed at length in a forthcoming paper . Expected Neutrino Signal: We consider black hole formation which occurs soon ($`1`$ s) after core collapse (other scenarios are considered in Ref. ). Black hole formation is triggered by accretion, which drives the proto-neutron star mass above the maximum stable neutron star mass. The neutrino signal expected in this scenario has been studied by Burrows and Mezzacappa and Bruenn . In these models, the neutrino luminosities were fairly constant at more than $`10^{52}`$ erg/s per flavor until abruptly terminated by black hole formation. In fact, the transition should have a nonzero duration, of order the light crossing time $`2R/c0.1`$ ms, as the proto-neutron star radius shrinks to that of the final black hole. During the transition, the gravitational redshift, originally $`10\%`$, rapidly diverges, truncating the neutrino signal. Using a singularity-avoiding code, Baumgarte et al. studied the transition and found its duration to be 0.5 ms. Thus, we can consider the neutrino fluxes to be sharply and simultaneously terminated. The results below assume a luminosity $`L_{BH}=10^{52}`$ erg/s per flavor at the cutoff time $`t_{BH}`$, and a distance $`D=10`$ kpc. We assume the following temperatures: $`T=3.5`$ MeV for $`\nu _e`$, $`T=5`$ MeV for $`\overline{\nu }_e`$, and $`T=8`$ MeV for $`\nu _\mu `$, $`\nu _\tau `$, $`\overline{\nu }_\mu `$, and $`\overline{\nu }_\tau `$. It will be shown that the necessary quantities can be measured in a realistic situation. Neutrino Mass Effects: At lowest order, a neutrino with mass $`m`$ (in eV) and energy $`E`$ (in MeV) will have an energy-dependent delay (in s) relative to a massless neutrino in traveling over a distance D (in 10 kpc): $$\mathrm{\Delta }t(E)=0.515\left(\frac{m}{E}\right)^2D.$$ (1) The distance is scaled by the approximate distance to the Galactic center, though a supernova may be detected from anywhere in the Galaxy. For the smallest detectable masses, the delay effects will be visible only after the sharp cutoff, where no events are otherwise expected. Since the delays are very small, the luminosities and temperatures can be taken as constant over the short interval before $`t_{BH}`$. The event rate for $`t>t_{BH}`$ is : $$\frac{dN}{dt}(t)=C\left[\frac{L_{BH}}{10^{51}\mathrm{erg}/\mathrm{s}}\right]_0^{E_{max}}𝑑Ef(E)\left[\frac{\sigma (E)}{10^{42}\mathrm{cm}^2}\right],$$ (2) where $`f(E)`$ is the neutrino energy spectrum and $`\sigma (E)`$ the cross section. The upper limit $`E_{max}`$ on the integral allows only delays as large as $`tt_{BH}`$, i.e., $$E_{max}=m\sqrt{\frac{0.515D}{tt_{BH}}},$$ (3) where the units are as in Eq. (1). Note that the time and neutrino mass dependence appear only through $`E_{max}`$. For $`t<t_{BH}`$, $`E_{max}\mathrm{}`$, and the rate is constant. If the neutrino energy can be measured, as for some charged-current reactions, then the event rates for different neutrino energies can easily be obtained. For an H<sub>2</sub>O detector, the constant $`C`$ is $$C_{\mathrm{H2O}}=(1.74/\mathrm{s})\left[\frac{M_D}{1\mathrm{kton}}\right]\left[\frac{10\mathrm{kpc}}{D}\right]^2\left[\frac{1\mathrm{MeV}}{E}\right].$$ (4) For a Fermi-Dirac spectrum, $`E=3.15T`$. The constant for a <sup>208</sup>Pb detector can be obtained by scaling by the relative number of targets/kton, i.e., 18/208. The expected number of delayed counts after $`t_{BH}`$ can be calculated using Eq. (2). This will be useful when $`t_{BH}`$ can be measured independently. It can be shown that this has the very simple form: $$N_{del}=\frac{dN}{dt}(t_{BH})\times 0.515\left(\frac{m}{E_c}\right)^2D,$$ (5) where the event rate is in s<sup>-1</sup>, and the other units are as in Eq. (1). This formula would obviously be true if only a single energy contributed and the sharp cutoff in the event rate were rigidly translated by the delay. But it is remarkable and very convenient that it is still true even when there is a spectrum of energies and the event rate develops a decaying tail past the cutoff (as in Figs. 1 and 2). The physical significance of the “central” energy $`E_c`$ is that it is (to an excellent approximation) simply the Gamow peak of the falling thermal spectrum and the rising cross section. As derived, this is an exact result. Electron Neutrino Mass: We first consider the measurement of $`t_{BH}`$ and $`m_{\nu _e}`$ using the $`\overline{\nu }_e+pe^++n`$ events in the 32-kton SK detector. For $`T=5`$ MeV, the thermally-averaged cross section (for the sum of the two protons in H<sub>2</sub>O) is $`44.\times 10^{42}`$ cm<sup>2</sup> . The event rate at or before $`t_{BH}`$ is thus $`1500`$ s<sup>-1</sup>. After $`t_{BH}`$, the rate is zero if $`m_{\nu _e}=0`$ and will develop a tail if $`m_{\nu _e}>0`$. For a sharp edge, the edge position can be determined with an error given by the reciprocal of the event rate before the edge, i.e., the event spacing . If we knew that $`m_{\nu _e}=0`$, then $`t_{BH}`$ would be determined to $`1`$ ms. More realistically, a mass as large as the laboratory bound, $`m_{\nu _e}3`$ eV , would cause delays as large as 40 ms, so that the extracted $`t_{BH}`$ would be too large. However, we can simultaneously measure $`m_{\nu _e}`$ and $`t_{BH}`$ by splitting the $`\overline{\nu }_e+pe^++n`$ data into different ranges of neutrino energy (using $`E_\nu E_e+1.3`$ MeV). These are defined in the caption of Fig. 1. The Low group must be excluded from consideration because these events have positron total energy less than 10 MeV, and can be confused with the 5 – 10 MeV gammas from neutral-current reactions on <sup>16</sup>. The High group has very little delay and will thus primarily be sensitive to $`t_{BH}`$. Then the Mid group will determine $`m_{\nu _e}`$, by counting events delayed past the $`t_{BH}`$ determined by the High group. In Fig. 1, we show a possible analysis for the case of $`m_{\nu _e}=1.8`$ eV. In the High group, the number of events in the tail is $`1`$, so the cutoff appears sharp and is specified to within $`2`$ ms. This uncertainty affects the expected number in the Mid group by $`2`$ events. Even so, one can still reliably see a few delayed counts after the measured $`t_{BH}`$, enough to establish a nonzero mass (the statistics are discussed in more detail below). A more sophisticated fit would improve our results somewhat, and we assume a final uncertainty on $`t_{BH}`$ of about 1 ms. For a supernova in which the neutrino fluxes are not truncated by black hole formation, SK could detect an electron neutrino mass as small as $`3`$ eV . Mu and Tau Neutrino Masses: We consider mu and tau neutrino detection in OMNIS, a proposed supernova neutrino detector based on lead and iron . Since their energies are below the charged-current thresholds, supernova mu and tau neutrinos can be detected only by their neutral-current interactions. On the other hand, due to the temperature hierarchy, they will dominate the neutral-current yields. In OMNIS, the dominant neutral-current reaction is the spallation of single neutrons from lead. The neutrons could be detected by capture in a gadolinium-doped liquid scintillator, which yields an 8-MeV gamma cascade in about 0.030 ms (much smaller than typical mass delays). For $`T=8`$ MeV, the thermally-averaged cross section for the sum of $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ (or $`\nu _\tau `$ and $`\overline{\nu }_\tau `$) on <sup>208</sup>Pb, including the 1-neutron spallation probability, is $`760\times 10^{42}`$ cm<sup>2</sup> . The cross sections on <sup>206</sup>Pb and <sup>207</sup>Pb, which together comprise 46% of natural lead, are expected to be similar . For a supernova at 10 kpc in which the neutrino fluxes are not cut off by black hole formation, we assume that OMNIS will have $`1000`$ 1-neutron neutral-current events due to $`\nu _\mu `$, $`\nu _\tau `$, $`\overline{\nu }_\mu `$, and $`\overline{\nu }_\tau `$ on lead (the events on iron are not included in our calculations). This goal could be met with a 2.2 kton lead detector with perfect neutron detection efficiency. A realistic design based on 4 kton of lead and 10 kton of iron, and with about this many events, is described by Boyd . In Fig. 2, we plot the relevant neutral-current rate for different values of the neutrino mass, calculated using Eq. (2). In Fig. 3, we plot the number of delayed events $`N_{del}`$ as a function of the neutrino mass, using Eq. (5) and by direct integration. Equation (5) is remarkable for its simplicity, and also because it is written in terms of measurable quantities. The cutoff time $`t_{BH}`$ will be measured in SK. The neutral-current event rate at or before $`t_{BH}`$ will be measured in OMNIS, as will $`N_{del}`$. The central energy $`E_c`$ depends on the mu and tau neutrino temperature, which can be estimated by the neutral-current yields on different targets . We assume that the distance $`D`$ can be determined by astronomical means. Given the measured value of $`N_{del}`$, Eq. (5) can be immediately solved for the best-fit neutrino mass. If $`N_{del}=0`$ is measured, then the best-fit mass is $`m=0`$, and an upper limit can be placed. An expectation of 2.3 counts fluctuates down to 0 counts only 10% of the time. Thus, setting $`N_{del}=2.3`$, an upper limit on the mass $`m_{lim}`$ is obtained. This is the largest mass, given the expected Poisson statistics, that could be confused with the massless case. For the present case, this is 6.1 eV. Since the fractional error on $`N_{del}`$ due to Poisson statistics is large ($`1/\sqrt{2.3}65\%`$), errors on other inputs are expected to be irrelevant. The uncertainty on $`t_{BH}`$ from SK is assumed to be about 1 ms. From Fig. (2), this uncertainty can be seen to change the expected number $`N_{del}`$ by $`0.2`$ events, which is negligible. Other possible errors, e.g., the detector background, the disregarded 0.5 ms tail of the luminosity, and $`\nu _e`$ and $`\overline{\nu }_e`$ events after $`t_{BH}`$, are even less important . For a supernova that does not have the sharp cutoff in the rate characteristic of black hole formation, the model-independent $`t`$ analysis yields an $`m_{lim}`$ that is independent of the distance $`D`$ and scales as $`1/M_D^{1/4}`$ . For the present case, $`m_{lim}`$ scales as: $$m_{lim}E_c\sqrt{\frac{ED}{\sigma _{eff}L_{BH}M_D}},$$ (6) where $`\sigma _{eff}`$ is the thermally-averaged cross section. In terms of absolute sensitivity, these techniques compare as 21 eV and 6 eV, respectively. These differences are consequences of the sharp cutoff in the neutrino flux. Conclusions: If a black hole forms early in a core-collapse supernova, then the fluxes of the various flavors of neutrinos will be abruptly and simultaneously terminated when the neutrinospheres are enveloped by the event horizon. For a massive neutrino, the cutoff in the arrival time will be delayed by $`\mathrm{\Delta }t(m/E)^2`$ relative to a massless neutrino. The Galactic core-collapse supernova rate is about 3/century or higher , and the work of Brown and Bethe suggests that black holes are formed about half of the time. In the work of Burrows and Mezzacappa and Bruenn , the neutrino luminosities just before black hole formation are very high. These results indicate that there is a reasonably good chance that such an event could be observed by the present and proposed supernova neutrino detectors . If so, there are important practical consequences. First, since SK can measure the neutrino energy of the $`\overline{\nu }_e+pe^++n`$ events, both $`t_{BH}`$ and $`m_{\nu _e}`$ can be measured by the arrival times for different neutrino energies. An electron neutrino mass as small as 1.8 eV can be detected. Second, although the mu and tau neutrino energies are not measured in their neutral-current detection reactions, their masses can be measured by counting the number of events after $`t_{BH}`$. In the proposed OMNIS detector, a mu and tau neutrino mass (assumed degenerate ) as small as 6 eV can be detected. This is the only known direct technique with eV-scale sensitivity for these masses. Third, these results scale with the distance, luminosity, and detector mass as $`\sqrt{D/L_{BH}M_D}`$. This favorable scaling with the detector mass suggests that it would be realistic to consider even larger detectors, in order to reach 1 or 2 eV for all three neutrino masses. J.F.B. was supported by a Sherman Fairchild Fellowship from Caltech. R.N.B. was supported by NSF grant PHY-9901241. A.M. is supported at the Oak Ridge National Laboratory, managed by UT-Battelle, LLC, for the U.S. Dept. of Energy under contract DE-AC05-00OR22725. We thank Felix Boehm, Steve Bruenn, Will Farr, Josh Grindlay, Manoj Kaplinghat, Gail McLaughlin, Alex Murphy, Yong-Zhong Qian, Petr Vogel, and Jerry Wasserburg for discussions.
warning/0006/cond-mat0006200.html	ar5iv	text	# Density Functional Theory of Magnetic Systems Revisited ## Abstract The Hohenberg-Kohn theorem of density functional theory (DFT) for the case of electrons interacting with an external magnetic field (that couples to spin only) is examined in more detail than previously. An unexpected generalization is obtained: in certain cases (which include half metallic ferromagnets and magnetic insulators) the ground state, and hence the spin density matrix, is invariant for some non-zero range of a shift in uniform magnetic field. In such cases the ground state energy is not a functional of the spin density matrix alone. The energy gap in an insulator or a half metal is shown to be a ground state property of the N-electron system in magnetic DFT. The half metallic state of a ferromagnet has been receiving greatly increasing attention since its prediction from band theory to be the ground state of important magnetic materials such as CrO<sub>2</sub>, NiMnSb, and Sr<sub>2</sub>FeMoO<sub>6</sub> and several other intermetallics and oxides, and its unusual physical properties. Such systems have become very attractive for magnetoelectronic applications, where control of the spin degree of freedom is already leading to new devices. Materials thought to be half metals have been connected with the phenomena of colossal magnetoresistance (CMR), large tunneling MR, and large, low-field intergrain MR, and they would be optimal for applications of spin valve systems for non-volatile magnetic memory and for high density magnetic storage. The half metallic state is, in a one-electron picture, a collinear magnetic state in which one spin direction is metallic while the other is gapped (‘insulating’). This state is half metallic in another sense: the absence of low energy spin flips leads to a vanishing magnetic susceptibility like an insulator, but its charge response (conductivity) is that of a metal. These properties combine to give a one-electron description of a spin-charge separated state. In fact, almost all understanding of half metals so far is based on the one-electron picture, which opens up questions such as (1) what is a half metal in many body context, and (2) are there other unusual possibilities in magnetic systems? One general characterization might be in terms of conductivity (charge response) and susceptibility (spin response) alluded to above: in an insulator both vanish, in a conventional (even ferromagnetic) metal both are non-zero, and in a half metal the conductivity is non-zero while the susceptibility vanishes. A clear many body formulation is however lacking. Since density functional theory (DFT) is a rigorous many body theory for (chosen) ground state properties, we revisit the foundations of DFT with magnetic properties in mind. The first Hohenberg-Kohn (HK) theorem, which is the basis for the DFT of spin-independent particles, demonstrates the existence of a unique map $$n(𝐫)v(𝐫)\mathrm{mod}(\mathrm{constant}),$$ (1) where $`v`$ is the external potential and $`n`$ is the ground state particle density. According to the second HK theorem, the ground state energy and density are obtained as the solution to a variational principle: $$E[v\mu ]=min_n\{F[n]+n(v\mu )d^3r\},$$ (2) with $`\mu `$ the chemical potential. Although the variational principle has been put on an independent, more general basis, the uniqueness of the map (1) remains an important issue regarding the existence and uniqueness of the functional derivative $`\delta F/\delta n=(v\mu )`$. DFT, as extended by Kohn and Sham and many others, forms the basis of our understanding of the electronic behavior of real materials. The theory has been extended to electrons with spin and also applied heavily, however the HK theorem for interacting particles with spin has repeatedly been stated to be analagous to the HK theorem, although this was already questioned in . Zero susceptibility, however, would imply that the ground state spin density does not change when an external magnetic field is changed. In this paper we construct a more revealing generalization of the HK theorem, obtain explicitly the conditions that allow half metallicity, and demonstrate some unexpected consequences. We consider the system in an external magnetic field $`𝐁(𝐫)`$ in the (commonly considered) non-relativistic limit, in which the field acts only on the electron spin and the dipolar interaction between spins is neglected. The potentials can be combined into 2$`\times `$2 spin matrix $$u_{ss^{}}(𝐫)=v(𝐫)\delta _{ss^{}}\mu _B𝐁(𝐫)\stackrel{}{\sigma }_{ss^{}}.$$ (3) The external field $`𝐁`$ may vary in magnitude and direction. Realizing that two different scalar potentials cannot lead to the same ground state $`\mathrm{\Psi }`$, the original derivation of HK concluded that if $`\mathrm{\Psi }v\mu `$ is unique then $`nv\mu `$ is unique. Following the original derivation of HK, we begin by supposing that there are two different potentials $`u,u^{}`$ which lead to the same ground state $`\mathrm{\Psi }`$. We show that $`\mathrm{\Psi }(v\mu ,\stackrel{}{B})`$ is not a unique mapping in general. The many-body Hamiltonian of the system is $$\widehat{H}=\widehat{T}+\widehat{W}+\widehat{U},$$ (4) where $`\widehat{T}`$ is the kinetic energy operator, $`\widehat{W}`$ is the Coulomb interaction energy, and $`\widehat{U}`$ is the interaction with the external potential. The fermionic many-particle Schrödinger equation is (in atomic units) $`[{\displaystyle \underset{i}{\overset{N}{}}}{\displaystyle \frac{_i^2}{2}}+{\displaystyle \underset{i<j}{\overset{N}{}}}w(𝐫_i,𝐫_j)]\mathrm{\Psi }(𝐫_1\alpha _1,..,𝐫_N\alpha _N)`$ (5) $`+{\displaystyle \underset{i}{\overset{N}{}}}{\displaystyle \underset{\beta _i}{}}u_{\alpha _i,\beta _i}(𝐫_i)\mathrm{\Psi }(𝐫_1\alpha _1,..,𝐫_i\beta _i,..,𝐫_N\alpha _N)`$ (6) $`=E\mathrm{\Psi }(𝐫_1\alpha _1,\mathrm{},𝐫_N\alpha _N),`$ (7) where $`𝐫_i,\alpha _i`$ are the space and spin coordinates of the $`i`$-th electron; $`w(𝐫,𝐫^{})=e^2/\|𝐫𝐫^{}\|`$ is the Coulomb interaction. Assume there are two external potentials $`u,u^{}`$ with energies $`E,E^{}`$ that have the same ground state wave function $`\mathrm{\Psi }(𝐫_1\alpha _1,\mathrm{},𝐫_N\alpha _N)`$. Subtracting the two many-particle Schrödinger equations leads to $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{\beta _i}{}}\mathrm{\Delta }u_{\alpha _i,\beta _i}(𝐫_i)\mathrm{\Psi }(𝐫_1\alpha _1,..,𝐫_i\beta _i,..,𝐫_N\alpha _N)=`$ (8) $`\mathrm{\Delta }E\mathrm{\Psi }(𝐫_1\alpha _1,\mathrm{},𝐫_N\alpha _N),`$ (9) where $`\mathrm{\Delta }u=uu^{},\mathrm{\Delta }E=EE^{}`$. Now, we perform a unitary spin rotation $`Q_{ss^{}}(𝐫)`$ at each point of space that diagonalizes the difference in potentials (i.e. rotates $`\stackrel{}{B}`$ to lie along the $`\widehat{z}`$ direction: $$\{Q(𝐫)[\mathrm{\Delta }u(𝐫)]Q^{}(𝐫)\}_{ss^{}}=\mathrm{\Delta }\stackrel{~}{u}_s(𝐫)\delta _{ss^{}}.$$ (10) The wavefunction is transformed according to $`{\displaystyle \underset{i}{\overset{N}{}}}Q_{\alpha _i\alpha _i^{}}(𝐫_i)\mathrm{\Psi }(𝐫_1\alpha _1^{},\mathrm{},𝐫_N\alpha _N^{})\stackrel{~}{\mathrm{\Psi }}(𝐫_1\alpha _1,\mathrm{},𝐫_N\alpha _N),`$ (11) where $`_i^NQ_{\alpha _i\alpha _i^{}}(𝐫_i)`$ is the operator that rotates each of the ($`\alpha _i`$). Collecting these results gives $`{\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{\Delta }\stackrel{~}{u}_{\alpha _i}(𝐫_i)\stackrel{~}{\mathrm{\Psi }}(𝐫_1\alpha _1,\mathrm{},𝐫_N\alpha _N)=\mathrm{\Delta }E\stackrel{~}{\mathrm{\Psi }}(𝐫_1\alpha _1,\mathrm{},𝐫_N\alpha _N).`$ (12) $`\stackrel{~}{\mathrm{\Psi }}`$ is some $`\{𝐫_i\}`$-dependent multi-component function of the $`2^N`$ possible spin configurations, at least one of which must be non-zero. Choose a non-zero component $`\stackrel{~}{\mathrm{\Psi }}_c`$ and denote by $`N_{}`$ the number of $`\alpha _i=`$ values in this component. Since $`\stackrel{~}{\mathrm{\Psi }}`$ is antisymmetric (as was $`\mathrm{\Psi }`$) with respect to permutations of ($`𝐫_i\alpha _i`$) with ($`𝐫_j\alpha _j`$), we may renumber the particle indices in such a way that $`\alpha _1=\alpha _2=\mathrm{}=\alpha _N_{}`$, $`\alpha _{N_{}+1}=\alpha _{N_{}+2}=\mathrm{}=\alpha _N.`$ This ordering lets us write $`\{{\displaystyle \underset{i=1}{\overset{N_{}}{}}}\mathrm{\Delta }\stackrel{~}{u}_{}(𝐫_i)+{\displaystyle \underset{i=N_{}+1}{\overset{N}{}}}\mathrm{\Delta }\stackrel{~}{u}_{}(𝐫_i)\}\stackrel{~}{\mathrm{\Psi }_c}(𝐫_1,\mathrm{},𝐫_N)`$ (13) $`=\mathrm{\Delta }E\stackrel{~}{\mathrm{\Psi }_c}(𝐫_1,\mathrm{},𝐫_N).`$ (14) This equation must hold for all values of ($`𝐫_1,\mathrm{},𝐫_N)`$. (We suppose $`u,u^{}`$ are analytic in $`𝐫`$ except possibly at isolated points, so that $`\stackrel{~}{\mathrm{\Psi }_c}`$ is non-zero almost everywhere.) By varying only $`𝐫_1`$, and then separately varying only $`𝐫_N`$, we obtain $$\mathrm{\Delta }\stackrel{~}{u}_{}=C_{},\mathrm{\Delta }\stackrel{~}{u}_{}=C_{},$$ (15) where $`C_{},C_{}`$ are constants. The special cases $`N_{}=0`$ or $`N_{}=N`$ do not lead to new consequences. For further analysis, we consider separate cases. Case A: impure spin states. Suppose that there are at least two components of $`\stackrel{~}{\mathrm{\Psi }}`$ with different values of $`N_{}`$ and hence $`N_{}=NN_{}`$. Then $$N_{}C_{}+(NN_{})C_{}=\mathrm{\Delta }E.$$ (16) Since this holds for two different values of $`N_{}`$, it follows that $`C_{}=C_{}C,`$ which leads to $`\mathrm{\Delta }\stackrel{~}{u}_{}=\mathrm{\Delta }\stackrel{~}{u}_{}`$ so $`\mathrm{\Delta }u_{\alpha ,\beta }=C\delta _{\alpha ,\beta }`$. By the ground state energy minimum principle, this recovers the usual Hohenberg-Kohn result $`n_{ss^{}}=n_{ss^{}}^{}\left(\begin{array}{c}v(𝐫)v^{}(𝐫)C,\\ 𝐁(𝐫)𝐁^{}(𝐫)0,\end{array}\right)`$ (19) implying a non-zero susceptibility. Case B: pure spin states. Suppose now that all non-zero components of $`\stackrel{~}{\mathrm{\Psi }}`$ have the same value of $`N_{}`$ and $`N_{}`$. These may be considered as “pure spin” states, eigenfunctions of the operator $`\widehat{S}^z`$ = $`_i\sigma _{\alpha _i\beta _i}^z/2`$ with eigenvalues $`S_z=N_{}(N/2)=(N_{}N_{})/2`$. Then $`C_{}`$ and $`C_{}`$ need not be equal and we can write $`\mathrm{\Delta }\stackrel{~}{u}`$ $`=`$ $`\left(\begin{array}{cc}C_{}& 0\\ 0& C_{}\end{array}\right)=\overline{C}\mathrm{𝟏}\mu _B\overline{B}\sigma ^z,`$ (22) where $`\overline{C}=(C_{}+C_{}`$)/2 and $`\mu _B\overline{B}=(C_{}C_{}`$)/2. Backtransforming according to the inverse of Eq. (10) gives $$\mathrm{\Delta }u_{\alpha \beta }(𝐫)=\overline{C}\delta _{\alpha \beta }+\mu _B\overline{B}[Q^{}(𝐫)\sigma _zQ(𝐫)]_{\alpha \beta }.$$ (23) The last term on the right is position-dependent, non-diagonal, and non-vanishing in general. In this case the conditions for identical ground state wavefunctions are $$\mathrm{\Psi }=\mathrm{\Psi }^{}\left(\begin{array}{ccc}v(𝐫)v^{}(𝐫)& =& \overline{C}\\ 𝐁(𝐫)𝐁^{}(𝐫)& =& \overline{B}\widehat{e}(𝐫).\end{array}\right).$$ (24) where $`\widehat{e}`$ is the unit vector $`\frac{1}{2}Tr\{\stackrel{}{\sigma }Q^{}(𝐫)\sigma _zQ(𝐫)\}`$. The result (19) is modified accordingly. This result is a highly non-trivial generalization of the HK theorem: two magnetic fields whose difference is constant in magnitude, but possibly is non-unidirectional, may give rise to the same ground state. Now we investigate the conditions on $`Q`$ for which $`\stackrel{~}{\mathrm{\Psi }}`$ is an eigenstate of $`\widehat{S}^z`$, i.e. $`\stackrel{~}{\mathrm{\Psi }}`$ describes a collinear spin arrangement. Considering the Hamiltonian Eq. (4), there must be an operator $$\widehat{U}_o=\underset{i=1}{\overset{N}{}}Q_{\alpha _i^{}\alpha _i}^{}(𝐫_i)\sigma _{\alpha _i^{}\beta _i^{}}^zQ_{\beta _i^{}\beta _i}(𝐫_i)$$ (25) that commutes with $`\widehat{T}+\widehat{U}+\widehat{W}`$. We now specialize to the particular case where one of the external fields, $`𝐁^{}`$, is zero. Since the interaction $`\widehat{W}`$ is spin independent, $`\widehat{U}_o`$ will commute with $`H^{}`$ if and only if it commutes with $`\widehat{T}`$. One can show that this necessitates that $`Q`$ be $`𝐫`$-independent, so that $`\mathrm{\Psi }`$ itself is an eigenstate of $`\widehat{S}^z`$, and hence is a collinear spin state. The second condition in Eq. (24) reduces to $`𝐁𝐁^{}=B\widehat{z}`$: a turning on of a uniform magnetic field leaves the ground state invariant. Restated: in the subspace of collinear magnetizations, the ground state determines the magnetic field only up to some codirectional uniform field. A direct corollary is that there is no longer any ground state energy functional $`E[n]`$ of the density $`n_{ss^{}}(𝐫)`$ alone (see Eq. (26) below). In the collinear situation, inserting $`\mathrm{\Delta }u=\mu _BB\sigma _z`$ into Eq. (6) yields $$\mathrm{\Delta }E(B)=(N_{}N_{})\mu _BB,$$ (26) which gives the well known dependence of energy vs. field for a system of fixed spin. Consider as a simple example a Be atom in a uniform magnetic field, with its ground state characterized as $`1s^22s^2`$ ($`N=4`$, $`N_{}=2`$). The lowest excited $`\widehat{S}_z`$-eigenstate is $`1s^22s2p`$ with $`N_{}=3`$. Its excitation energy is that of a $`2s2p`$ promotion. There is another excited state $`1s2s^22p`$ with the same $`N_{}`$, but the much higher excitation energy of a $`1s2p`$ core excitation. The energetically lowest $`N_{}=4`$ state is $`1s2s2p^2`$ whose excitation energy is roughly the sum of the previous two. The situation is sketched in Fig. 1(a), where the lines with positive slopes correspond to states with all spins reversed. Since states with $`N_{}=N/2\pm n`$ are degenerate for $`B=0`$, Fig. 1(a) may be supplemented symmetrically to the vertical axis. Hence, for $`\|B\|<B_0`$ the groundstate is $`1s^22s^2`$ with energy $`E_0`$, for $`B_0<B<B_1`$ the ground state is $`1s^22s2p`$ with energy $`E_12\mu _BB`$, and for $`BB_1`$ the ground state is $`1s2s2p^2`$ with energy $`E_34\mu _BB`$. The ground state does not change with field except at certain isolated values. In an extended system, say a non-magnetic insulator with gap $`\mathrm{\Delta }_g`$, there is a continuum above $`\mathrm{\Delta }_g`$ (one excited electron with reversed spin), another continuum above $`2\mathrm{\Delta }_g`$ (two excited electrons) and so on, as illustrated in Fig. 1(b). In an extended system one would prefer to consider the intensive quantity $$\frac{\mathrm{\Delta }E(B)}{N}=\mu _BB\left(\frac{N_{}N_{}}{N}\right)$$ (27) instead of $`\mathrm{\Delta }E`$ itself. Then, one finds that for $`\mu _BB<\mathrm{\Delta }_g`$ the groundstate is independent of $`B`$, beyond which the state changes and $`\mathrm{\Delta }E/N`$ veers off. Thus while the gap $`\mathrm{\Delta }_g`$ is not a ground state property of the $`N`$ particle system in paramagnetic DFT (it involves the N$`\pm `$1 particle ground states), it is a ground state property in the presence of a uniform field. For a stoichiometric half metal with moment per cell $`\mu _B(`$ an integer) the picture is related, except there is an overall bias – a slope of -$`\mu _B`$ in the energy per cell – and the positive and negative $`B`$ directions are not symmetric. The situation that is sketched in Fig. 1(c) has a gap $`\mathrm{\Delta }_v+\mathrm{\Delta }_c`$ for $``$ spin states, with no gap for $``$ spin. The chemical potential $`\mu `$ corresponds to the energy to remove an $``$ spin, and the quantities $`\mathrm{\Delta }_v=\mu _BB_v`$, $`\mathrm{\Delta }_c=\mu _B\|B_c\|`$ represent the energy, or field, required to flip a spin from $``$ to $``$, or vice versa. Note again that the interval of $`B`$ for which the state does not change, which is the gap in the $``$ spectrum, is a ground state property of the $`N`$ particle system in an external magnetic field. It is useful to consider the form of constrained DFT in which $`N_{}`$ and $`N_{}`$ are specified, which leads to two associated chemical potentials $`\mu _{}`$, $`\mu _{}`$. Then as $`N_s`$ is changed to $`N_s\pm 1`$, $`\mu _s`$ may vary only to order 1/$`N_s`$ (metallic behavior) or it may jump discontinuously across a gap, just as is the case for insulators. The half metal is defined as that situation in which one and only one of $`\mu _s`$ (we have choosen $``$) is discontinuous upon addition of one electron. For an insulator, there is a discontinuity in $`\mu `$ for both spins. We now consider the KS eigenvalue spectrum. As long as the external field shifts the bands sufficiently little not to disturb the half metallicity ($`B_c<B<B_v`$), the ground state, and hence the charge density in each spin channel, remain unchanged. Using the same arguments as were applied to establish the discontinuity in $`v_{xc}(N(\mu ))`$ for an insulator as $`\mu `$ crosses the gap (the kinetic energy is discontinuous across the gap), one finds that there is a discontinuity in $`v_{xc,}(N_{},N_{})`$ if the filling with $`N_{}`$ moves $`\mu _{}`$ across the gap. The Kohn-Sham gap $`\epsilon _g`$ is smaller than the true (quasiparticle) gap $`\mathrm{\Delta }_g=\mathrm{\Delta }_c+\mathrm{\Delta }_v`$. When the magnetic field is large enough that $`\mu `$ reaches the KS band minimum $`\epsilon _c(N_{})`$, the occupation of that channel becomes N$`{}_{}{}^{}+ϵ`$ (with $`ϵ0`$). This is the point of the discontinuity, where the KS conduction eigenvalue (in fact, the entire $``$ spectrum) jumps upward. By comparison with Dyson’s equation, and the fact that the system’s ground state spin densities must be the same whether obtained from DFT or the quasiparticle Greens function, this jump must be such as to make $`\epsilon _c(N_{},N_{}+ϵ)`$ $`\mathrm{\Delta }_c`$, the quasiparticle conduction band edge, for $`ϵ`$0. It is apparent then that the KS gap in the insulating channel is not equal to the true gap in that channel, and that $`\epsilon _c(N_{},N_{})\mu `$ is not the true spin flip energy (which is $`\mathrm{\Delta }_c`$ \- $`\mu `$). By our definition, as the reverse field is applied and $`\mu `$ is driven toward the valence band maximum $`\epsilon _v`$, there is no discontinuity, and the other spin flip energy – a true excitation energy – is given correctly by DFT. Needless to say, an approximation such as the local density approximation that interpolates across the discontinuity, will fail to predict both $`\mathrm{\Delta }_c`$ and $`\mathrm{\Delta }_v`$. We now summarize. We have presented new, rigorous results for the Hohenberg-Kohn mapping in a magnetic field. We obtain conditions that characterize half metals: (1) two collinear systems in different uniform magnetic fields may have the same half metallic (or magnetic insulating) ground state; (2) exactly one of the chemical potentials $`\mu _s`$ is discontinuous upon particle addition to a half metal. We have pointed out other consequences, primary among them being that the ground state energy of a system is no longer a unique functional of the density $`n_{ss^{}}`$ when magnetic fields are allowed (although the ground state itself is), and that the gap in a half metal is a ground state property of the $`N`$ particle system. These results are only exact in the non-relativistic ($`c\mathrm{}`$) limit. For $`c`$ finite, half metallicity is an approximate notion due to orbital currents and orbital moments and spin-orbit coupling that mixes them, and the general theory probably restores the conventional theorems of DFT. Still, the notion of half metallicity will be an important model limit. We acknowledge stimulating discussions with I. I. Mazin throughout the course of this work, and helpful interaction with W. Kohn during the early stages. This study was begun when the authors were at the Institute of Theoretical Physics at the University of California at Santa Barbara, which is supported by National Science Foundation Grant PHY-9407194. W.E.P. was supported by National Science Foundation Grant DMR-9802076.
warning/0006/gr-qc0006002.html	ar5iv	text	# Rotating Relativistic Thin Disks ## 1 Introduction Axially symmetric solutions of Einstein field equations corresponding to disk like configurations of matter are of great astrophysical interest and have been extensively studied. These solutions can be static or stationary and with or without radial pressure. Solutions for static disks without radial pressure were first studied by Bonnor and Sackfield , and Morgan and Morgan , and with radial pressure by Morgan and Morgan , and, in connection with gravitational collapse, by Chamorro, Gregory and Stewart . Disks with radial tension has been recently considered in . Also models of disks with, electric fields , magnetic fields , and both magnetic and electric fields have been introduced recently . Using a solitonic technique Neugebauer and Meinel found solutions representig rigidly rotating disks of dust . Several classes of exact solutions of the Einstein field equations that represent disks were obtained by different authors -. The stability of static models with no radial pressure can be explained by either assuming the existence of hoop stresses or that the particles on the disk plane move under the action of their own gravitational field in such a way that as many particles move clockwise as counterclockwise. This last interpretation is frequently made since it can be invoked to mimic true rotational effects. Even though, this interpretation can be seen as a device, there are observational evidence of disks made of streams of rotating and counter-rotating matter . For disks without tension, i.e., usual fluid disks, the rotation is a necessary ingredient to have stability. Despite its relevance the literature on exact relavistic disks solutions to the Einstein field equations is scarce, we only found a discussion of a rotating disk as a source of Kerr metric . The aim of this work is to study models of rotating relativistic disks in some detail. The models are based in Taub-NUT and Kerr metrics, that are some of the simplest axially symmetric stationary solutions of the vacuum Einstein equations . To construct the disk solutions we use the well-known “displace, cut and reflect” method. We find that for disks built from a generic stationary axially symmetric metric the “sound velocity”, $`(pressure/density)^{1/2}`$, is equal to the geometric mean of the prograde and retrograde geodesic circular velocities of a test particle moving on the disk. We also found that the generic disk may have zones with heat flow. For the two families of models studied the boundaries that separates the zones with and without heat flow are not stable against radial perturbations (ring formation). In Sec. II we study the generation of a disk model using the “displace, cut and reflect” method from a vacuum solution of the Einstein field equations for a stationary axially symmetric metric, in particular, we study: the general form of the associated energy-momentum tensor, disk velocity, and angular momentum. In the next section, Sec. III, we discuss the relation between the geodesic circular velocities on the disk and the pressure and density. In Sec. IV a class of models of rotating disks based on the Taub-NUT metric is presented. In particular, we study the apparition of heat flow and stability against radial perturbations. In Sec. V, another class of models of rotating disks is presented, this time based in the Kerr metric. Also the stability against radial perturbation is considered. The material presented here is complementary to the work of reference . In the last section, Sec VI , we summarize our main results and made some comments. ## 2 Rotating Relativistic Disks In this section we present a summary of the main quantities associated to the disk, we follow closely our reference . A sufficiently general metric for our purposes is Weyl-Lewis-Papapetrou line element, $$ds^2=e^{2\mathrm{\Phi }}(dt+𝒜d\phi )^2+e^{2\mathrm{\Phi }}[r^2d\phi ^2+e^{2\mathrm{\Lambda }}(dr^2+dz^2)],$$ (1) where $`\mathrm{\Phi }`$, $`\mathrm{\Lambda }`$ and $`𝒜`$ are functions of $`r`$ and $`z`$ only. Assuming the existence of the second derivatives of the functions $`\mathrm{\Phi }`$, $`\mathrm{\Lambda }`$ and $`𝒜`$, the Einstein vacuum equations for this metric are equivalent to, $`\mathrm{\Phi }_{,rr}+{\displaystyle \frac{\mathrm{\Phi }_{,r}}{r}}+\mathrm{\Phi }_{,zz}+{\displaystyle \frac{e^{4\mathrm{\Phi }}}{2r^2}}(𝒜_{,r}^2+𝒜_{,z}^2)`$ $`=`$ $`0,`$ (2a) $`𝒜_{,rr}{\displaystyle \frac{𝒜_{,r}}{r}}+𝒜_{,zz}+4(\mathrm{\Phi }_{,r}𝒜_{,r}+\mathrm{\Phi }_{,z}𝒜_{,z})`$ $`=`$ $`0,`$ (2b) $`\mathrm{\Lambda }_{,r}`$ $`=`$ $`r(\mathrm{\Phi }_{,r}^2\mathrm{\Phi }_{,z}^2){\displaystyle \frac{e^{4\mathrm{\Phi }}}{4r}}(𝒜_{,r}^2𝒜_{,z}^2),`$ (2c) $`\mathrm{\Lambda }_{,z}`$ $`=`$ $`2r\mathrm{\Phi }_{,r}\mathrm{\Phi }_{,z}{\displaystyle \frac{e^{4\mathrm{\Phi }}}{2r}}𝒜_{,r}𝒜_{,z}.`$ (2d) Now if the first derivatives of the metric tensor are not continuous on the plane $`z=0`$ with discontinuity functions, $$b_{ab}=g_{ab,z}\|_{_{z=0^+}}g_{ab,z}\|_{_{z=0^{}}},$$ the Einstein equations yield an energy-momentum tensor (EMT) $`T_{ab}=Q_{ab}\delta (z)`$, where $`\delta (z)`$ is the usual Dirac function with support on the disk and $$Q_b^a=\frac{1}{2}\{b^{az}\delta _b^zb^{zz}\delta _b^a+g^{az}b_b^zg^{zz}b_b^a+b_c^c(g^{zz}\delta _b^ag^{az}\delta _b^z)\}.$$ is the distributional energy-momentum tensor. The “true” surface energy-momentum tensor of the disk can be written as $`S_{ab}=e^{\mathrm{\Lambda }\mathrm{\Phi }}Q_{ab}`$. For the metric (1) we obtain $`S_{00}`$ $`=2e^{3\mathrm{\Phi }\mathrm{\Lambda }}\left[2\mathrm{\Phi }_{,z}\mathrm{\Lambda }_{,z}\right],`$ (3a) $`S_{01}`$ $`=e^{3\mathrm{\Phi }\mathrm{\Lambda }}\left[4𝒜\mathrm{\Phi }_{,z}2𝒜\mathrm{\Lambda }_{,z}+𝒜_{,z}\right]`$ (3b) $`S_{11}`$ $`=2e^{3\mathrm{\Phi }\mathrm{\Lambda }}\left[(r^2e^{4\mathrm{\Phi }}𝒜^2)\mathrm{\Lambda }_{,z}+2𝒜^2\mathrm{\Phi }_{,z}+𝒜𝒜_{,z}\right]`$ (3c) where all the quantities are evaluated at $`z=0^+`$. The eigenvalue problem for the energy-momentum tensor (3a) - (3c), $$S_b^a\xi ^b=\lambda \xi ^a,$$ (4) has the solutions $$\lambda _\pm =\frac{1}{2}\left(T\pm \sqrt{D}\right),$$ (5) and $`\lambda _r=\lambda _z=0`$, where $$T=S_0^0+S_1^1,D=(S_1^1S_0^0)^2+4S_1^0S_0^1.$$ (6) We can write $$g_{ab}=V_aV_b+W_aW_b+X_aX_b+Y_aY_b,$$ (7) and the canonical form of the EMT, $$S_{ab}=\sigma V_aV_b+PW_aW_b+K\left(V_aW_b+W_aV_b\right),$$ (8) with a orthonormal basis $`V^a`$ $`=`$ $`N_0(1,\mathrm{\Omega },0,0),`$ (9a) $`W^a`$ $`=`$ $`N_1(\mathrm{\Delta },1,0,0),`$ (9b) $`X^a`$ $`=`$ $`e^{\mathrm{\Phi }\mathrm{\Lambda }}(0,0,1,0),`$ (9c) $`Y^a`$ $`=`$ $`e^{\mathrm{\Phi }\mathrm{\Lambda }}(0,0,0,1),`$ (9d) where $`N_0`$ and $`N_1`$ are normalization factors, and $$\mathrm{\Omega }=\{\begin{array}{ccc}(\lambda _{}S_0^0)/S_1^0& ,& D0,\\ & & \\ (S_1^1S_0^0)/2S_1^0& ,& D0,\end{array}$$ (10) $$\mathrm{\Delta }=\{\begin{array}{ccc}(\lambda _+S_1^1)/S_0^1& ,& D0,\\ & & \\ 0& ,& D0,\end{array}$$ (11) The energy density, the azimuthal pressure, and the heat flow function are, respectively, $$\sigma =\{\begin{array}{ccc}\lambda _{}& ,& D0,\\ & & \\ T/2& ,& D0,\end{array}$$ (12) $$P=\{\begin{array}{ccc}\lambda _+& ,& D0,\\ & & \\ T/2& ,& D0,\end{array}$$ (13) $$K=\{\begin{array}{ccc}0& ,& D0,\\ & & \\ \sqrt{D}/2& ,& D0.\end{array}$$ (14) The orthonormal basis $`\{V^a,W^a,X^a,Y^a\}`$ is comoving with the disk. Thus the time-like vector $`V^a`$ will define the velocity vector of the disk: $$V^a=(V^0,V^1,0,0)=V^0(1,\mathrm{\Omega },0,0),$$ (15) where $$V^0=\frac{e^\mathrm{\Phi }}{\sqrt{1V^2}},$$ (16) and $$V=\sqrt{\frac{g_{11}\mathrm{\Omega }^2+2g_{01}\mathrm{\Omega }}{g_{00}}}$$ (17) is the tangential velocity of the disk. The specific angular momentum of a particle of the disk, with mass $`\mu `$, is given by $$h=\frac{p_\phi }{\mu }=g_{\phi a}V^a.$$ (18) Thus, $$h=\frac{g_{11}\mathrm{\Omega }+g_{01}}{\sqrt{g_{00}2g_{01}\mathrm{\Omega }g_{11}\mathrm{\Omega }^2}}.$$ (19) A condition of stability under radial perturbations is $$\frac{d(h^2)}{dr}=2h\frac{dh}{dr}>0,$$ (20) For $`h>0`$, we have stability when the specific angular momentum is an increasing function of $`r`$. This criteria is an extension of Rayleigh criteria of stability of a fluid at rest in a gravitational field, see for instance . ## 3 The counter-rotating Model In this section we analyze the model of disks made of particles moving in opposite directions. When the discriminant $`D`$ is positive there is not heat flow. Furthermore, when $`P>0`$ an observer comoving with the disks can consider the energy-momentum tensor as representing two streams of collisionless particles that circulate in opposite directions. Let be $`u^a=(u^0,u^1,0,0)=u^0(1,\omega ,0,0)`$ the velocity vector of the stream. The angular velocity $`\omega `$ can be obtained from the geodesic equation for a test particle, we get $$g_{11,r}\omega ^2+2g_{01,r}\omega +g_{00,r}=0,$$ (21) with solutions, $$\omega _\pm =\frac{g_{01,r}\pm \sqrt{g_{01,r}^2g_{00,r}g_{11,r}}}{g_{11,r}}.$$ (22) Therefore, in general, the two streams circulate with different velocities. We can compute the tangential velocity of the streams by projecting the velocity vector $`u^a`$ onto the comoving tetrad, $`e_{\widehat{a}}^{}{}_{}{}^{b}`$ = $`\{V^b,W^b,X^b,Y^b\}`$, $$u^{\widehat{a}}=e_{}^{\widehat{a}}{}_{b}{}^{}u^b=\eta ^{\widehat{a}\widehat{c}}e_{\widehat{c}b}u^b.$$ (23) We get, $$U_\pm =\left\|\frac{u^{\widehat{1}}}{u^{\widehat{0}}}\right\|=\left\|\frac{W_0+W_1\omega _\pm }{V_0+V_1\omega _\pm }\right\|.$$ (24) By performing the product of $`U_+`$ and $`U_{}`$ we obtain $$U_+U_{}=\left\|\frac{W_{0}^{}{}_{}{}^{2}+W_0W_1(\omega _++\omega _{})+W_{1}^{}{}_{}{}^{2}\omega _+\omega _{}}{V_{0}^{}{}_{}{}^{2}+V_0V_1(\omega _++\omega _{})+V_{1}^{}{}_{}{}^{2}\omega _+\omega _{}}\right\|.$$ (25) From (21) we have the relations, $$\omega _++\omega _{}=\frac{2g_{01,r}}{g_{11,r}},\omega _+\omega _{}=\frac{g_{00,r}}{g_{11,r}}.$$ (26) Thus, $$U_+U_{}=\left\|\frac{g_{11,r}W_{0}^{}{}_{}{}^{2}2g_{01,r}W_0W_1+g_{00,r}W_{1}^{}{}_{}{}^{2}}{g_{11,r}V_{0}^{}{}_{}{}^{2}2g_{01,r}V_0V_1+g_{00,r}V_{1}^{}{}_{}{}^{2}}\right\|.$$ (27) By using (7) and (8), we can write $$U_+U_{}=\left\|\frac{A+\sigma B}{APB}\right\|,$$ (28) where $`A`$ $`=`$ $`g_{11,r}S_{00}2g_{01,r}S_{01}+g_{00,r}S_{11},`$ (29a) $`B`$ $`=`$ $`g_{00}g_{11,r}2g_{01}g_{01,r}+g_{00,r}g_{11}.`$ (29b) Using the Einstein equations (2a) - (2d) and the expressions (3a) - (3c) for the energy-momentum tensor (with $`K=0`$) we can show that $$A=2r(S_0^0+S_1^1).$$ (30) Also, using the metric (1) we get, $$B=2r.$$ (31) Hence, assuming that $`P>0`$ and $`\sigma >0`$, we obtain $`U_+U_{}`$ $`=`$ $`\left\|{\displaystyle \frac{S_0^0+S_1^1+\sigma }{S_0^0+S_1^1P}}\right\|`$ (32) $`=`$ $`{\displaystyle \frac{P}{\sigma }}.`$ (33) In other words, for disks built from a generic stationary axially symmetric metrics the “sound velocity” $`(P/\sigma )^{1/2}`$ is equal to the geometric mean of the prograde and retrograde geodesic circular velocities of a test particle moving on the disk. ## 4 Taub-NUT Disks The simplest stationary axially symmetric solution of the Einstein equations is the Taub-NUT solution that can be written as (1) with $`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}\left[{\displaystyle \frac{x^21}{x^2+2px+1}}\right],`$ (34a) $`\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}\left[{\displaystyle \frac{x^21}{x^2y^2}}\right],`$ (34b) $`𝒜`$ $`=`$ $`2\alpha qy,`$ (34c) where $`p=m/\alpha `$ , $`q=l/\alpha `$ , with $`\alpha ^2=m^2+l^2`$, so that $`p^2+q^2=1`$. $`m`$ and $`l`$ are the mass and the NUT parameter, respectively. $`x`$ and $`y`$ are the prolate spheroidal coordinates, given by $`2\alpha x`$ $`=`$ $`\sqrt{r^2+(z+b+\alpha )^2}+\sqrt{r^2+(z+b\alpha )^2},`$ (35a) $`2\alpha y`$ $`=`$ $`\sqrt{r^2+(z+b+\alpha )^2}\sqrt{r^2+(z+b\alpha )^2},`$ (35b) with $`b>\alpha >0`$. Note that we have displaced the origin of the $`z`$-axis in $`b`$. From the above expressions we can compute the trace ($`T`$) and discriminant ($`D`$) of the energy-momentum tensor. We found $`\stackrel{~}{T}=\alpha T`$ and $`\stackrel{~}{D}=\alpha ^2D`$, with $`\stackrel{~}{T}`$ $`=`$ $`{\displaystyle \frac{4\overline{y}[\overline{y}^2(2\overline{x}^3+3p\overline{x}^2p)+\overline{x}(p\overline{x}^33p\overline{x}2)]}{[(\overline{x}^2\overline{y}^2)(\overline{x}^2+2p\overline{x}+1)]^{3/2}}}`$ (36a) $`\stackrel{~}{D}`$ $`=`$ $`{\displaystyle \frac{16[\overline{y}^2(p\overline{x}^2+2\overline{x}+p)^2q^2\overline{x}^2(\overline{x}^21)(1\overline{y}^2)]}{(\overline{x}^2\overline{y}^2)(\overline{x}^2+2p\overline{x}+1)^3}}.`$ (36b) $`\overline{x}`$ and $`\overline{y}`$ are given by $`2\overline{x}`$ $`=`$ $`\sqrt{\stackrel{~}{r}^2+(\kappa +1)^2}+\sqrt{\stackrel{~}{r}^2+(\kappa 1)^2},`$ (37a) $`2\overline{y}`$ $`=`$ $`\sqrt{\stackrel{~}{r}^2+(\kappa +1)^2}\sqrt{\stackrel{~}{r}^2+(\kappa 1)^2},`$ (37b) where $`\stackrel{~}{r}=r/\alpha `$ and $`\kappa =b/\alpha `$. Even though we have an exact solution and all the relevant quantities of the disk can be explicitly computed the resulting expressions are cumbersome and not very illuminating, the adoption of the graphic method is more adequate in this case. The first quantity to be considered is the discriminant D that will give us the canonical form of the EMT. In Fig. 1 we show $`\stackrel{~}{D}`$ for Taub-NUT disks with different values of $`p`$ and $`\kappa `$. First we plot $`\stackrel{~}{D}`$ for disks with $`\kappa =1.7`$ and $`p=0.9`$, $`0.7`$, $`0.5`$, $`0.3`$, and $`0.1`$. Then we plot $`\stackrel{~}{D}`$ for disks with $`p=0.8`$ and $`\kappa =2`$, $`2.6`$, $`3.2`$, $`3.8`$, and $`4.4`$. We also computed $`\stackrel{~}{D}`$ for many other values of $`p`$ and $`\kappa `$. We found that in all the cases when $`p1`$, $`D`$ is not positive definite and it has only one root $`\stackrel{~}{r}_0>0`$. Therefore the Taub-NUT disks always have heat flow beginning at $`\stackrel{~}{r}=\stackrel{~}{r}_0`$. In Fig 2 we depict the density $`\stackrel{~}{\sigma }=\alpha \sigma `$ and the pressure $`\stackrel{~}{P}=\alpha P`$ for Taub-NUT disks with different values of $`p`$ and $`\kappa `$. First we plot $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{P}`$ (scaled by a factor $`10`$) for Taub-NUT disks with $`\kappa =1.4`$ and $`p=1`$, $`0.8`$, and $`0.6`$ as functions of $`\stackrel{~}{r}=r/m`$. Then we plot $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{P}`$ (scaled by a factor $`10`$) for Taub-NUT disks with $`\kappa =1.7`$ and $`p=1`$, $`0.8`$, and $`0.6`$ as functions of $`\stackrel{~}{r}=r/m`$. We also computed $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{P}`$ for many other values of $`p`$ and $`\kappa `$, in all the cases, we obtain a similar behaviour. The density $`\stackrel{~}{\sigma }`$ is always positive, $`\stackrel{~}{\sigma }>0`$; whereas the pressure becomes negative (tension) for a value of $`\stackrel{~}{r}<\stackrel{~}{r}_0`$. We can also see that $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{P}`$ presents a non smooth behaviour at $`\stackrel{~}{r}=\stackrel{~}{r}_0`$. For $`\stackrel{~}{r}>\stackrel{~}{r}_0`$ we have heat flow. The heat flow function $`\stackrel{~}{K}=\alpha K`$ is represented in Fig. 3. We plot $`\stackrel{~}{K}`$ for $`\kappa =1.4`$ (upper two curves) and $`\kappa =1.7`$ (lower two curves), with $`p=0.8`$ and $`0.6`$ as functions of $`\stackrel{~}{r}`$. Also in Fig. 3, in order to see the change of behaviour of $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{P}`$ at $`\stackrel{~}{r}=\stackrel{~}{r}_0`$ and the relation between $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{K}`$, we plot $`\stackrel{~}{\sigma }`$, $`\stackrel{~}{P}`$ and $`\stackrel{~}{K}`$ for $`\kappa =1.1`$ and $`p=0.8`$ in the interval $`2.5\stackrel{~}{r}3`$. We can see that $`\stackrel{~}{K}>\stackrel{~}{\sigma }`$ for $`\stackrel{~}{r}\stackrel{~}{r}_1>\stackrel{~}{r}_0`$. Thus, there is not causal propagation of heat for $`\stackrel{~}{r}>\stackrel{~}{r}_1`$. In Fig. 4 we show the tangential disk velocity $`V`$ and the disk angular momentum $`h`$ for Taub-NUT disks with $`p=0.8`$ and $`\kappa =1.1`$, $`1.4`$, $`1.7`$ and $`2`$, as functions of $`\stackrel{~}{r}`$. We see that there is a strong change in the slope of $`V`$ and $`h`$ at $`\stackrel{~}{r}=\stackrel{~}{r}_0`$. That means that there is a strong instability at this value of $`\stackrel{~}{r}`$. We also show the stream angular momenta $`h_+`$ and $`h_{}`$ in order to compare with the counter-rotating model. We also found instability, but in a different place. We computed also $`V`$ and $`h`$ for a wide range of the parameters $`p`$ and $`\kappa `$, we found always the same behaviour. In order to compare the counter-rotating model of the disk with the true disk rotation, we plot in Fig. 5 the tangential velocities of the counter-rotating streams, $`U_+`$ (top) and $`U_{}`$, and the product $`U_+U_{}`$ (bottom full line) and the $`P/\sigma `$ (full dots) for $`\kappa =1.1`$ and $`p=0.8`$, as functions of $`\stackrel{~}{r}`$. We can see the exact matching of the two curves, $`U_+U_{}`$ and $`P/\sigma `$, in perfect agreement with (33). Also in Fig. 5, in order to show the behavior of the counter-rotating model for the different disks models, we plot the $`P/\sigma `$ relation for $`\kappa =1.4`$ (upper curves) and $`\kappa =1.7`$ (lower curves), with $`p=1`$, $`0.8`$, and $`0.6`$ as functions of $`\stackrel{~}{r}`$. ## 5 Kerr Disks In this section we use the “displace, cut, and reflect” method to built disk solutions using the Kerr solution, that can be written as (1) with $`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}\left[{\displaystyle \frac{p^2x^2+q^2y^21}{(px+1)^2+q^2y^2}}\right],`$ (38a) $`\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}\left[{\displaystyle \frac{p^2x^2+q^2y^21}{p^2(x^2y^2)}}\right],`$ (38b) $`𝒜`$ $`=`$ $`{\displaystyle \frac{2\alpha q}{p}}\left[{\displaystyle \frac{(1y^2)(px+1)}{p^2x^2+q^2y^21}}\right],`$ (38c) where $`p=\alpha /m`$ , $`q=a/m`$ , with $`\alpha ^2=m^2a^2`$, so that $`p^2+q^2=1`$. $`m`$ and $`a`$ are the mass and the Kerr parameter, respectively. $`x`$ e $`y`$ are, again, the prolate spheroidal coordinates, given by (35a) and (35b). With the above expressions we can compute the trace and the discriminant of the energy-momentum tensor, that now can be written as $`\stackrel{~}{T}=mT`$ and $`\stackrel{~}{D}=m^2D`$, with $`\stackrel{~}{T}`$ $`=`$ $`{\displaystyle \frac{4\overline{x}\overline{y}(\overline{x}^21)(1\overline{y}^2)}{p(\overline{x}^2\overline{y}^2)^2(p^2\overline{x}^2+q^2\overline{y}^21)}}`$ $``$ $`2\overline{y}\left\{{\displaystyle \frac{2q^2\overline{x}(p\overline{x}+1)(1\overline{y}^2)+p(\overline{x}^21)[(p\overline{x}+1)^2q^2\overline{y}^2]}{p(\overline{x}^2\overline{y}^2)(p^2\overline{x}^2+q^2\overline{y}^21)[(p\overline{x}+1)^2+q^2\overline{y}^2]}}\right\}`$ $`\stackrel{~}{D}`$ $`=`$ $`4\overline{y}^2\left\{{\displaystyle \frac{2q^2\overline{x}(p\overline{x}+1)(1\overline{y}^2)+p(\overline{x}^21)[(p\overline{x}+1)^2q^2\overline{y}^2]}{p(\overline{x}^2\overline{y}^2)(p^2\overline{x}^2+q^2\overline{y}^21)[(p\overline{x}+1)^2+q^2\overline{y}^2]}}\right\}^2`$ (39b) $``$ $`{\displaystyle \frac{16q^2\overline{y}^2(\overline{x}^21)(1\overline{y}^2)[3p^2\overline{x}^2+4p\overline{x}q^2\overline{y}^2+1]^2}{p^2(\overline{x}^2\overline{y}^2)^2(p^2\overline{x}^2+q^2\overline{y}^21)^3[(p\overline{x}+1)^2+q^2\overline{y}^2]}},`$ where $`\overline{x}`$ and $`\overline{y}`$ are given by $`2p\overline{x}`$ $`=`$ $`\sqrt{\stackrel{~}{r}^2+(\kappa +p)^2}+\sqrt{\stackrel{~}{r}^2+(\kappa p)^2},`$ (40a) $`2p\overline{y}`$ $`=`$ $`\sqrt{\stackrel{~}{r}^2+(\kappa +p)^2}\sqrt{\stackrel{~}{r}^2+(\kappa p)^2},`$ (40b) where $`\stackrel{~}{r}=r/m`$ and $`\kappa =b/m`$. In Fig. 6 we show $`\stackrel{~}{D}`$ for Kerr disks with different values of $`p`$ and $`\kappa `$. We first plot $`\stackrel{~}{D}`$ for disks with $`\kappa =1.5`$ and $`p=1`$, $`0.9`$, $`0.8`$, $`0.6`$, and $`0.3`$. Then we show $`\stackrel{~}{D}`$ for disks with $`p=0.8`$ and $`\kappa =1.7`$, $`2`$, $`2.3`$, $`2.9`$, and $`3.8`$. We also computed $`\stackrel{~}{D}`$ for different values of $`p`$ and $`\kappa `$, we find that for some values of $`p`$ and $`\kappa `$, $`\stackrel{~}{D}>0`$ everywhere on the disk, whereas for the complementary case $`\stackrel{~}{D}`$ have two roots, $`0<\stackrel{~}{r}_1<\stackrel{~}{r}_2`$, so that $`\stackrel{~}{D}<0`$ when $`0<\stackrel{~}{r}_1<\stackrel{~}{r}_2`$. Accordingly we have two kinds of Kerr disks, ones with heat flow in an annular region ($`0<\stackrel{~}{r}_1<\stackrel{~}{r}_2`$) and others without heat flow. We first analyze some cases when there is no heat flow. In Fig. 7 we show the density $`\stackrel{~}{\sigma }=m\sigma `$ and the pressure $`\stackrel{~}{P}=mP`$ for these Kerr disks. First we plot $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{P}`$ (scaled by a factor $`100`$) for a disk with $`p=0.9`$ and $`\kappa =2.4`$, $`2.5`$, and $`3`$ as functions of $`\stackrel{~}{r}=r/m`$. Then we plot $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{P}`$ (scaled by a factor $`10`$) for disks with $`\kappa =2.5`$ and $`p=1`$, $`0.9`$, and $`0.8`$. We also computed $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{P}`$ for a variety of $`p`$ and $`\kappa `$, we found a similar behaviour. The density $`\stackrel{~}{\sigma }`$ is positive definite everywhere on the disk, falling to zero at infinity. There are values of $`p`$ and $`\kappa `$ such that the pressure $`\stackrel{~}{P}`$ is positive definite, also falling to zero at infinity, whereas in other cases the pressure $`\stackrel{~}{P}`$ is negative (tension) at the central region of the disks becoming positive for a value $`\stackrel{~}{r}_0>0`$; then falling to zero at infinity. The rotational motion of these disks models is shown in Fig. 8. We first plot $`V`$ for $`p=0.9`$ with $`\kappa =2.4`$, $`2.5`$, and $`3`$ and for $`p=0.8`$ with $`\kappa =2.5`$ (top curve) and then we plot the disk angular momentum $`h`$ for the same values of $`p`$ and $`\kappa `$. We can see that the disk velocity increases initially and then falls to zero at infinity and always is less than the light velocity. Also we see that the disk angular momentum is an increasing monotonic function of $`\stackrel{~}{r}`$. Thus these disks models are stable. Finally, we analyze a Kerr disk with heat flow, with $`p=0.9`$ and $`\kappa =1.5`$. In Fig. 9 we show $`\stackrel{~}{\sigma }`$, $`\stackrel{~}{P}`$ and $`\stackrel{~}{K}`$ for this values of $`p`$ and $`\kappa `$. We see that the density $`\stackrel{~}{\sigma }`$ is positive definite everywhere on the disk, falling to zero at infinity and that presents a change of slope at the values of $`\stackrel{~}{r}`$ when $`\stackrel{~}{D}=0`$, $`\stackrel{~}{r}_1`$ and $`\stackrel{~}{r}_2`$. The pressure $`\stackrel{~}{P}`$ is negative at the central region of the disk, becoming positive for a value of $`\stackrel{~}{r}>\stackrel{~}{r}_2`$. Also, $`\stackrel{~}{P}`$ presents a change of slope at $`\stackrel{~}{r}_1`$ and $`\stackrel{~}{r}_2`$. The heat flow function $`\stackrel{~}{K}`$ increases rapidly and becomes larger than $`\stackrel{~}{\sigma }`$ for a value of $`\stackrel{~}{r}>\stackrel{~}{r}_1`$, $`\stackrel{~}{K}`$ falls to zero, becoming less than $`\stackrel{~}{\sigma }`$ for a value of $`\stackrel{~}{r}<\stackrel{~}{r}_2`$. So, there is not causal propagation of heat on this region of the disk. We also plot in Fig. 9 the disk angular momentum, $`h`$. We see that the angular momentum presents two instabilities, at the values of $`\stackrel{~}{r}`$ when $`\stackrel{~}{D}=0`$, $`\stackrel{~}{r}_1`$ and $`\stackrel{~}{r}_2`$, and so the region of the disk where there is heat flow is highly instable. In order to see that these instabilities are present also in the counter-rotating model, we plot also the stream angular momenta, $`h_+`$ and $`h_{}`$ (scaled by a factor $`20`$). We analyzed several other cases of Kerr disks with heat flow, we found the same behaviour. ## 6 Discussion We presented two families of rotating disks based on the Taub-NUT and Kerr solutions. For the first family we find that the presence of heat flow is a generic property of this models, also generic is the radial instability produced by this flow. For the second family we find two sub classes one with and the other without heat flow. The heat flow is concentrated in an annular region that is highly instable. The relation found between the geodesic circular velocities and the “speed of sound” found is a generalization for the stationary case of the corresponding relation found by Morgan and Morgan that is the base of the counter-rotating model. The inclusion of current in the rotating models to have a rotating “hot” disk as the ones studied in is under consideration. The rotation add mass to the otherwise massless charges. These models are quite nontrivial. Also, the inclusion of radial pressure or tension is being considered. Finally, we want to mention that all the computation of this work was performed using the algebraic programming system Reduce . Acknowledgments We want to thank CNPq and FAPESP for financial support. Also G.A.G is grateful for the warm hospitality of the DMA-IMECC-UNICAMP where the main part of this work was performed.
warning/0006/hep-ph0006281.html	ar5iv	text	# Is 𝑈_{𝑒⁢3} really related to the solar neutrino solutions? ## Abstract It has been said that the measurements of $`U_{e3}`$ in the lepton flavor mixing matrix would help discriminate between the possible solar neutrino solutions under the natural conditions with the neutrino mass hierarchies of $`m_1m_2m_3`$ and $`m_1m_2m_3`$, where $`m_i`$ is the $`i`$-th generation neutrino absolute mass. However, it is not true, and the relation between $`\mathrm{sin}^22\theta _{12}`$ and $`U_{e3}`$ obtained by Akhmedov, Branco, and Rebelo is trivial in actual. We show in this paper that the value of $`U_{e3}`$ cannot predict the solar neutrino solutions without one additional nontrivial condition. hep-ph/0006281 DPNU-00-22 Recent neutrino oscillation experiments suggest the strong evidences of tiny neutrino masses and lepton flavor mixings. Studies of the lepton flavor mixing matrix, which is so-called Maki-Nakagawa-Sakata (MNS) matrix, will give us important cues of the physics beyond the standard model. The mixing angle between the second and the third generations is expected to be almost maximal, and the large mixing between the first and the second generations is also favored as the large angle MSW (MSW-L) solution. On the other hand, the mixing between the first and the third generations, which corresponds to $`U_{e3}`$ in the MNS matrix, is small as the present upper bound of CHOOZ experiments show $`U_{e3}<0.16`$. It is very interesting if value of $`U_{e3}`$ is related to the solar neutrino solutions. In Ref., Akhmedov, Branco, and Rebelo said that the measurements of $`U_{e3}`$ would help discriminate between the possible solar neutrino solutions under the natural conditions with the neutrino mass hierarchies of $`m_1m_2m_3`$ and $`m_1m_2m_3`$, where $`m_i`$ is the $`i`$-th generation neutrino absolute mass. However, it is not true, and the relation between $`\mathrm{sin}^22\theta _{12}`$ and $`U_{e3}`$ obtained in Ref. is trivial in actual. We will show in this paper that the value of $`U_{e3}`$ cannot predict the solar neutrino solutions without one additional nontrivial condition. This is because we know only four parameters, $`\mathrm{sin}^22\theta _{12},\mathrm{sin}^22\theta _{23},\mathrm{\Delta }m_{sol}^2`$, and $`\mathrm{\Delta }m_{ATM}^2`$, from experiments, although five parameters must be needed in order to obtain the MNS matrix, and its element $`U_{e3}`$. Let us start our discussions with each type of neutrino mass hierarchy. Neutrino mass spectra can be classified in three types as, Type A: $`m_1m_2m_3`$, Type B: $`m_1m_2m_3`$, and Type C: $`m_1m_2m_3`$. It is expected that $`\mathrm{\Delta }m_{ATM}^2\|m_3^2m_2^2\|`$ and $`\mathrm{\Delta }m_{sol}^2\|m_2^2m_1^2\|`$. By using $`\theta _{23}=\pi /4`$ and $`U_{e3}=ϵ(1)`$ according to the data of the Super-Kamiokande and the CHOOZ experiments, respectively, the MNS matrix is given by $$U=\left(\begin{array}{ccc}c& s& ϵ\\ \frac{1}{\sqrt{2}}(s+cϵ)& \frac{1}{\sqrt{2}}(csϵ)& \frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}(scϵ)& \frac{1}{\sqrt{2}}(c+sϵ)& \frac{1}{\sqrt{2}}\end{array}\right),$$ (1) where $`c\mathrm{cos}\theta _{12}`$ and $`s\mathrm{sin}\theta _{12}`$. The Majorana mass matrix of neutrino in the diagonal base of the charged lepton mass matrix is given by $`M_\nu `$ $`=`$ $`Udiag.(m_1,m_2,m_3)U^T,`$ (2) $`=`$ $`\left(\begin{array}{ccc}\mu & \frac{1}{\sqrt{2}}[ϵ(m_3\mu )+m_{}cs]& \frac{1}{\sqrt{2}}[ϵ(m_3\mu )m_{}cs]\\ \frac{1}{\sqrt{2}}[ϵ(m_3\mu )+m_{}cs]& \frac{1}{2}(m_3+\mu ^{}2m_{}csϵ)& \frac{1}{2}(m_3\mu ^{})\\ \frac{1}{\sqrt{2}}[ϵ(m_3\mu )m_{}cs]& \frac{1}{2}(m_3\mu ^{})& \frac{1}{2}(m_3+\mu ^{}2m_{}csϵ)\end{array}\right),`$ (6) where $$\mu m_1c^2+m_2s^2,\mu ^{}m_1s^2+m_2c^2,m_{}m_2m_1.$$ (7) In Type A with the mass hierarchy of $`m_1m_2m_3`$, the neutrino mass matrix of Eq.(3) is written by $`M_\nu `$ $`=`$ $`m_0\left(\begin{array}{ccc}\kappa & \alpha +\beta & \alpha \beta \\ \alpha +\beta & 1+\delta \delta ^{}& 1\delta \\ \alpha \beta & 1\delta & 1+\delta +\delta ^{}\end{array}\right),`$ (11) where we just normalize Eq.(3) by $`m_3`$ as $`m_0m_3/2`$, $`\kappa 2\mu /m_3`$, $`\alpha \sqrt{2}ϵ(1\mu /m_3)`$, $`\beta \sqrt{2}m_{}cs/m_3`$, $`\delta \mu ^{}/m_3`$, and $`\delta ^{}(2m_{}cs/m_3)ϵ.`$ The values of $`m_0,\kappa ,\beta `$, and $`\delta `$ are determined by the atmospheric and the solar neutrino solutions. Only $`\alpha `$ and $`\delta ^{}`$ are unknown parameters, since they have the free parameter $`ϵ`$. Equation (11) induces the mixing angles of $`\mathrm{tan}2\theta _{12}={\displaystyle \frac{\sqrt{2}\beta }{\delta \frac{\kappa }{2}}},\mathrm{sin}\theta _{13}={\displaystyle \frac{\alpha }{\sqrt{2}\left(1\frac{\kappa }{2}\right)}}.`$ (12) By using the approximations of $`\mathrm{\Delta }m_{ATM}^2m_3^2`$ and $`\mathrm{\Delta }m_{sol}^2m_2^2`$, they become $`\mathrm{tan}2\theta _{12}{\displaystyle \frac{\sqrt{2}\beta }{\mathrm{cos}2\theta _{12}\sqrt{\frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}}}},\mathrm{sin}\theta _{13}{\displaystyle \frac{\alpha }{\sqrt{2}\left(1\mathrm{sin}^2\theta _{12}\sqrt{\frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}}\right)}}.`$ (13) Here we must notice that the value of $`\beta `$ is determined by the atmospheric and the solar neutrino solutions. Only $`ϵ=\mathrm{sin}\theta _{13}`$ is the free parameter with $`U_{e3}(=\mathrm{sin}\theta _{13})<0.16`$, which makes the value of $`\alpha `$ be also free parameter. If $`O(\alpha )O(\beta )`$, which dose not have physical meaning, Eqs.(13) induce $$U_{e3}\frac{1}{\sqrt{2}}\mathrm{sin}2\theta _{12}\sqrt{\frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}}.$$ (14) The right-hand side of this equation<sup></sup><sup></sup>* Equation (14) is not the same as the result in Ref. $`\mathrm{sin}\theta _{13}\frac{1}{2}\frac{\mathrm{tan}2\theta _{12}}{(1+\mathrm{tan}^22\theta _{12})^{1/4}}\left(\frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}\right)^{1/2},`$ which can not apply to the large angle solutions. Our result of Eq.(14) can apply not only to the small angle solution but also to the large angle solutions. gives the following values of $`U_{e3}`$ corresponding to the solar neutrino solutions as $`U_{e3}`$ $``$ $`10^{1.5}(\mathrm{MSW}\mathrm{L}),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}10}^{3.5}(\mathrm{VO}),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}10}^3(\mathrm{MSW}\mathrm{S}).`$ (15) These results are the same as those of Ref.. It seems that the measurements of $`U_{e3}`$ can predict the solar neutrino solutions from Eq.(15). However, we must notice that the relation of Eq.(15) is satisfied just only in the case of $`O(\alpha )O(\beta )`$. This is the trivial condition, since $`\alpha `$ is the free parameter which has nothing to do with $`\beta `$ at all. In Ref., they have denoted $`\epsilon \alpha +\beta `$ and $`\epsilon ^{}\alpha \beta `$, and said that $`\epsilon +\epsilon ^{}`$ and $`\epsilon \epsilon ^{}`$ are expected to be of the same order if there are no accidental cancellations. However, the condition $`\epsilon +\epsilon ^{}\epsilon \epsilon ^{}`$ means $`\alpha \beta `$, which is not the natural condition but just the trivial assumption. Any relations between $`\alpha `$ and $`\beta `$ are considerable, and for example, if we take $`O(\alpha )O(\beta )`$, Eq.(14) becomes $$\mathrm{sin}\theta _{13}\frac{1}{\sqrt{2}}\mathrm{sin}2\theta _{12}\sqrt{\frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}}.$$ (16) We stress again that Eqs.(14) and (16) are just the trivial relations which are induced by the trivial assumptions of $`O(\alpha )O(\beta )`$ and $`O(\alpha )O(\beta )`$, respectively. We can similarly analyze the case of Type B Type B has two patterns of (b1) and (b2) according to the relative sign assignments of mass eigenvalues. The stability of the mixing angles against the quantum corrections strongly depends on the relative assignments of mass eigenvalues. with the mass hierarchy of $`m_1m_2m_3`$. In the case of (b1), which is $`M_\nu ^ddiag.(m_1,m_2,0)`$ in the first order, with the notation of $`m_2=m_1+d`$ $`(\|d\|m_1)`$ the neutrino mass matrix $`M_\nu `$ is given by $`M_\nu `$ $`=`$ $`m_1\left(\begin{array}{ccc}1+\kappa & \alpha +\beta & \alpha \beta \\ \alpha +\beta & \frac{1}{2}+\delta +\gamma +\delta ^{}& \frac{1}{2}+\delta \gamma \\ \alpha \beta & \frac{1}{2}+\delta \gamma & \frac{1}{2}+\delta +\gamma +\delta ^{}\end{array}\right),`$ (20) where $`\kappa ds^2/m_1`$, $`\alpha (ϵ/\sqrt{2})\left[1+(m_3ds^2)/m_1\right]`$, $`\beta dcs/\sqrt{2}m_1`$, $`\delta m_3/2m_1`$, $`\gamma dc^2/2m_1`$, and $`\delta ^{}(dcs/m_1)ϵ.`$ The values of $`\kappa ,\beta ,\delta `$, and $`\gamma `$ are determined by the atmospheric and the solar neutrino solutions. Only $`\alpha `$ and $`\delta ^{}`$ are free parameters, since they contain $`ϵ`$. We can easily obtain mixing angles $`\mathrm{tan}2\theta _{12}={\displaystyle \frac{\sqrt{2}\beta }{\gamma +\frac{\kappa }{2}}},\mathrm{sin}\theta _{13}={\displaystyle \frac{\sqrt{2}\alpha }{1+\kappa 2\delta }},`$ (21) from Eq.(20). By using the approximations of $`\mathrm{\Delta }m_{ATM}^2m_1^2,`$ and $`\mathrm{\Delta }m_{sol}^24m_1^2(\gamma +\frac{\kappa }{2})(1+\mathrm{tan}^2\theta _{12})^{1/2},`$ we can obtain $`\alpha {\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}\theta _{13},\beta {\displaystyle \frac{1}{4\sqrt{2}}}\mathrm{sin}2\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}}\right).`$ (22) Here we must notice that the value of $`\beta `$ is determined by the atmospheric and the solar neutrino solutions, and $`\alpha `$ $`(\mathrm{sin}\theta _{13})`$ is the free parameter. If $`O(\alpha )O(\beta )`$, which does not have physical meaning, Eqs.(22) induce $`U_{e3}{\displaystyle \frac{1}{4}}\mathrm{sin}2\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}}\right).`$ (23) This equation is the same as that of Ref.. As we have shown in the case of Type A, Eq.(23) is the trivial equation which is obtained from the trivial assumption of $`O(\alpha )O(\beta )`$, since $`\alpha `$ is the free parameter which has nothing to do with $`\beta `$. Similar discussions can be applied to the case of (b2), which is $`M_\nu ^ddiag.(m_1,m_2,0)`$ in the first order. With the notation of $`m_2=m_1+d`$ $`(\|d\|m_1)`$, the neutrino mass matrix $`M_\nu `$ is given by $`M_\nu `$ $`=`$ $`m_1\left(\begin{array}{ccc}(c^2s^2)+\kappa & \sqrt{2}cs+\alpha +\beta & \sqrt{2}cs+\alpha \beta \\ \sqrt{2}cs+\alpha +\beta & \frac{1}{2}(c^2s^2)+\delta +\gamma \delta ^{}& \frac{1}{2}(c^2s^2)+\delta \gamma \\ \sqrt{2}cs+\alpha \beta & \frac{1}{2}(c^2s^2)+\delta \gamma & \frac{1}{2}(c^2s^2)+\delta +\gamma +\delta ^{}\end{array}\right)`$ (27) where $`\kappa ds^2/m_1`$, $`\alpha (ϵ/\sqrt{2})[(c^2s^2)m_3/m_1+ds^2/m_1]`$, $`\beta dcs/\sqrt{2}m_1`$, $`\delta m_3/2m_1`$, $`\gamma dc^2/2m_1`$, $`\delta ^{}[(d2m_1)cs/m_1]ϵ`$ . The values of $`\kappa ,\beta ,\delta `$, and $`\gamma `$ are determined by the atmospheric and the solar neutrino solutions. The parameters $`\alpha `$ and $`\delta ^{}`$ are free, which contain $`ϵ`$. The mixing angles are induced from Eq.(LABEL:typeB-2) as $`\mathrm{tan}2\theta _{12}={\displaystyle \frac{\sqrt{2}(\beta \sqrt{2}cs)}{(c^2s^2)\gamma +\frac{\kappa }{2}}},\mathrm{sin}\theta _{13}={\displaystyle \frac{\sqrt{2}\alpha }{(c^2s^2)+\kappa 2\delta }}.`$ (29) By using the approximations of $`\mathrm{\Delta }m_{ATM}^2m_1^2[(c^2s^2\gamma +\kappa /2)^2+2(\sqrt{2}cs\beta )^2]`$ and $`\mathrm{\Delta }m_{sol}^24m_1^2(\gamma +\kappa /2)[(c^2s^2\gamma +\kappa /2)^2+2(\sqrt{2}cs\beta )^2]^{1/2},`$ the mixing angles in Eqs.(29) become $`\mathrm{tan}2\theta _{12}\sqrt{2}\beta ,\mathrm{sin}\theta _{13}\sqrt{2}\alpha ,(\mathrm{small}\mathrm{mixing}),`$ (30) $`\mathrm{sin}2\theta _{12}1\sqrt{2}\beta ,\mathrm{sin}\theta _{13}\sqrt{2}\alpha \left({\displaystyle \frac{d}{2m_1}}{\displaystyle \frac{m_3}{m_1}}\right)^1,(\mathrm{large}\mathrm{mixing}).`$ (31) If $`O(\alpha )O(\beta )`$, Eqs.(17) and (31) induce $`\mathrm{sin}\theta _{13}\mathrm{tan}\theta _{12},(\mathrm{small}\mathrm{mixing}),`$ (32) $`\mathrm{sin}\theta _{13}(1\mathrm{sin}2\theta _{12})\left({\displaystyle \frac{m_3}{m_1}}{\displaystyle \frac{d}{2m_1}}\right)^1,(\mathrm{large}\mathrm{mixing}).`$ (33) Similarly this is also the trivial relation. In Type C with the mass hierarchy of $`m_1m_2m_3`$, we show the case of (c4), for example, which is $`M_\nu ^ddiag.(m_1,m_2,m_3)`$ in the first order. With the notation of $`m_2=m_1+d`$, $`m_3=m_1+D`$, and $`\|d\|\|D\|m_1`$, the neutrino mass matrix $`M_\nu `$ is given by $`M_\nu `$ $`=`$ $`m_1\left(\begin{array}{ccc}1+\kappa & \alpha +\beta & \alpha \beta \\ \alpha +\beta & 1+\delta +\gamma \delta ^{}& \delta \gamma \\ \alpha \beta & \delta \gamma & 1+\delta +\gamma +\delta ^{}\end{array}\right)`$ (37) where $`\kappa ds^2/m_1,\alpha (ϵ/\sqrt{2})[D/m_1ds^2/m_1],\beta dcs/\sqrt{2}m_1,\delta D/2m_1,\gamma dc^2/2m_1,`$ and $`\delta ^{}(dcs/m_1)ϵ.`$ Equation (LABEL:typeC-4) induces the mixing angles as $`\mathrm{tan}\theta _{12}={\displaystyle \frac{\sqrt{2}\beta }{\gamma \frac{\kappa }{2}}},\mathrm{sin}\theta _{13}={\displaystyle \frac{\sqrt{2}\alpha }{2\delta \kappa }}.`$ (39) By using the approximations of $`\mathrm{\Delta }m_{sol}^24m_1^2(1+\gamma +\kappa /2)[(\gamma \kappa /2)^2+2\beta ^2]^{1/2}`$ and $`\mathrm{\Delta }m_{ATM}^24\delta m_1^2,`$ we can obtain $`{\displaystyle \frac{\beta }{\delta }}{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}2\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}}\right),{\displaystyle \frac{\alpha }{\delta }}\sqrt{2}\mathrm{sin}\theta _{13}.`$ (40) Under the trivial assumption of $`O(\alpha )O(\beta )`$, Eqs.(40) induce the trivial relation $`\mathrm{sin}\theta _{13}{\displaystyle \frac{1}{2}}\mathrm{sin}2\theta _{12}{\displaystyle \frac{\mathrm{\Delta }m_{sol}^2}{\mathrm{\Delta }m_{ATM}^2}}.`$ (41) It has been said that the measurements of $`U_{e3}`$ in the lepton flavor mixing matrix would help discriminate between the possible solar neutrino solutions under the natural conditions with the neutrino mass hierarchies of $`m_1m_2m_3`$ and $`m_1m_2m_3`$ . However, it is not true, and the relation between $`\mathrm{sin}^22\theta _{12}`$ and $`U_{e3}`$ obtained in Ref. is trivial in actual. Why can not we obtain the relations between the value of $`U_{e3}`$ and the solar neutrino solutions? This is easily understood as follows. Neglecting the $`CP`$ phases in the lepton sector, the number of independent parameters in the Majorana mass matrix of neutrino are six. Five parameters are enough to determine the MNS matrix, since overall factor in the neutrino mass matrix does not contribute to the MNS matrix. Thus we need five input parameters in order to determine the MNS matrix, and its element $`U_{e3}`$. Since the neutrino oscillation experiments except for the CHOOZ give us only four input parameters $`\mathrm{\Delta }m_{\mathrm{ATM}}^2`$, $`\mathrm{\Delta }m_{\mathrm{sol}}^2`$, $`\mathrm{sin}^22\theta _{12}`$, and $`\mathrm{sin}^22\theta _{23}`$, the value of $`U_{e3}`$ remains as an unknown parameter, which we only know the upper bound from CHOOZ experiments as $`U_{e3}<0.16`$. Therefore the value of $`U_{e3}`$ cannot predict the solar neutrino solutions without one additional nontrivial condition. The work of NH is supported by the Grant-in-Aid for Science Research, Ministry of Education, Science and Culture, Japan (No. 12740146, No. 12014208).
warning/0006/nucl-th0006086.html	ar5iv	text	# Charm meson scattering cross sections by pion and rho meson ## I INTRODUCTION Since charm quarks may lose appreciable energies in a quark-gluon plasma via gluon radiations, study of the charm meson spectrum in heavy ion collisions is expected to provide useful information on the properties of the quark-gluon plasma formed in these collisions . However, charm mesons may interact strongly with hadrons during later stage of heavy ion collisions, and this may also lead to changes in their final spectrum. To use charm mesons as a probe of the properties of the initial quark-gluon plasma thus requires the understanding of their interactions with hadrons. In a previous study , we have evaluated the charm meson scattering cross sections with pion and rho meson in a simple hadronic model that includes only the pseudoscalar-pseudoscalar-vector (PPV) meson interactions. In that study, we have also neglected the exchange of charm mesons as it is expected to be unimportant due to the large charm meson mass. Including form factors at the interaction vertices, we have obtained a thermally averaged total cross section of about 10 mb in the temperature range of interest. In a schematic model for the dynamics of heavy ion collisions, we have found that the inverse slope of the charm meson transverse momentum spectrum is significantly enhanced by their scatterings in the hadronic matter. As a result, the invariant mass spectrum of the dileptons from the decay of charm meson pairs is expected to be modified, which has been suggested as a possible explanation for the observed enhancement of intermediate-mass dileptons in heavy ion collisions at SPS energies . In Ref. , Matinyan and Müller have used a similar hadronic Lagrangian to evaluate the cross sections of charmonium absorption in hadronic matter. In contrast with the charm meson scattering cross sections, they have obtained very small charmonium absorption cross sections. The model has been extended in Refs. by using the local flavor SU(4) gauge invariance to include also interactions among three vector mesons and among four particles. Because of these additional interactions, the charmonium absorption cross sections are increased by an order-of-magnitude. In this paper, we shall use this extended hadronic Lagrangian to study the charm meson scattering cross sections by pion and rho meson. The paper is organized as follows. In Sec. II A, we introduce the hadronic Lagrangian based on the local flavor SU(4) gauge symmetry. The interaction Lagrangians that are relevant to charm meson scattering with pion and rho meson are then given in Sec. II B. We also derive in this section the scattering amplitudes for these processes and give their explicit expressions in Appendix A. Constraints on the scattering amplitudes as a result of the conservation of SU(4) flavor current are then discussed in Sec. II C, and an example is shown in detail in Appendix B. After addressing in Sec. II D the problem of on-shell divergence in some of the amplitudes, we fix the coupling constants in Sec. II E and introduce in Sec. II F the form factors at interaction vertices. Numerical results for the charm meson scattering cross sections are presented in Sec. III. In Sec. IV, we compare our results with previous ones obtained using the PPV coupling and including only diagrams with light meson exchanges. Finally, a summary is given in Sec. V. ## II CHARM MESON INTERACTIONS WITH HADRONS ### A hadronic Lagrangian with SU(4) Symmetry We have previously introduced a hadronic Lagrangian with SU(4) symmetry for studying the charmonium absorption cross sections by hadrons . It starts from the free Lagrangian for pseudoscalar and vector mesons, $`_0=\mathrm{Tr}\left(_\mu P^{}^\mu P\right){\displaystyle \frac{1}{2}}\mathrm{Tr}\left(F_{\mu \nu }^{}F^{\mu \nu }\right),`$ (1) where $`F_{\mu \nu }=_\mu V_\nu _\nu V_\mu `$, and $`P`$ and $`V`$ denote, respectively, the $`4\times 4`$ pseudoscalar and vector meson matrices in SU(4) . Introducing the minimal substitution, $`_\mu P`$ $``$ $`𝒟_\mu P=_\mu P{\displaystyle \frac{ig}{2}}[V_\mu ,P],`$ (2) $`F_{\mu \nu }`$ $``$ $`_\mu V_\nu _\nu V_\mu {\displaystyle \frac{ig}{2}}[V_\mu ,V_\nu ],`$ (3) leads to the following Lagrangian for the interacting hadrons: $``$ $`=`$ $`_0+ig\mathrm{Tr}\left(^\mu P[P,V_\mu ]\right){\displaystyle \frac{g^2}{4}}\mathrm{Tr}\left([P,V_\mu ]^2\right)`$ (4) $`+`$ $`ig\mathrm{Tr}\left(^\mu V^\nu [V_\mu ,V_\nu ]\right)+{\displaystyle \frac{g^2}{8}}\mathrm{Tr}\left([V_\mu ,V_\nu ]^2\right).`$ (5) Since hadron masses explicitly break the SU(4) symmetry, mass terms based on the experimentally determined values are added to Eq. (5). ### B scattering amplitudes The above Lagrangian yields the following processes for charm meson scattering by $`\pi `$ and $`\rho `$ mesons, $`\pi D\rho D^{},\pi D\pi D,\pi D^{}\pi D^{},\pi D^{}\rho D,\rho D\rho D,\rho D^{}\rho D^{}.`$ (6) There are also similar processes for anti-charm mesons. Fig. 1 shows the diagrams for the eight processes in Eq. (6). Expanding the Lagrangian in Eq. (5) using the $`4\times 4`$ matrices for $`P`$ and $`V`$, we obtain the following interaction Lagrangians that are relevant to charm meson scattering: $`_{\rho \pi \pi }`$ $`=`$ $`g_{\rho \pi \pi }\stackrel{}{\rho }^\mu \left(\stackrel{}{\pi }\times _\mu \stackrel{}{\pi }\right),`$ (7) $`_{\rho \rho \rho }`$ $`=`$ $`g_{\rho \rho \rho }_\mu \stackrel{}{\rho _\nu }\left(\stackrel{}{\rho ^\mu }\times \stackrel{}{\rho ^\nu }\right),`$ (8) $`_{\pi DD^{}}`$ $`=`$ $`ig_{\pi DD^{}}D^\mu \stackrel{}{\tau }\left(\overline{D}_\mu \stackrel{}{\pi }_\mu \overline{D}\stackrel{}{\pi }\right)+\mathrm{H}.\mathrm{c}.,`$ (9) $`_{\rho DD}`$ $`=`$ $`ig_{\rho DD}\left(D\stackrel{}{\tau }_\mu \overline{D}_\mu D\stackrel{}{\tau }\overline{D}\right)\stackrel{}{\rho }^\mu ,`$ (10) $`_{\rho D^{}D^{}}`$ $`=`$ $`ig_{\rho D^{}D^{}}[(_\mu D^\nu \stackrel{}{\tau }\overline{D}_\nu ^{}D^\nu \stackrel{}{\tau }_\mu \overline{D}_\nu ^{})\stackrel{}{\rho }^\mu `$ (11) $`+`$ $`(D^\nu \stackrel{}{\tau }_\mu \stackrel{}{\rho }_\nu _\mu D^\nu \stackrel{}{\tau }\stackrel{}{\rho }_\nu )\overline{D}^\mu +D^\mu (\stackrel{}{\tau }\stackrel{}{\rho }^\nu _\mu \overline{D}_\nu ^{}\stackrel{}{\tau }_\mu \stackrel{}{\rho }^\nu \overline{D}_\nu ^{})],`$ (12) $`_{\pi \rho DD^{}}`$ $`=`$ $`g_{\pi \rho DD^{}}D^\mu \left(2\stackrel{}{\tau }\stackrel{}{\pi }\stackrel{}{\tau }\stackrel{}{\rho _\mu }\stackrel{}{\tau }\stackrel{}{\rho _\mu }\stackrel{}{\tau }\stackrel{}{\pi }\right)\overline{D}+\mathrm{H}.\mathrm{c}.,`$ (13) $`_{\pi \pi D^{}D^{}}`$ $`=`$ $`g_{\pi \pi D^{}D^{}}\left({\displaystyle \frac{\stackrel{}{\pi }\stackrel{}{\pi }}{2}}\right)D^\mu \overline{D}_\mu ^{},`$ (14) $`_{\rho \rho DD}`$ $`=`$ $`g_{\rho \rho DD}\left({\displaystyle \frac{\stackrel{}{\rho _\mu }\stackrel{}{\rho ^\mu }}{2}}\right)D\overline{D},`$ (15) $`_{\rho \rho D^{}D^{}}`$ $`=`$ $`g_{\rho \rho D^{}D^{}}D^\mu \left(2\stackrel{}{\tau }\stackrel{}{\rho _\nu }\stackrel{}{\tau }\stackrel{}{\rho _\mu }\stackrel{}{\tau }\stackrel{}{\rho _\mu }\stackrel{}{\tau }\stackrel{}{\rho _\nu }\stackrel{}{\rho _\gamma }\stackrel{}{\rho ^\gamma }g_{\mu \nu }\right)\overline{D}^\nu .`$ (16) In the above, $`\stackrel{}{\tau }`$ are Pauli matrices; $`\stackrel{}{\pi }`$ and $`\stackrel{}{\rho }`$ denote the pion and rho meson isospin triplets, respectively; while $`D`$ and $`D^{}`$ denote the pseudoscalar and vector charm meson isospin doublets, respectively. Using the above interacting Lagrangians, we have derived the amplitudes for all diagrams in Fig. 1, and they are given in Appendix A. In general, the amplitude for a process $`n`$, before summing and averaging over external spins and isospins, is given by the coherent sum of all individual amplitudes that contributing to the process, i.e., $`_n=\left({\displaystyle \underset{i}{}}_{ni}^{\lambda _k\mathrm{}\lambda _l}\right)ϵ_{k\lambda _k}\mathrm{}ϵ_{l\lambda _l}_n^{\lambda _k\mathrm{}\lambda _l}ϵ_{k\lambda _k}\mathrm{}ϵ_{l\lambda _l}`$ (17) where $`i`$ runs through $`a,b,c`$ for process 1 and $`a,b,c,d`$ for all other processes, and $`ϵ_{j\lambda _j}`$ denotes the polarization vector of external vector meson $`j`$. ### C current conservation Since the Lagrangian in Eq. (5) is generated from the free Lagrangian in Eq. (1) by the minimal substitution, it is invariant under the local flavor SU(4) gauge transformation, thus it is also invariant under the global flavor SU(4) gauge transformation. This invariance remains valid after including degenerate pseudoscalar and degenerate vector meson mass terms, leading to the conservation of a SU(4) flavor current. As a result, the scattering amplitude for any process satisfies the following condition: $`_n^{\lambda _k\mathrm{}\lambda _l}p_{j\lambda _j}=0,`$ (18) where $`p_{j\lambda _j}`$ is the momentum of external vector meson $`j`$. As an example, we show explicitly in Appendix B that the condition $`_2^{\lambda \omega }p_{4\omega }=0`$ is indeed satisfied by the amplitude for process 2, $`\pi D\rho D^{}`$. ### D on-shell divergence The amplitudes for diagrams 3b, 4a, and 6a become singular when the exchanged mesons are on-shell. Since the on-shell process describes a two-step process, their contribution needs to be subtracted from the cross section. Several methods have been proposed to treat such a singularity . Since we are interested in charm meson scattering in hadronic matter, the exchanged meson is expected to acquire an imaginary self-energy due to collisional broadening. The one-step process then corresponds to keeping only the real part of the propagator for the exchanged meson. However, a consistent evaluation of this effect also requires the inclusion of vertex corrections due to the medium, which has not been carried out even for light meson scattering in hadronic matter. We thus follow Ref. by adding an imaginary part of 50 MeV to the self-energy of the exchanged meson in the above three diagrams. Since the width of $`D^{}`$ in vacuum is very small (about 44 KeV) , the amplitude for diagram 1c can also be very large when the center-of-mass energy of the initial pion and charm meson is close to the $`D^{}`$ mass. We thus also add an imaginary part of 50 MeV to the self-energy of the $`D^{}`$ meson in diagram 1c. In Sec. III we shall show that thermal averages of these cross sections do not change much for values of imaginary self-energy between $`5`$ and $`500`$ MeV. ### E coupling constants For the coupling constants in the interaction Lagrangians, we shall use empirical values if they are available, i.e. $`g_{\rho \pi \pi }=6.1`$, $`g_{\pi DD^{}}=4.4`$ . $`g_{\rho DD}=g_{\rho D^{}D^{}}=2.52`$ . Since there is little empirical information on other coupling constants, we use the SU(4) relations to determine their values, i.e., $`g_{\rho \rho \rho }=g_{\rho \pi \pi },g_{\pi \rho DD^{}}=g_{\pi DD^{}}g_{\rho DD},g_{\pi \pi D^{}D^{}}=2g_{\pi DD^{}}^2,g_{\rho \rho DD}=2g_{\rho DD}^2,g_{\rho \rho D^{}D^{}}=g_{\rho D^{}D^{}}^2.`$ (19) We note that the SU(4) symmetry gives the following relations among couplings constants: $`{\displaystyle \frac{g_{\rho \pi \pi }}{2}}(3.0)=g_{\pi DD^{}}(4.4)=g_{\rho DD}(2.5)={\displaystyle \frac{g_{\rho \rho \rho }}{2}}=g_{\rho D^{}D^{}}={\displaystyle \frac{g}{4}};`$ (20) $`g_{\pi \rho DD^{}}={\displaystyle \frac{g_{\pi \pi D^{}D^{}}}{2}}={\displaystyle \frac{g_{\rho \rho DD}}{2}}=g_{\rho \rho D^{}D^{}}={\displaystyle \frac{g^2}{16}}.`$ (21) The empirical values given in the parentheses are seen to agree reasonably with those predicted by the SU(4) symmetry. ### F form factors Because of the finite size of hadrons, form factors are needed at interaction vertices. In the present study, we take them to be the usual mono-pole form for vertices in the $`t`$ and $`u`$ channel processes, i.e., $`f_3(t\mathrm{or}u)={\displaystyle \frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2+𝐪^2}},`$ (22) where $`\mathrm{\Lambda }`$ is a cutoff parameter, and $`𝐪^2`$ is the squared three momentum transfer in the center-of-mass frame, given by $`(𝐩_\mathrm{𝟏}𝐩_\mathrm{𝟑})_{\mathrm{cm}}^2`$ and $`(𝐩_\mathrm{𝟏}𝐩_\mathrm{𝟒})_{\mathrm{cm}}^2`$, respectively, for the $`t`$ and $`u`$ channel processes. These form factors are different from that used in Ref. , where it is given by $`f(t)=(\mathrm{\Lambda }^2m^2)/(\mathrm{\Lambda }^2t)`$, since the latter is not suitable for diagrams involving the charmed meson exchange that has a large invariant four momentum transfer $`t`$. As in Ref. , form factors at $`s`$ channel vertices are taken as $`f_3(s)={\displaystyle \frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2+𝐤^2}},`$ (23) with $`𝐤`$ denoting the three momentum of either the incoming or outgoing particles in the center-of-mass, i.e., $`𝐤^2=p_{i,\mathrm{cm}}^2`$ or $`p_{f,\mathrm{cm}}^2`$. After introducing the form factors at three-point vertices, the form factors at four-point vertices can in principle be determined by requiring the total amplitude for a given process satisfies the current conservation condition of Eq. (18) . Since the uncertainty of form factors involving charm mesons is already large for three-point vertices and the gauge invariance is not valid once we use empirical vector meson masses, we choose not to follow this more involved approach. Instead, we simply take the form factors at four-point vertices to be $`f_4=\left({\displaystyle \frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2+\overline{𝐪^2}}}\right)^2,`$ (24) where $`\overline{𝐪^2}`$ is the average value of the squared three momenta in the form factors for the $`s`$, $`t`$, and $`u`$ channels, i.e., $`\overline{𝐪^2}={\displaystyle \frac{5}{6}}\left(p_{i,\mathrm{cm}}^2+p_{f,\mathrm{cm}}^2\right).`$ (25) Since there is no empirical information on form factors involving charm mesons, we shall use for simplicity the same value for all cutoff parameters and choose $`\mathrm{\Lambda }`$ as either $`1`$ or $`2`$ GeV to study the uncertainties of our results due to form factors. ### G thermally averaged cross sections For charm meson scattering in hadronic matter, it is useful to study the thermal average of their cross sections. For a hadronic matter at temperature $`T`$, this is given by $`\sigma v`$ $`=`$ $`{\displaystyle \frac{_{z_0}^{\mathrm{}}𝑑z\left[z^2(\alpha _1+\alpha _2)^2\right]\left[z^2(\alpha _1\alpha _2)^2\right]K_1(z)\sigma (s=z^2\mathrm{T}^2)}{4\alpha _1^2K_2(\alpha _1)\alpha _2^2K_2(\alpha _2)}},`$ (26) where $`\alpha _i=m_i/\mathrm{T}`$, $`z_0=\mathrm{max}(\alpha _1+\alpha _2,\alpha _3+\alpha _4)`$, $`K_n`$’s are modified Bessel functions, and $`v`$ is the relative velocity of initial-state particles in their collinear frame . ## III NUMERICAL RESULTS We first consider the case without form factors at interaction vertices, i.e., $`\mathrm{\Lambda }=\mathrm{}`$. In Fig. 2, the solid lines show the energy dependence of the total charm meson scattering cross section for a given initial state, i.e., processes 1 and 2 for $`\pi D`$ scattering, processes 3 and 4 for $`\pi D^{}`$ scattering, processes 5 and 6 for $`\rho D`$ scattering, and processes 7 and 8 for $`\rho D^{}`$ scattering. We have also evaluated the thermal average of these cross sections, and they are shown by the solid lines in Fig. 3 as functions of temperature. To study the effects due to form factors, we take the value for the cutoff parameter $`\mathrm{\Lambda }`$ as either $`2`$ or $`1`$ GeV. The results are shown in Figs. 2 and 3 by the dashed and dash-dotted curves, respectively. As expected, magnitude of the cross sections decreases with decreasing cutoff parameter. For the cutoff parameters used here, the cross sections for $`\pi D,\pi D^{},\rho D`$ and $`\rho D^{}`$ scatterings are all roughly between 10 and 20 mb. Compared to the case without form factors, we see that form factors only suppress modestly (by a factor of two of less) the total cross sections and their thermal averages. This is due to the dominance of elastic processes, which involve small momentum transfer near the threshold. In contrast, form factors suppress significantly the cross section for $`J/\psi `$ absorption by pion, e.g., the process $`\pi \psi D^{}\overline{D}`$ is reduced by as much as a factor of 8 due to its large threshold . In Figs. 4 and 5, we show by the dashed and dot-dashed curves the cross sections for individual processes and their thermal averages for a cutoff parameter of $`1`$ GeV. It is seen that cross sections for both processes 2 ($`\pi D\rho D^{}`$) and 4 ($`\pi D^{}\rho D`$) increase from zero at their respective threshold, while cross sections for processes 6 ($`\rho D\pi D^{}`$) and 8 ($`\rho D^{}\pi D`$) diverge near threshold because they are exothermic. For the four elastic processes 1, 3, 5, and 7, their cross sections are finite at threshold. We also note that elastic processes are much more important than corresponding inelastic processes. Fig. 6 shows the thermal averages of cross sections for processes 1, 3, 4 and 6 at $`T=150`$ MeV as functions of the imaginary self-energy $`\mathrm{\Gamma }`$ of the exchanged meson. It is seen that they do not vary much for values of $`\mathrm{\Gamma }`$ between 5 and 500 MeV. ## IV Comparison with Previous Results In our previous study of charm meson scattering cross sections , we have considered only the pseudoscalar-pseudoscalar-vector meson interactions, which lead to diagrams 1a-c, 2a-b, 3b-c, 4a, 4c, 5b-c, 6a, 6c, and 8a-b in Fig. 1. Since diagrams involving charm mesons are expected to be less important, we have evaluated only those diagrams that involve the exchange of light mesons ($`\pi `$ or $`\rho `$ meson), i.e., diagrams (a) in processes 1, 2, 4, 6, and 8. These results are shown by circles in Figs. 4 and 5. Except for process 1 ($`\pi D\pi D`$) near the threshold, they are close to present results that include also charm meson exchanges and contact terms. This comparison thus demonstrates explicitly that diagrams with light meson exchanges dominate the charm meson scattering cross sections. This is in contrast to the charmonium absorption cross section in hadronic matter, where inclusion of additional couplings among three vector mesons and among four particles increases the $`J/\psi `$ absorption cross section by pion by an order of magnitude . This is due to the absence of light meson exchanges in charmonium absorption by hadrons. We note that these additional interactions among three vector mesons and among four particles yield new processes with light meson exchanges, i.e., diagrams 3a, 5a, and 7a. We have checked that these diagrams also dominate the cross sections for these processes. As shown in Fig. 4 and 5, these elastic processes are more important than corresponding inelastic processes and thus increase significantly the total scattering cross sections of charm mesons by pion and rho meson. ## V SUMMARY In summary, we have studied the scattering cross sections of charm mesons by pion and rho meson using a gauge invariant hadronic Lagrangian generated from the SU(4) symmetry. This leads to interaction Lagrangians not only among two pseudoscalar mesons and one vector meson but also among three vector mesons as well as among four particles. We have found that the charm meson scattering cross sections are dominated by diagrams with light meson exchanges. For the processes considered previously based only on the pseudoscalar-pseudoscalar-vector meson interactions, these additional interaction Lagrangians do not introduce new diagrams with light meson exchanges, and their cross sections are thus not much affected. However, the interaction Lagrangians involving three vector mesons or four particles lead to new processes with light meson exchanges besides those considered previously and thus increase the total charm meson scattering cross sections by hadrons. Therefore, we expect as in the previous study that the charm meson spectra in heavy ion collisions can be significantly modified by hadronic scattering. ## ACKNOWLEDGMENTS This work was supported in part by the National Science Foundation under Grant No. PHY-9870038, the Welch Foundation under Grant No. A-1358, and the Texas Advanced Research Program under Grant Nos. FY97-010366-0068 and FY99-010366-0081. ## Appendix A In this appendix, we give the explicit expressions for the amplitudes of all diagrams in Fig. 1. We show only the reduced amplitudes without the polarization vectors of external vector mesons and before summing and averaging over external spins and isospins. For process 1, $`\pi D\pi D`$, we have $`_{1a}`$ $`=`$ $`g_{\rho \pi \pi }g_{\rho DD}\left(ıϵ_{ijk}\tau ^k\right)_{\alpha \beta }\left({\displaystyle \frac{1}{tm_\rho ^2}}\right)(su),`$ (A1) $`_{1b}`$ $`=`$ $`g_{\pi DD^{}}^2(\tau ^j\tau ^i)_{\alpha \beta }\left({\displaystyle \frac{1}{um_D^{}^2}}\right)\left[st{\displaystyle \frac{(m_D^2m_\pi ^2)^2}{m_D^{}^2}}\right],`$ (A2) $`_{1c}`$ $`=`$ $`g_{\pi DD^{}}^2(\tau ^i\tau ^j)_{\alpha \beta }\left({\displaystyle \frac{1}{sm_D^{}^2}}\right)\left[t+u{\displaystyle \frac{(m_D^2m_\pi ^2)^2}{m_D^{}^2}}\right].`$ (A3) For process $`2`$, $`\pi D\rho D^{}`$, we have $`_{2a}^{\lambda \omega }`$ $`=`$ $`g_{\rho \pi \pi }g_{\pi DD^{}}\left(ıϵ_{ijk}\tau ^k\right)_{\alpha \beta }(2p_1+p_3)^\lambda \left({\displaystyle \frac{1}{tm_\pi ^2}}\right)(p_1p_2p_3)^\omega ,`$ (A4) $`_{2b}^{\lambda \omega }`$ $`=`$ $`g_{\pi DD^{}}g_{\rho DD}(\tau ^j\tau ^i)_{\alpha \beta }(p_1+p_2+p_4)^\lambda \left({\displaystyle \frac{1}{um_D^2}}\right)(2p_1+p_4)_\omega ,`$ (A5) $`_{2c}^{\lambda \omega }`$ $`=`$ $`g_{\pi DD^{}}g_{\rho D^{}D^{}}(\tau ^i\tau ^j)_{\alpha \beta }(p_1p_2)^\gamma \left({\displaystyle \frac{1}{sm_D^{}^2}}\right)\left[g_{\gamma \gamma ^{}}{\displaystyle \frac{(p_1+p_2)_\gamma (p_1+p_2)_\gamma ^{}}{m_D^{}^2}}\right]`$ (A6) $`\times `$ $`\left[(p_1p_2p_4)^\lambda g^{\gamma ^{}\omega }+(p_3+p_4)^\gamma ^{}g^{\lambda \omega }+(p_1+p_2+p_3)^\omega g^{\gamma ^{}\lambda }\right],`$ (A7) $`_{2d}^{\lambda \omega }`$ $`=`$ $`g_{\pi \rho DD^{}}(\tau ^i\tau ^j2\tau ^j\tau ^i)_{\alpha \beta }g^{\lambda \omega }.`$ (A8) In the above, $`p_n`$ denotes the momentum of particle $`n`$. Our convention is such that particles $`1`$ and $`2`$ represent initial-state mesons while particles $`3`$ and $`4`$ represent final-state mesons on the left and right side of the diagrams shown in Fig. 1, respectively. For vector mesons, the indices $`\mu ,\nu ,\lambda `$, and $`\omega `$ denote the polarization components of external mesons while the indices $`\gamma `$ and $`\gamma ^{}`$ denote those of the exchanged meson. The indices $`i`$ and $`j`$ represent the isospin state of isospin-triplet mesons on the left of a diagram, while the indices $`\alpha `$ and $`\beta `$ represent those of isospin-doublet mesons on the right of a diagram. For the isospin-triplet meson in the propagator the index $`k`$ represents its isospin state. The amplitudes for processes 4, 6 and 8 are related to above amplitudes for process 2 by the crossing symmetry. For process 4, $`\pi D^{}\rho D`$, we then have $`_{4a}^{\nu \lambda }`$ $`=`$ $`\widehat{\mathrm{T}_4}_{2a}^{\lambda \omega },_{4b}^{\nu \lambda }=\widehat{\mathrm{T}_4}_{2c}^{\lambda \omega },_{4c}^{\nu \lambda }=\widehat{\mathrm{T}_4}_{2b}^{\lambda \omega },_{4d}^{\nu \lambda }=\widehat{\mathrm{T}_4}_{2d}^{\lambda \omega },`$ (A9) where $`\widehat{\mathrm{T}_4}`$ represents the replacement, $`p_2p_4,\nu \omega `$, and $`ij`$. For process 6, $`\rho D\pi D^{}`$, we have $`_{6a}^{\mu \omega }`$ $`=`$ $`\widehat{\mathrm{T}_6}_{2a}^{\lambda \omega },_{6b}^{\mu \omega }=\widehat{\mathrm{T}_6}_{2c}^{\lambda \omega },_{6c}^{\mu \omega }=\widehat{\mathrm{T}_6}_{2b}^{\lambda \omega },_{6d}^{\mu \omega }=\widehat{\mathrm{T}_6}_{2d}^{\lambda \omega },`$ (A10) where $`\widehat{\mathrm{T}_6}`$ represents the replacement, $`p_1p_3,\mu \lambda `$, and $`ij`$. For process 8, $`\rho D^{}\pi D`$, we have $`_{8a}^{\mu \nu }`$ $`=`$ $`\widehat{\mathrm{T}_8}_{2a}^{\lambda \omega },_{8b}^{\mu \nu }=\widehat{\mathrm{T}_8}_{2b}^{\lambda \omega },_{8c}^{\mu \nu }=\widehat{\mathrm{T}_8}_{2c}^{\lambda \omega },_{8d}^{\mu \nu }=\widehat{\mathrm{T}_8}_{2d}^{\lambda \omega },`$ (A11) where $`\widehat{\mathrm{T}_8}`$ represents the replacement, $`p_1p_3,p_2p_4,\mu \lambda `$, and $`\nu \omega `$. For process 3, $`\pi D^{}\pi D^{}`$, we have $`_{3a}^{\nu \omega }`$ $`=`$ $`g_{\rho \pi \pi }g_{\rho D^{}D^{}}\left(ıϵ_{ijk}\tau ^k\right)_{\alpha \beta }\left({\displaystyle \frac{1}{tm_\rho ^2}}\right)`$ (A12) $`\times `$ $`\left[(us)g^{\nu \omega }+4(p_1^\nu p_3^\omega p_3^\nu p_1^\omega )+p_2^\nu (p_1+p_3)^\omega +(p_1+p_3)^\nu p_4^\omega \right],`$ (A13) $`_{3b}^{\nu \omega }`$ $`=`$ $`g_{\pi DD^{}}^2(\tau ^j\tau ^i)_{\alpha \beta }(p_1p_3+p_4)^\nu \left({\displaystyle \frac{1}{um_D^2}}\right)(2p_1p_4)^\omega ,`$ (A14) $`_{3c}^{\nu \omega }`$ $`=`$ $`g_{\pi DD^{}}^2(\tau ^i\tau ^j)_{\alpha \beta }(2p_1+p_2)^\nu \left({\displaystyle \frac{1}{sm_D^2}}\right)(p_1p_2p_3)^\omega ,`$ (A15) $`_{3d}^{\nu \omega }`$ $`=`$ $`g_{\pi \pi D^{}D^{}}\delta _{ij}\delta _{\alpha \beta }g^{\nu \omega }.`$ (A16) For process 5, $`\rho D\rho D`$, we have $`_{5a}^{\mu \lambda }`$ $`=`$ $`g_{\rho \rho \rho }g_{\rho DD}\left(ıϵ_{ijk}\tau ^k\right)_{\alpha \beta }\left({\displaystyle \frac{1}{tm_\rho ^2}}\right)`$ (A17) $`\times `$ $`\left[(us)g^{\mu \lambda }+4(p_2^\mu p_4^\lambda p_4^\mu p_2^\lambda )+p_1^\mu (p_2+p_4)^\lambda +3(p_2+p_4)^\mu p_3^\lambda \right],`$ (A18) $`_{5b}^{\mu \lambda }`$ $`=`$ $`g_{\rho DD}^2(\tau ^j\tau ^i)_{\alpha \beta }(p_1+2p_4)^\mu \left({\displaystyle \frac{1}{um_D^2}}\right)(p_1+p_2+p_4)^\lambda ,`$ (A19) $`_{5c}^{\mu \lambda }`$ $`=`$ $`g_{\rho DD}^2(\tau ^i\tau ^j)_{\alpha \beta }(p_1+2p_2)^\mu \left({\displaystyle \frac{1}{sm_D^2}}\right)(p_1+p_2+p_4)^\lambda ,`$ (A20) $`_{5d}^{\mu \lambda }`$ $`=`$ $`g_{\rho \rho DD}\delta _{ij}\delta _{\alpha \beta }g^{\mu \lambda }.`$ (A21) For process 7, $`\rho D^{}\rho D^{}`$, we have $`_{7a}^{\mu \nu \lambda \omega }`$ $`=`$ $`g_{\rho \rho \rho }g_{\rho D^{}D^{}}\left(ıϵ_{ijk}\tau ^k\right)_{\alpha \beta }\left({\displaystyle \frac{1}{tm_\rho ^2}}\right)`$ (A22) $`\times `$ $`\left[(p_1+p_3)^\gamma g^{\mu \lambda }+(p_12p_3)^\mu g^{\gamma \lambda }+(2p_1+p_3)^\lambda g^{\mu \gamma }\right]g_{\gamma \gamma ^{}}`$ (A23) $`\times `$ $`\left[(p_2p_4)^\gamma ^{}g^{\nu \omega }+(p_1+p_2+p_3)^\omega g^{\gamma ^{}\nu }+(p_1p_3+p_4)^\nu g^{\gamma ^{}\omega }\right],`$ (A24) $`_{7b}^{\mu \nu \lambda \omega }`$ $`=`$ $`g_{\rho D^{}D^{}}^2\left(\tau ^j\tau ^i\right)_{\alpha \beta }\left({\displaystyle \frac{1}{um_D^{}^2}}\right)`$ (A25) $`\times `$ $`\left[(p_1+p_4)^\gamma g^{\mu \omega }+(p_12p_4)^\mu g^{\gamma \omega }+(2p_1+p_4)^\omega g^{\gamma \mu }\right]\left[g_{\gamma \gamma ^{}}{\displaystyle \frac{(p_1p_4)_\gamma (p_1p_4)_\gamma ^{}}{m_D^{}^2}}\right]`$ (A26) $`\times `$ $`\left[(p_1+p_3p_4)^\nu g^{\gamma ^{}\lambda }(p_2+p_3)^\gamma ^{}g^{\nu \lambda }+(p_1+p_2+p_4)^\lambda g^{\gamma ^{}\nu }\right],`$ (A27) $`_{7c}^{\mu \nu \lambda \omega }`$ $`=`$ $`g_{\rho D^{}D^{}}^2\left(\tau ^i\tau ^j\right)_{\alpha \beta }\left({\displaystyle \frac{1}{sm_D^{}^2}}\right)`$ (A28) $`\times `$ $`\left[(2p_1+p_2)^\nu g^{\gamma \mu }(p_1+2p_2)^\mu g^{\nu \gamma }+(p_1+p_2)^\gamma g^{\mu \nu }\right]\left[g_{\gamma \gamma ^{}}{\displaystyle \frac{(p_1+p_2)_\gamma (p_1+p_2)_\gamma ^{}}{m_D^{}^2}}\right]`$ (A29) $`\times `$ $`\left[(p_3+p_4)^\gamma ^{}g^{\lambda \omega }+(p_1p_2p_4)^\lambda g^{\gamma ^{}\omega }+(p_1+p_2+p_3)^\omega g^{\gamma ^{}\lambda }\right],`$ (A30) $`_{7d}^{\mu \nu \lambda \omega }`$ $`=`$ $`g_{\rho \rho D^{}D^{}}\left[(2\tau ^j\tau ^i\tau ^i\tau ^j)_{\alpha \beta }g^{\mu \nu }g^{\lambda \omega }+(2\tau ^i\tau ^j\tau ^j\tau ^i)_{\alpha \beta }g^{\mu \omega }g^{\nu \lambda }2\delta _{ij}\delta _{\alpha \beta }g^{\mu \lambda }g^{\nu \omega }\right].`$ (A31) After averaging (summing) over initial (final) spins and isopsins, the cross section for a process is given by $`{\displaystyle \frac{d\sigma _n}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{64\pi sp_{i,\mathrm{cm}}^2I_sI_i}}_n^{\lambda _k\mathrm{}\lambda _l}_n^{\lambda _k^{}\mathrm{}\lambda _l^{}}\left(g_{\lambda _k\lambda _k^{}}{\displaystyle \frac{p_{k\lambda _k}p_{k\lambda _k^{}}}{m_k^2}}\right)\mathrm{}\left(g_{\lambda _l\lambda _l^{}}{\displaystyle \frac{p_{l\lambda _l}p_{l\lambda _l^{}}}{m_l^2}}\right),`$ (A32) with $`s,t,u`$ being the standard Mandelstam variables, and $`p_{i,\mathrm{cm}}^2={\displaystyle \frac{\left[s(m_1+m_2)^2\right]\left[s(m_1m_2)^2\right]}{4s}}`$ (A33) is the squared momentum of initial-state mesons in the center-of-momentum frame. The factors $`I_s`$ and $`I_i`$ are due to averaging over initial spin and isospins, respectively. Values of $`I_s`$ are 1, 1, 3, 3, 3, 3, 9 and 9, respectively, for processes $`1`$ to $`8`$ in Fig. 1, while $`I_i`$ is 6 for all processes. ## Appendix B In this appendix, we show as an example that the scattering amplitude for process 2, $`\pi D\rho D^{}`$, satisfies the condition of Eq. (18) as a result of the SU(4) flavor current conservation. In particular, we shall prove that $`_2^{\lambda \omega }p_{4\omega }=0`$. Starting from Eq. (A8), we obtain $`_{2a}^{\lambda \omega }p_{4\omega }`$ $`=`$ $`g_{\rho \pi \pi }g_{\pi DD^{}}\left(ıϵ_{ijk}\tau ^k\right)_{\alpha \beta }\left({\displaystyle \frac{tm_D^2}{tm_\pi ^2}}\right)(2p_1)^\lambda ,`$ (B1) $`_{2b}^{\lambda \omega }p_{4\omega }`$ $`=`$ $`g_{\pi DD^{}}g_{\rho DD}\left(\tau ^j\tau ^i\right)_{\alpha \beta }\left({\displaystyle \frac{um_\pi ^2}{um_D^2}}\right)(p_1+p_2+p_4)^\lambda ,`$ (B2) $`_{2c}^{\lambda \omega }p_{4\omega }`$ $`=`$ $`g_{\pi DD^{}}g_{\rho D^{}D^{}}\left(\tau ^i\tau ^j\right)_{\alpha \beta }\left({\displaystyle \frac{1}{sm_D^{}^2}}\right)[(sm_\rho ^2)(p_1p_2)^\lambda `$ (B3) $`+`$ $`{\displaystyle \frac{(m_D^2m_\pi ^2)(m_D^{}^2m_\rho ^2)}{m_D^{}^2}}p_4^\lambda ],`$ (B4) $`_{2d}^{\lambda \omega }p_{4\omega }`$ $`=`$ $`g_{\pi \rho DD^{}}\left(\tau ^i\tau ^j2\tau ^j\tau ^i\right)_{\alpha \beta }p_4^\lambda .`$ (B5) In arriving at the above, we have discarded all terms with $`p_3^\lambda `$ as they vanish after contracting with the polarization vector $`ϵ_{3\lambda }`$. Using the SU(4) relation for the coupling constants shown in Eq. (21) then gives $`_2^{\lambda \omega }p_{4\omega }`$ $`=`$ $`{\displaystyle \frac{g^2}{8}}\left(\tau ^j\tau ^i\right)_{\alpha \beta }\left(m_D^2m_\pi ^2\right)\left({\displaystyle \frac{p_1^\lambda }{tm_\pi ^2}}+{\displaystyle \frac{p_2^\lambda }{um_D^2}}\right)`$ (B6) $`+`$ $`{\displaystyle \frac{g^2}{16}}\left(\tau ^i\tau ^j\right)_{\alpha \beta }[2\left({\displaystyle \frac{m_D^2m_\pi ^2}{tm_\pi ^2}}\right)p_1^\lambda +\left({\displaystyle \frac{m_D^{}^2m_\rho ^2}{sm_D^{}^2}}\right)(p_1p_2)^\lambda `$ (B7) $`+`$ $`{\displaystyle \frac{(m_D^2m_\pi ^2)(m_D^{}^2m_\rho ^2)}{m_D^{}^2(sm_D^{}^2)}}(p_1+p_2)^\lambda ].`$ (B8) With degenerate pseudoscalar meson masses and degenerate vector meson masses, the above expression then reduces to zero. For other amplitudes shown in Eqs. (A9)-(A31), the current conservation condition can be similarly proved. We note that $`_2^{\lambda \omega }p_{3\lambda }=0`$ holds for any masses.
warning/0006/hep-ph0006235.html	ar5iv	text	# University of Wisconsin - Madison MADPH-00-1180May 2000 Neutrino Mixing Schemes ## 1 Introduction A revolution in our understanding of the neutrino sector is underway, driven by observations that are interpreted in terms of changes in neutrino flavors as they propagate. Since neutrino oscillations occur only if neutrinos are massive, these phenomena indicate physics beyond the Standard Model. With the present evidence for oscillations from atmospheric, solar and accelerator data, we are already able to begin to make strong inferences about the mass spectrum and mixings of neutrinos. Theoretical efforts to achieve a synthesis have produced a variety of models with differing testable consequences. A combination of particle physics, nuclear physics and astrophysics is needed for a full determination of the fundamental properties of neutrinos. This article reviews what has been achieved thus far and the future prospects for understanding the nature of neutrino masses and mixing. ## 2 Two-Neutrino Analyses In a model with two neutrinos, the probability for a given neutrino flavor $`\nu _\alpha `$ oscillating into $`\nu _\beta `$ in a vacuum is $$P(\nu _\alpha \nu _\beta )=\mathrm{sin}^22\theta \mathrm{sin}^21.27\frac{\delta m^2L}{E},$$ (1) where $`\theta `$ is the two-neutrino mixing angle, $`\delta m^2`$ is the mass-squared difference of the two mass eigenstates in eV<sup>2</sup>, $`L`$ is the distance from the neutrino source to the detector in kilometers, and $`E`$ is the neutrino energy in GeV. ### 2.1 Atmospheric neutrinos The atmospheric neutrino experiments determine the ratios $`{\displaystyle \frac{N_\mu }{N_\mu ^0}}`$ $`=`$ $`\alpha [P(\nu _\mu \nu _\mu )+rP(\nu _e\nu _\mu ],`$ (2) $`{\displaystyle \frac{N_e}{N_e^0}}`$ $`=`$ $`\alpha [P(\nu _e\nu _e+r^1P(\nu _\mu \nu _e],`$ (3) where $`N_e^0`$ and $`N_\mu ^0`$ are the expected numbers of atmospheric $`e`$ and $`\mu `$ events, respectively, in the absence of oscillations, $`r\frac{N_e^0}{N_\mu ^0}`$, $``$ indicates an average over the neutrino spectrum, and $`\alpha `$ is an overall neutrino flux normalization. Atmospheric neutrino data have generally indicated that the expected number of muons detected is suppressed relative to the expected number of electrons atmos ; this suppression can be explained via neutrino oscillations oldatmos . The atmospheric data also indicate that $`N_e/N_e^0`$ is relatively flat with zenith angle, while $`N_\mu /N_\mu ^0`$ decreases with increasing zenith angle (i.e., with longer oscillation distance). Assuming $`\nu _\mu \nu _\tau `$ oscillations, the favored two-neutrino parameters are SuperKatmos $`\delta m^2`$ $`=`$ $`3.5\times 10^3\mathrm{eV}^2(1.5\text{}7\times 10^3\mathrm{eV}^2),`$ (4) $`\mathrm{sin}^22\theta `$ $`=`$ $`1.00(0.80\text{}1.00);`$ (5) the 90% C.L. allowed ranges are given in parentheses. The absolute normalization of the electron data indicates $`\alpha 1.18`$, which is within the theoretical uncertainties atmosth . The flatness versus $`L`$ of the electron data implies that simple $`\nu _\mu \nu _e`$ oscillations are strongly disfavored. Large amplitude ($`\mathrm{sin}^22\theta >0.2`$) $`\nu _\mu \nu _e`$ oscillations are also excluded by the CHOOZ reactor data CHOOZ for $`\delta m_{\mathrm{atm}}^2>10^3\mathrm{eV}^2`$. It is also possible that atmospheric $`\nu _\mu `$ are oscillating into sterile neutrinos. However, measurements of the upgoing zenith angle distribution and $`\pi ^0`$ production disfavor this possibility nusatm ; jglchap . ### 2.2 Solar neutrinos For the <sup>37</sup>Cl Cl and <sup>71</sup>Ga Ga experiments, the expected number of neutrino events is $$N=\sigma P(\nu _e\nu _e)(\beta \varphi _\mathrm{B}+\varphi _{\mathrm{non}\mathrm{B}})𝑑E,$$ (6) where we allow an arbitrary normalization factor $`\beta `$ for the <sup>8</sup>B neutrino flux since its normalization is not certain. For the Kamiokande Kam and Super-Kamiokande SuperKsolar experiments the interaction is $`\nu e\nu e`$ and the outgoing electron energy is measured. The number of events per unit of electron energy is then $`{\displaystyle \frac{dN}{dE_e}}`$ $`=`$ $`\beta {\displaystyle \left\{\frac{d\sigma _{CC}}{dE_e^{}}P(\nu _e\nu _e)+\frac{d\sigma _{NC}}{dE_e^{}}\left[1P(\nu _e\nu _e)\right]\right\}}`$ (7) $`\times G(E_e^{},E_e)\varphi _\mathrm{B}dE_\nu dE_e^{},`$ where $`d\sigma _{CC}/dE_e^{}(d\sigma _{NC}/dE_e^{})`$ is the charged-current (neutral-current) differential cross section for an incident neutrino energy $`E_\nu `$ and $`G(E_e^{},E_e)`$ is the probability that an electron of energy $`E_e^{}`$ is measured as having energy $`E_e`$. The neutrino fluxes are taken from the standard solar model (SSM) SSM . If $`\nu _e`$ oscillates into a sterile neutrino, $`\sigma _{NC}=0`$. For two-neutrino vacuum oscillations (VO) of $`\nu _e`$ into either $`\nu _\mu `$ or $`\nu _\tau `$ vacuum , the solar data indicate bw98vac ; bks98 $`\delta m^2`$ $`=`$ $`7.5\times 10^{11}\mathrm{eV}^2,`$ (8) $`\mathrm{sin}^22\theta `$ $`=`$ $`0.91,`$ (9) although there are also regions near $`\delta m^2=2.5\times 10^{10},4.4\times 10^{10},\mathrm{and}6.4\times 10^{10}\mathrm{eV}^2`$ that also give acceptable fits. For two-neutrino MSW oscillations MSW of $`\nu _e`$ into either $`\nu _\mu `$ or $`\nu _\tau `$, there are three possible solution regions bks98 . The small-angle MSW (SAM) solution is $`\delta m^2`$ $`=`$ $`7.5\times 10^6\mathrm{eV}^2,`$ (10) $`\mathrm{sin}^22\theta `$ $`=`$ $`0.01,`$ (11) and the large-angle MSW (LAM) solution is $`\delta m^2`$ $``$ $`10^5\mathrm{eV}^2,`$ (12) $`\mathrm{sin}^22\theta `$ $`=`$ $`0.6\text{}0.9.`$ (13) There is also a low $`\delta m^2`$ MSW (LOW) solution $`\delta m^2`$ $``$ $`10^7\mathrm{eV}^2,`$ (14) $`\mathrm{sin}^22\theta `$ $`=`$ $`0.6\text{}0.9`$ (15) although it is less favored. Note that all solutions except for small-angle MSW have large mixing in the solar sector. Two-neutrino oscillations of $`\nu _e`$ into a sterile neutrino are somewhat disfavored because sterile neutrinos do not have a NC interaction, which tends to make oscillation predictions for <sup>37</sup>Cl and SuperK similar, contrary to experimental evidence. Future measurements of NC interactions in the SNO detector SNO will provide a more thorough test for oscillations of solar $`\nu _e`$ into sterile neutrinos. SuperK and SNO will also provide an improved measurement of the <sup>8</sup>B neutrino spectrum in the near future, which should help distinguish the various solar scenarios. The Borexino borexino and ICARUS icarus experiments will provide a measurement of the <sup>7</sup>Be neutrinos, and could detect the strong seasonal dependence that exists for <sup>7</sup>Be neutrinos in many VO models seasonal . ### 2.3 LSND There are also indications of neutrino oscillations from the LSND accelerator experiment LSND . Their data suggest $`\nu _\mu \nu _e`$ oscillations with two-neutrino parameters, $$0.3\mathrm{eV}^2\delta m^2=\frac{0.03\mathrm{eV}^2}{(\mathrm{sin}^22\theta )^{0.7}}2.0\mathrm{eV}^2.$$ (16) Values of $`\delta m^2`$ above 10 eV<sup>2</sup> are also allowed for $`\mathrm{sin}^22\theta 0.0025`$caldwell . The future MiniBooNE experiment miniboone is expected to either confirm or refute the LSND result. ## 3 Global Analyses A complete description of all neutrino data requires a model that can explain all of the oscillation phenomena simultaneously. Since each of the three types of oscillation evidence (atmospheric, solar, LSND) requires a distinct $`\delta m^2`$ scale, and since $`N`$ neutrinos have only $`N1`$ independent mass-squared differences, four neutrinos are needed to completely explain all of the data. If one of the types of neutrino data is ignored, then it is in principle possible to explain the remaining data with a three-neutrino model. Because the LSND results have yet to be confirmed by another experiment, three-neutrino models are generally used in the context of describing the atmospheric and solar data. ### 3.1 Three neutrinos If the atmospheric and solar data are to be described by a three-neutrino model, there are two distinct mass-squared difference scales \[compare Eq. (4) with Eqs. (8), (10), (12) and (14)\]. The two possible patterns of three-neutrino masses are illustrated in Fig. 1. We assume without loss of generality that $`\|m_1\|<\|m_2\|<\|m_3\|`$, and that $`\delta m_{21}^2\delta m_{\mathrm{sun}}^2`$ and $`\delta m_{31}^2\delta m_{32}^2\delta m_{\mathrm{atm}}^2\delta m_{\mathrm{sun}}^2`$. Then the off-diagonal vacuum oscillation probabilities in a three-neutrino model are bww2 $`P(\nu _e\nu _\mu )`$ $`=`$ $`4\left\|U_{e3}U_{\mu 3}^{}\right\|^2\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}4e\left\{U_{e1}U_{e2}^{}U_{\mu 1}^{}U_{\mu 2}\right\}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{sun}}`$ (17) $`2J\mathrm{sin}2\mathrm{\Delta }_{\mathrm{sun}},`$ $`P(\nu _e\nu _\tau )`$ $`=`$ $`4\left\|U_{e3}U_{\tau 2}^{}\right\|^2\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}4e\left\{U_{e1}U_{e2}^{}U_{\tau 1}^{}U_{\tau 2}\right\}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{sun}}`$ (18) $`+2J\mathrm{sin}2\mathrm{\Delta }_{\mathrm{sun}},`$ $`P(\nu _\mu \nu _\tau )`$ $`=`$ $`4\left\|U_{\mu 3}U_{\tau 3}^{}\right\|^2\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}4e\left\{U_{\mu 1}U_{\mu 2}^{}U_{\tau 1}^{}U_{\tau 2}\right\}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{sun}}`$ (19) $`2J\mathrm{sin}2\mathrm{\Delta }_{\mathrm{sun}},`$ where $`U`$ is the neutrino mixing matrix (in the basis where the charged-lepton mass matrix is diagonal), $`\mathrm{\Delta }_j1.27(\delta m_j^2/\mathrm{eV}^2)(L/\mathrm{km})(E/\mathrm{GeV})`$ and $`J=m\left\{U_{e2}U_{e3}^{}U_{\mu 2}^{}U_{\mu 3}\right\}`$ is the CP-violating invariant jarlskog . For a discussion of CP violating effects in neutrino oscillations see Refs. 4nuCP ; deruj ; cpv The matrix elements $`U_{\alpha j}`$ are the mixing between flavor ($`\alpha =e,\mu ,\tau `$) and the mass ($`j=1,2,3`$) eigenstates. The CP-odd term changes sign under reversal of the oscillating flavors, or if neutrinos are replaced by anti-neutrinos. For either Dirac or Majorana neutrinos we choose the following parametrization for $`U`$ to describe neutrino oscillations $$U=\left(\begin{array}{ccc}c_{13}c_{12}& c_{13}s_{12}& s_{13}e^{i\delta }\\ c_{23}s_{12}s_{13}s_{23}c_{12}e^{i\delta }& c_{23}c_{12}s_{13}s_{23}s_{12}e^{i\delta }& c_{13}s_{23}\\ s_{23}s_{12}s_{13}c_{23}c_{12}e^{i\delta }& s_{23}c_{12}s_{13}c_{23}s_{12}e^{i\delta }& c_{13}c_{23}\end{array}\right),$$ (20) where $`c_{jk}\mathrm{cos}\theta _{jk}`$ and $`s_{jk}\mathrm{sin}\theta _{jk}`$. For Majorana neutrinos, $`U`$ contains two further phase factors, but these do not enter into oscillation phenomena. For the oscillation of neutrinos in atmospheric and long-baseline experiments with $`L/E>10^2`$km/GeV, the $`\mathrm{\Delta }_{\mathrm{sun}}`$ terms are negligible and the relevant vacuum oscillation probabilities are $`P(\nu _e\nu _\mu )`$ $`=`$ $`s_{23}^2\mathrm{sin}^22\theta _{13}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}},`$ (21) $`P(\nu _e\nu _\tau )`$ $`=`$ $`c_{23}^2\mathrm{sin}^22\theta _{13}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}},`$ (22) $`P(\nu _\mu \nu _\tau )`$ $`=`$ $`c_{13}^4\mathrm{sin}^22\theta _{23}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}.`$ (23) The diagonal oscillation probabilities for a given neutrino species can be found by subtracting all the off-diagonal probabilities involving that species from unity. For neutrinos from the sun, $`L/E10^{10}`$km/GeV, and the $`\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}`$ terms oscillate very rapidly, averaging to 1/2. Then the observable oscillation probability in a vacuum is $$P(\nu _e\nu _e)=1\frac{1}{2}\mathrm{sin}^22\theta _{13}c_{13}^4\mathrm{sin}^22\theta _{12}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{sun}}.$$ (24) When $`\theta _{13}=0`$ (i.e., $`U_{13}=0`$), the oscillations of atmospheric and long-baseline neutrinos decouple from those of solar neutrinos: at the $`\delta m_{\mathrm{atm}}^2`$ scale, there are pure $`\nu _\mu \nu _\tau `$ oscillations with amplitude $`\mathrm{sin}^22\theta _{23}`$, with no admixture of $`\nu _e`$, and at the $`\delta m_{\mathrm{sun}}^2`$ scale the $`\nu _e`$ oscillations are described by a simple two-neutrino formula with amplitude $`\mathrm{sin}^22\theta _{12}`$. Then the two-neutrino parameters for atmospheric and solar oscillations can be adopted directly in the three-neutrino case. If $`\theta _{13}0`$, then $`\nu _e`$ will participate in atmospheric and long-baseline oscillations. As discussed earlier, pure $`\nu _\mu \nu _e`$ oscillations at the $`\delta m_{\mathrm{atm}}^2`$ scale are strongly disfavored by the atmospheric data, but some small amount of $`\nu _\mu \nu _e`$ is still allowed. The CHOOZ reactor experiment measures $`\overline{\nu }_e`$ survival, and is sensitive to oscillations between $`\overline{\nu }_e`$ and $`\overline{\nu }_\mu `$ for $`\delta m_{\mathrm{atm}}^2>10^3\mathrm{eV}^2`$. The combined data from atmospheric experiments and CHOOZ favor $`\mathrm{sin}\theta _{13}=0`$ (i.e., pure $`\nu _\mu \nu _\tau `$ oscillations in the atmosphere) and suggest that $`\mathrm{sin}\theta _{13}<0.3`$ bww2 . The solar data also allow solar neutrinos to mix maximally, or nearly maximally. If we require both atmospheric and solar oscillations to be maximal, there is a unique three-neutrino solution to the neutrino mixing matrix bimaximal . This “bimaximal” mixing corresponds to $`\theta _{13}=0`$ and $`\|\theta _{12}\|=\|\theta _{23}\|=\pi /4`$, and is a special case of the decoupled solution for atmospheric and solar neutrinos. One interesting aspect of the bimaximal solution is that the solar $`\nu _e`$ oscillations are 50% into $`\nu _\mu `$ and 50% into $`\nu _\tau `$, although the flavor content of the $`\nu _e`$ oscillation is not observable in solar experiments. Further properties of the bimaximal and nearly bimaximal solutions and models that give rise to such solutions are discussed in Ref. bimaximal . ### 3.2 Four neutrinos As discussed earlier, four neutrinos are necessary to completely describe the atmospheric, solar and LSND results. A fourth light neutrino must be sterile since it is not observed in $`Z`$ decays zwidth . There must be three separate mass–squared scales which must satisfy the hierarchy $`\delta m_{sun}^2\delta m_{atm}^2\delta m_{LSND}^2`$. In the three–neutrino case the relation $`\delta m_{sun}^2\delta m_{atm}^2`$ leads to only one fundamental type of mass structure, in which one heavier mass is separated from two lighter, nearly degenerate states (or vice versa). In the four–neutrino case, however, there are two distinct types of mass structures: one heavier mass separated from three lighter, nearly degenerate states, or vice versa (which we will refer to as the $`1+3`$ spectrum), or two nearly degenerate mass pairs separated from each other (the $`2+2`$ spectrum). In each case, the largest separation scale is determined by the LSND scale. It can be shown BGG ; 4nuvar that only the $`2+2`$ spectrum is completely consistent with the positive oscillation signals in the solar, atmospheric and LSND experiments, and the negative results from the BUGEY reactor Bugey and CDHS accelerator CDHS experiments. Therefore our discussions below assume the $`2+2`$ case, which is illustrated in Fig. 2. Here we will assume that the mass splitting of the heavier pair drives the atmospheric oscillations, the mass splitting of the lighter pair drives the solar oscillations, and the separation of the mass pairs drives the LSND oscillations. Sterile neutrinos are also of interest in explaining $`r`$-process nucleosynthesis via supernova explosions (see e.g. Ref. caldwell ). The vacuum neutrino flavor oscillation probabilities, for an initially produced $`\nu _\alpha `$ to a finally detected $`\nu _\beta `$, can be written 4nuCP $$P(\nu _\alpha \nu _\beta )=\delta _{\alpha \beta }\underset{j<k}{}\left[4\mathrm{Re}(W_{\alpha \beta }^{jk})\mathrm{sin}^2\mathrm{\Delta }_{kj}2\mathrm{Im}(W_{\alpha \beta }^{jk})\mathrm{sin}2\mathrm{\Delta }_{kj}\right],$$ (25) where $$W_{\alpha \beta }^{jk}=U_{\alpha j}U_{\alpha k}^{}U_{\beta j}^{}U_{\beta k},$$ (26) and $`\mathrm{\Delta }_{kj}1.27\delta m_{kj}^2L/E`$. We assume that there are four mass eigenstates $`m_0,m_1,m_2,m_3,m_4`$, which are most closely associated with the flavor states $`\nu _s`$, $`\nu _e`$, $`\nu _\mu `$, and $`\nu _\tau `$, respectively. The solar oscillations are driven by $`\delta m_{01}^2`$, the atmospheric oscillations by $`\delta m_{32}^2`$, and the LSND oscillations by $`\delta m_{02}^2\delta m_{03}^2\delta m_{12}^2\delta m_{13}^2`$. Hence, the solar oscillations are primarily $`\nu _e\nu _s`$ and the atmospheric oscillations are primarily $`\nu _\mu \nu _\tau `$. For oscillations of solar $`\nu _e`$ to sterile neutrinos, the solar data disfavors large mixing; hence the most likely solar solution is MSW small mixing. It is also possible to reverse the roles of $`\nu _s`$ and $`\nu _\tau `$; however, current data disfavor oscillations of atmospheric $`\nu _\mu `$ to sterile neutrinos nusatm . The $`4\times 4`$ mixing matrix $`U`$ may be parametrized by 6 mixing angles ($`\theta _{01}`$, $`\theta _{02}`$, $`\theta _{03}`$, $`\theta _{12}`$, $`\theta _{13}`$, $`\theta _{23}`$) and 6 phases ($`\delta _{01}`$, $`\delta _{02}`$, $`\delta _{03}`$, $`\delta _{12}`$, $`\delta _{13}`$, $`\delta _{23}`$); only three of the phases are observable in neutrino oscillations. A complete analysis of the four-neutrino mixing matrix is quite complicated. However, the smallness of the mixing indicated by the LSND results suggests that $`\nu _e`$ does not mix strongly with the two heavier states, i.e., that $`\theta _{12}`$ and $`\theta _{13}`$ are small. Similarly, one can assume that the other light state does not mix strongly with the heavier states, i.e., $`\theta _{02}`$ and $`\theta _{03}`$ are also small. This situation occurs naturally in the explicit four-neutrino models in the literature. Then after dropping terms second order in small mixing angles, $`U`$ takes the form 4nuCP $$U=\left(\begin{array}{cccc}c_{01}& s_{01}^{}& s_{02}^{}& s_{03}^{}\\ & & & \\ s_{01}& c_{01}& s_{12}^{}& s_{13}^{}\\ & & & \\ c_{01}(s_{23}^{}s_{03}+c_{23}s_{02})& s_{01}^{}(s_{23}^{}s_{03}+c_{23}s_{02})& c_{23}& s_{23}^{}\\ +s_{01}(s_{23}^{}s_{13}+c_{23}s_{12})& c_{01}(s_{23}^{}s_{13}+c_{23}s_{12})& & \\ & & & \\ c_{01}(s_{23}s_{02}c_{23}s_{03})& s_{01}^{}(s_{23}s_{02}c_{23}s_{03})& s_{23}& c_{23}\\ s_{01}(s_{23}s_{12}c_{23}s_{13})& +c_{01}(s_{23}s_{12}c_{23}s_{13})& & \end{array}\right),$$ (27) where $`c_{jk}\mathrm{cos}\theta _{jk}`$ and $`s_{jk}\mathrm{sin}\theta _{jk}e^{i\delta jk}`$. We see that under these conditions, $`U`$ has approximately block diagonal form. The mixing of solar neutrinos is due to $`\theta _{01}`$ and the mixing of atmospheric neutrinos is due to $`\theta _{23}`$; the values for these mixing angles are essentially given by the two-neutrino fits in Sec. 2. Both vacuum and MSW solar oscillation solutions are allowed; for MSW oscillations to occur in the sun requires $`m_0>m_1`$. The mixing that drives the LSND oscillations is due to $`\theta _{12}`$ and $`\theta _{13}`$; the effective amplitude of the $`\nu _\mu \nu _e`$ oscillations in LSND is $$\mathrm{sin}^22\theta _{LSND}4\|s_{12}c_{23}+s_{13}s_{23}^{}\|^2.$$ (28) The assumption that $`\theta _{02}`$ and $`\theta _{03}`$ are small is not required by current data. However, most explicit models have the approximate form given by Eq. (27); see Sec. 4.2. If in fact $`\theta _{02}`$ and $`\theta _{03}`$ are not small, $`\nu _e`$ oscillations to the flavor associated with the other light state ($`\nu _s`$ or $`\nu _\tau `$) are possible at the LSND $`L/E`$ scale with an amplitude of the same order as the LSND oscillation amplitude. ## 4 Consequences for Masses and Mixings ### 4.1 Three-neutrino models The atmospheric and solar data put restrictions on the neutrino mixing angles and mass-squared differences, but do not at all constrain the absolute neutrino mass scale, which must be determined by other methods. The freedom to choose the mass scale allows a wide variety of possible mass matrix structures, even for the same mass-squared differences and mixing. The matrix $`U`$ that relates neutrino flavor eigenstates to neutrino mass eigenstates depends in general on mixing in both the neutrino and charged lepton sector. If $`U_{\mathrm{}}`$ diagonalizes the charged lepton mass matrix and $`U_\nu `$ the neutrino mass matrix, then $`U=U_{\mathrm{}}^{}U_\nu `$. Except where stated otherwise, in our discussions here we will work in the basis where the charged lepton mass matrix is diagonal, so that $`U=U_\nu `$. In general, the neutrino mass matrix in the flavor basis can then be written $`M_{\alpha \beta }=_jU_{\alpha j}m_jU_{\beta j}`$ for Majorana neutrinos or $`M_{\alpha \beta }=_jU_{\alpha j}m_jU_{\beta j}^{}`$ for Dirac neutrinos (these are the same if $`CP`$ is conserved, i.e., when $`U`$ is real). As an example of the different mass matrix structures that are possible, we consider the case when at least one of the masses is much smaller than the largest mass. Then there is one type of mass matrix of the form $`M=M_0+O(\delta m_{jk}^2)`$ (up to trivial sign changes) that can lead to maximal mixing of atmospheric neutrinos: $$M_0=\frac{m}{2}\left(\begin{array}{ccc}0& 0& 0\\ 0& a& b\\ 0& b& a\end{array}\right),$$ (29) where $`a,b1`$. If $`a=b`$, then there is only one large mass ($`m_1,m_2m_3m`$) and the form of Eq. (29) automatically fixes $`\theta _{23}\pi /4`$ and $`\theta _{13}0`$; the $`O(\delta m_{jk}^2)`$ terms determine $`\theta _{12}`$. Bimaximal mixing ($`\theta _{12}\pi /4`$) is obtained if the leading $`O(\delta m_{jk}^2)`$ terms have the form $$\mathrm{\Delta }M=ϵ\left(\begin{array}{ccc}0& 1& 1\\ 1& 0& 0\\ 1& 0& 0\end{array}\right),$$ (30) where $`ϵm`$; subleading $`O(\delta m_{jk}^2)`$ terms are then needed to split $`m_1`$ and $`m_2`$. If $`ab`$ in Eq. (29), then there are two large masses with $`\theta _{23}\pi /4`$ and $`\theta _{13}0`$. Small $`\theta _{12}`$, appropriate for the small-angle MSW solar solution, is achieved if the leading $`O(\delta m_{jk}^2)`$ terms then have the form $$\mathrm{\Delta }M=ϵ\left(\begin{array}{ccc}0& 1& 1\\ 1& 0& 0\\ 1& 0& 0\end{array}\right);$$ (31) see Ref. ramond for an example of a GUT model that has this form. In the situation where all masses are approximately degenerate, $`m\|m_1\|\|m_2\|\|m_3\|\delta m_{jk}^2`$, there are three different types of mass matrices of the form $`M=M_0+O(\delta m_{jk}^2)`$ (up to trivial changes in sign) that can lead to bimaximal mixing, depending on the relative signs of the $`m_j`$: $$M_0=m\left(\begin{array}{ccc}0& \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}& \frac{1}{2}& \frac{1}{2}\\ \frac{1}{\sqrt{2}}& \frac{1}{2}& \frac{1}{2}\end{array}\right),\mathrm{or}m\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right),\mathrm{or}m\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right).$$ (32) The requirement from neutrinoless double beta decay that the $`\nu _e\nu _e`$ element of the neutrino mass matrix be small (described below) implies that only the first case is allowed for Majorana neutrinos. The form of $`M_0`$ gives three degenerate neutrinos of mass $`m`$ and fixes two combinations of mixing angles ($`c_{13}s_{12}1/\sqrt{2}`$ and $`c_{23}c_{12}s_{13}s_{23}s_{12}1/2`$), while the remaining degree of freedom among the mixing angles and the neutrino mass splittings are determined by the $`O(\delta m_{jk}^2)`$ terms. If the leading $`O(\delta m_{jk}^2)`$ terms are proportional to the mass matrix in Eq. (29) with $`a=b`$, $`m_3`$ is split from $`m_1`$ and $`m_2`$, $`\theta _{13}0`$ and bimaximal mixing is obtained. Subleading $`O(\delta m_{jk}^2)`$ terms then provide the splitting between $`m_1`$ and $`m_2`$. A different mixing scheme occurs if the neutrino mass matrix is approximately proportional to unity and the charged lepton mass matrix is close to the so-called democratic form demo $$M_{\mathrm{}}=\frac{m_{\mathrm{}}}{3}\left(\begin{array}{ccc}1& 1& 1\\ 1& 1& 1\\ 1& 1& 1\end{array}\right);$$ (33) then there is one large eigenvalue $`m_{\mathrm{}}m_\tau `$ and two constraints on the flavor mixing angles, $`c_{13}c_{23}1/\sqrt{3}`$ and $`s_{23}s_{12}s_{13}c_{23}c_{12}1/\sqrt{3}`$. If a small perturbation is added to the lower right element of $`M_{\mathrm{}}`$, the muon gets mass and $`\theta _{13}`$ is constrained to be approximately zero; then there is maximal mixing in the solar sector ($`\mathrm{sin}^22\theta _{12}=1`$) and nearly maximal mixing in the atmospheric sector ($`\mathrm{sin}^22\theta _{23}=8/9`$demo . Additional perturbations to the diagonal elements of the charged lepton and neutrino mass matrices are then needed to give the electron-muon and neutrino mass splittings, respectively. An $`SO(10)`$ SUSY GUT model with a minimal Higgs sector can provide large $`\nu _\mu `$$`\nu _\tau `$ mixing for atmospheric neutrinos, and can accommodate either small or large mixing of solar neutrinos albright . Other possible neutrino mass matrix textures are discussed in Refs. textures , mohap and king-albright . Although neutrino oscillations are not sensitive to the overall neutrino mass scale, there are other processes that do depend on absolute masses. For example, studies of the tritium beta decay spectrum put an upper limit on the effective mass of the electron neutrino $$m_{\nu _e}\underset{j}{}\|U_{ej}\|^2m_j,$$ (34) of about 2.5 eV tritium ; in a three-neutrino model this implies an upper limit of 2.5 eV on the heaviest neutrino bwwnumass . The current limit on the magnitude of the $`\nu _e`$$`\nu _e`$ element of the neutrino mass matrix for Majorana neutrinos from neutrinoless double beta ($`0\nu \beta \beta `$) decay doi is of order 0.5 eV baudis , taking into account the imprecise knowledge of the nuclear matrix element and the sensitivity of a background fluctuation analysisprivate . For the three-neutrino model this implies $`\|M_{\nu _e\nu _e}\|`$ $`=`$ $`\|{\displaystyle \underset{j}{}}U_{ej}m_jU_{ej}\|`$ $`=`$ $`\|(c_{13}c_{12})^2m_1+(c_{13}s_{12})^2m_2e^{i\varphi _2}+s_{13}^2m_3e^{i\varphi _3}\|M_{max}=0.2\mathrm{eV},`$ where $`\varphi _2`$ and $`\varphi _3`$ are extra phases present for Majorana neutrinos that are not observable in neutrino oscillations. For models with one large mass $`m_1,m_2m_3\sqrt{\delta m_{atm}^2}0.05`$ eV, and Eq. (4.1) does not provide any additional constraint. However, if all three masses are nearly degenerate ($`m_1m_2m_3m`$), the $`0\nu \beta \beta `$ decay limit becomes $`s_{12}^2[12s_{13}^2(M_{max}/m)]/(2c_{13}^2)`$, which in turn implies that the solar $`\nu _e\nu _e`$ oscillation amplitude (see Eq. (24)) has the constraint bwmaj $$4c_{13}^4s_{12}^2c_{12}^21\left(\frac{M_{max}}{m}\right)^22s_{13}^2\left(1+\frac{M_{max}}{m}\right).$$ (36) For any value of $`m>M_{max}/(12s_{13}^2)`$ there will be a lower limit on the size of the solar $`\nu _e\nu _e`$ oscillation amplitude; for example, given the current limit on $`s_{13}`$, the small-angle MSW solar solution is ruled out for nearly degenerate Majorana neutrinos if $`m>0.25`$ eV bwmaj . Large-angle MSW and vacuum solar solutions, which have large mixing, are still allowed; any solar solution with maximal mixing can never be excluded by this bound. Neutrino mass also provides an ideal hot dark matter component hdm ; the contribution of neutrinos to the mass density of the universe is given by $`\mathrm{\Omega }_\nu =m_\nu /(h^293`$ eV), where $`h`$ is the Hubble expansion parameter in units of 100 km/s/Mpcexpansion . For $`h=0.65`$, the model with three nearly degenerate neutrinos has $`\mathrm{\Omega }_\nu 0.075(m/`$eV). In three-neutrino models with hierarchical masses, in which the largest mass is of order $`\sqrt{\delta m_{atm}^2}`$, $`\mathrm{\Omega }_\nu `$ is much smaller, on the order of 0.001. Future measurements of the hot dark matter component may be sensitive to $`m_\nu `$ down to 0.4 eV eht . ### 4.2 Four-neutrino models As described in Sec. 3.2, four-neutrino models must have the $`2+2`$ mass structure, i.e., two nearly-degenerate pairs separated from each other by the LSND scale. One possible class of mass matrices that can give this situation is $$M=m\left(\begin{array}{cccc}ϵ_1& ϵ_2& 0& 0\\ ϵ_2& 0& 0& ϵ_3\\ 0& 0& ϵ_4& 1\\ 0& ϵ_3& 1& ϵ_5\end{array}\right),$$ (37) presented in the ($`\nu _s,\nu _e,\nu _\mu ,\nu _\tau `$) basis (i.e., the basis in which the charged lepton mass matrix is diagonal). When $`ϵ_5=ϵ_4`$, the mass matrix in Eq. (37) can accommodate any of the three solar solutions, depending on the hierarchy of the mass matrix elements 4nuvar $`\mathrm{SAM}:`$ $`ϵ_2ϵ_1,ϵ_4ϵ_31,`$ (38) $`\mathrm{LAM}:`$ $`ϵ_1,ϵ_2,ϵ_4ϵ_31,`$ (39) $`\mathrm{VO}:`$ $`ϵ_1ϵ_2ϵ_4ϵ_31.`$ (40) In each case, the mass eigenvalues have the hierarchy $`m_1<m_0m_2,m_3`$, as required for the MSW solution. The mass matrix contains five parameters, just enough to incorporate the required three mass-squared differences and the oscillation amplitudes for the solar and LSND neutrinos. The large amplitude for atmospheric oscillations does not require an additional parameter, since the mass matrix naturally gives nearly-maximal mixing of $`\nu _\mu `$ with $`\nu _\tau `$. For the VO case, $`ϵ_1`$ does not contribute to the phenomenology and can be set to zero, so that there are only four parameters; the large amplitude for solar oscillations also occurs naturally in this case 4nuvar ; roy . Another variant with five parameters is $`ϵ_5=ϵ_4`$ and $`ϵ_2ϵ_1ϵ_3ϵ_41`$ gibbons . In this case, $`ϵ_3`$ determines both the amplitude of the LSND oscillations and causes the splitting between $`m_2`$ and $`m_3`$ that drives the atmospheric oscillations. Two testable consequences of this model are that $`\delta m_{atm}^21.3\times 10^3`$ eV<sup>2</sup> and that there should be observable $`\nu _e\nu _\tau `$ oscillations in short-baseline experiments. For both of the previous cases ($`ϵ_5=ϵ_41`$ and $`ϵ_5=ϵ_41`$), the heavier mass pair is much heavier than the lighter mass pair, i.e., $`m_1<m_0m_2,m_3`$. Yet another possibility is to have $`ϵ_1=ϵ_5=0`$ and $`ϵ_3,ϵ_4ϵ_2<1`$, where $`ϵ_2`$ is not small compared to unity roy . In this case, the two degenerate pairs of masses have mass eigenvalues that are the same order of magnitude; there are only four parameters as both the solar and atmospheric mixings are naturally close to maximal. However, the $`m_0`$$`m_1`$ splitting in this model can give only MSW solar oscillations, which for large mixing are disfavored when $`\nu _e\nu _s`$. Other four-neutrino mass matrix ansatzes have been presented in the literature other4nu , but they generally have characteristics similar to those discussed here. In all of the explicit four-neutrino models discussed above, since the $`\nu _e`$$`\nu _e`$ element of the neutrino mass matrix vanishes, there is no contribution to neutrinoless double beta decay. However, because $`m_3`$ and $`m_4`$ are always of order 1 eV or more (to provide the necessary mass-squared splitting for LSND oscillations), these models always contribute to the hot dark matter in the universeprimack . ## 5 Long-baseline experiments Long-baseline experiments (with $`L/E_\nu 10^2`$$`10^3`$ km/GeV) are expected to confirm the $`\nu _\mu \nu _\mu `$ disappearance oscillations at the $`\delta m_{\mathrm{atm}}^2`$ scale. The K2K experimentk2k from KEK to SuperK, with a baseline of $`L250`$ km and a mean neutrino energy of $`E_\nu 1.4`$ GeV is underway. The MINOS experiment from Fermilab to Soudanminos , with a longer baseline $`L730`$ km and higher mean energies $`E_\nu =3`$, 6 or 12 GeV, is under construction and the ICANOEicanoe and OPERAopera experiments, with similar baselines from CERN to Gran Sasso, have been approved. These experiments with dominant $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ beams will securely establish the oscillation phenomena and may measure $`\delta m_{\mathrm{atm}}^2`$ to a precision of order 10%thesis . They will also measure the dominant oscillation amplitude $`\mathrm{sin}^22\theta _{23}`$. Further tests of the neutrino parameters will likely require higher intensity neutrino beams, and $`\nu _e`$, $`\overline{\nu }_e`$ beams as well as $`\nu _\mu `$, $`\overline{\nu }_\mu `$, such as those generated in a neutrino factory deruj ; nufact1 ; nufact2 ; nufact3 . The $`\nu _e`$, $`\overline{\nu }_e`$ components of the beam allow one to search for $`\nu _e\nu _\mu `$ and $`\overline{\nu }_e\overline{\nu }_\mu `$ appearance, which will occur in the leading $`\delta m_{atm}^2`$ oscillation if $`\mathrm{sin}^22\theta _{13}0`$. Depending on the machine parameters, $`\delta m_{\mathrm{atm}}^2`$ and $`\mathrm{sin}^22\theta _{23}`$ can be measured to an accuracy of 1–2%, and $`\mathrm{sin}^22\theta _{13}`$ can be tested down to 0.01 or below nufact3 . If the baseline is long enough ($`L>1000`$ km), matter effects in the Earth will also play an important role in an appearance experiment: for $`\delta m_{\mathrm{atm}}^2>0`$ ($`\delta m_{\mathrm{atm}}^2<0`$), $`\nu _e\nu _\mu `$ oscillations are enhanced (suppressed) and $`\overline{\nu }_e\overline{\nu }_\mu `$ oscillations are suppressed (enhanced). Therefore by comparing $`\nu _e\nu _\mu `$ with $`\overline{\nu }_e\overline{\nu }_\mu `$ oscillations it may be possible to determine the sign of $`\delta m_{\mathrm{atm}}^2`$ nufact3 . The enhancement due to matter of either $`\nu _e\nu _\mu `$ or $`\overline{\nu }_e\overline{\nu }_\mu `$ may also improve the $`\mathrm{sin}^22\theta _{13}`$ sensitivity for appropriate choices of $`L`$ and $`E_\nu `$ nufact3 . In a four-neutrino model, both short-baseline experiments with $`L/E_\nu 1`$ km/GeV (which probe $`\delta m_{\mathrm{LSND}}^2`$) and long-baseline experiments will be required to obtain maximal information on the neutrino mixing parameters 4nuCP . $`CP`$-violating effects due to the phase $`\delta `$ only become appreciable at the subleading $`\delta m^2`$ scale, and only then if the mixing angles are large enough CP3 . In a three-neutrino model with $`\delta m_{\mathrm{sun}}^2`$ and $`\delta m_{\mathrm{atm}}^2`$, $`CP`$ violation will be observable only for the large-angle MSW solar solution; a long-baseline experiment with a high-intensity neutrino beam from a neutrino factory may be able to give information on $`\delta `$ in this case CPlong . In a four-neutrino model, potentially large $`CP`$-violating effects are possible at the $`\delta m_{\mathrm{atm}}^2`$ scale 4nuCP ; CPlarge . ## 6 Summary and outlook In a three-neutrino world, it may be possible to completely determine the neutrino mixing matrix and two independent mass-squared differences using existing and planned neutrino oscillation experiments. Future measurements of solar neutrinos should pin down the neutrino mass and mixing parameters $`\delta m_{21}^2`$ and $`\mathrm{sin}^22\theta _{12}`$ that are predominantly responsible for the solar neutrino deficit. Long-baseline experiments can provide a more precise determination of $`\delta m_{32}^2`$ and $`\mathrm{sin}^22\theta _{23}`$, which drive the atmospheric neutrino anomaly, measure the size of $`\mathrm{sin}^22\theta _{13}`$, and determine the sign of $`\delta m_{32}^2`$. If the solar solution is large-angle MSW, long-baseline experiments may also be sensitive to the $`CP`$-violating phase $`\delta `$. Future measurements of beta decay spectra, double-beta decay and hot dark matter may then help determine the last remaining neutrino mass parameter, the absolute neutrino mass scale. In a four-neutrino world, there are three additional mixing angles and three additional $`CP`$ phases. Since the extra neutrino is sterile, it may be difficult to determine some of the additional mixing angles, especially if they are small, as is the case in many existing models. Short baseline experiments that probe the $`\delta m_{\mathrm{LSND}}^2`$ oscillation will be helpful in making sense of the larger parameter space. Furthermore, in four-neutrino models $`CP`$ violation may become observable at the $`\delta m_{\mathrm{atm}}^2`$ scale (rather than the $`\delta m_{\mathrm{sun}}^2`$ scale, as is the case in three-neutrino models), and there will be a contribution to hot dark matter on the order of $`m_\nu 1`$ eV.
warning/0006/hep-ph0006107.html	ar5iv	text	# Model–independent analysis of the exclusive semileptonic 𝐵→𝜋⁺⁢𝜋⁻⁢ℓ⁺⁢ℓ⁻ decay ## 1 Introduction Rare $`B`$ meson decays, induced by flavor–changing neutral current (FCNC) $`bs`$ or $`bd`$ transitions, occur at loop level in the standard model (SM) and is a potential precision testing ground for the SM, as well as attracting the theoretical interest as a promising tool for establishing new physics beyond SM. New physics effects show themselves in rare $`B`$ meson decays in two different ways, namely, through the Wilson coefficients which could be distinctly different from their SM counterparts or through the new structures in the effective Hamiltonian (see for example ). One of the main goal of the working $`B`$ factories and of the future hadron colliders is the study of the rare $`B`$ meson decays. Therefore theoretical and experimental investigation of the rare decays of $`B`$ mesons receive special attention. The aim of this paper to investigate the rare decays $`\overline{B}K^{}\pi ^+\mathrm{}^+\mathrm{}^{}`$ and $`\overline{B}\pi ^+\pi ^{}\mathrm{}^+\mathrm{}^{}`$ when $`K\pi `$ and $`\pi \pi `$ systems are the decay products of $`K^{}`$ and $`\rho `$ mesons, in a model–independent manner. The inclusive $`bs\mathrm{}^+\mathrm{}^{}`$ and exclusive $`BK^{}\mathrm{}^+\mathrm{}^{}`$ decays were analyzed in a model–independent way in and , respectively. The present work is an extension of the previous studies of the $`BK^{}(\rho )\mathrm{}^+\mathrm{}^{}`$ decay. The distribution in $`\mathrm{cos}\theta _P`$, where $`\theta _P`$ is the angle of the $`K^{}(\pi ^{})`$ in the $`K^{}\pi ^+(\pi ^{}\pi ^+)`$ center of mass frame, and the dependence on the azimuthal angle $`\phi `$ between the $`\mathrm{}^+\mathrm{}^{}`$ and $`K^{}\pi ^+`$ or $`\pi ^{}\pi ^+`$ planes, which does not exist in the $`BK^{}(\rho )\mathrm{}^+\mathrm{}^{}`$ decay, can provide additional new information. This information is sensitive to the polarization state of the vector meson $`K^{}(\rho )`$ and therefore opens new possibility in probing the structure of the effective Hamiltonian. Angular distributions and CP asymmetries in the $`\overline{B}K^{}\pi ^+\mathrm{}^+\mathrm{}^{}`$ and $`\overline{B}\pi ^+\pi ^{}\mathrm{}^+\mathrm{}^{}`$ decays in the SM were studied in . The paper is organized as follows. In section 2, we present the general, model independent expression of the decay distribution for the $`BK\pi (\pi \pi )\mathrm{}^+\mathrm{}^{}`$ decay, in terms of the helicity amplitudes, including the non–zero lepton mass effects. In section 3, we present our numerical analysis for different angular distribution together with a brief concluding remark of our results. ## 2 Theoretical background The matrix element for the $`BK^{}(\rho )\mathrm{}^+\mathrm{}^{}`$ decay at quark level is described by the $`bf\mathrm{}^+\mathrm{}^{}(f=s,d)`$ transition. Following the works , the matrix elements of the $`bf\mathrm{}^+\mathrm{}^{}`$ decay can be written as the sum of the SM and new physics contributions as $`_1`$ $`=`$ $`{\displaystyle \frac{G\alpha }{\sqrt{2}\pi }}V_{tb}V_{tf}^{}\{C_{SL}\overline{f}i\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{q^2}}Lb\overline{\mathrm{}}\gamma _\mu \mathrm{}+C_{BR}\overline{f}i\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{q^2}}Rb\overline{\mathrm{}}\gamma _\mu \mathrm{}`$ (1) $`+`$ $`C_{LL}^{tot}\overline{f}_L\gamma _\mu b_L\overline{\mathrm{}}_L\gamma ^\mu \mathrm{}_L+C_{LR}^{tot}\overline{f}_L\gamma _\mu b_L\overline{\mathrm{}}_R\gamma ^\mu \mathrm{}_R+C_{RL}\overline{f}_R\gamma _\mu b_R\overline{\mathrm{}}_L\gamma ^\mu \mathrm{}_L`$ $`+`$ $`C_{RR}\overline{f}_R\gamma _\mu b_R\overline{\mathrm{}}_R\gamma ^\mu \mathrm{}_R+C_{LRLR}\overline{f}_Lb_R\overline{\mathrm{}}_L\mathrm{}_R+C_{RLLR}\overline{f}_Rb_L\overline{\mathrm{}}_L\mathrm{}_R`$ $`+`$ $`C_{LRRL}\overline{f}_Lb_R\overline{\mathrm{}}_R\mathrm{}_L+C_{RLRL}\overline{f}_Rb_L\overline{\mathrm{}}_R\mathrm{}_L+C_T\overline{f}\sigma _{\mu \nu }b\overline{\mathrm{}}\sigma ^{\mu \nu }\mathrm{}`$ $`+`$ $`iC_{TE}ϵ^{\mu \nu \alpha \beta }\overline{f}\sigma _{\mu \nu }b\overline{\mathrm{}}\sigma _{\alpha \beta }\mathrm{}\},`$ the first four terms contain the SM contribution too, while all the others describe the new physics effects and where $`L(R)={\displaystyle \frac{1\gamma _5}{2}}\left({\displaystyle \frac{1+\gamma _5}{2}}\right).`$ In further analysis we will neglect the tensor type interactions (terms $`C_T`$ and $`C_{TE}`$ since physically measurable quantities are not sensitive to the presence of tensor type interactions, as is shown in ). The matrix elements for the precess $`BPP^{}\mathrm{}^+\mathrm{}^{}`$ can be obtained from the matrix element $`BV\mathrm{}^+\mathrm{}^{}(V=K^{},\rho )`$ in the following way $`=_{BV\mathrm{}^+\mathrm{}^{}}\mathrm{\Pi }(s_M)\pi ^+(p_+)\pi ^{}(p_{})\|V(p_V,\lambda ),`$ (2) where we assume that the resonance contribution of the intermediate vector meson can be implemented by the Breit–Wigner form $`\mathrm{\Pi }(s_M)={\displaystyle \frac{\sqrt{m_V\mathrm{\Gamma }_V/\pi }}{s_Mm_V^2+im_V\mathrm{\Gamma }_V}},`$ (3) where $`s_M=(p_++p_{})^2`$ and $`m_V`$ and $`\mathrm{\Gamma }_V`$ are the mass and width of the vector meson ($`K^{},\rho `$). For the decay part of the vector meson we use $`P(p_+)P^{}(p_{})\|V(p_V,\lambda )=\sqrt{BR}Y_{\lambda _{max}}^{\lambda _i}(\theta _P,\phi ),`$ (4) where $`Y_{\mathrm{}}^m(\theta _P,\phi )`$ are the $`J=\mathrm{}`$ spherical harmonics, and the angles $`\theta _P`$ and $`\phi `$ belong to those of the final $`P`$ meson in the vector meson’s rest frame. The coupling of $`VPP^{}`$ decay is effectively taken into account by the branching ratios. From (2) we see that in order to calculate the matrix element $`BPP^{}\mathrm{}^+\mathrm{}^{}`$, the matrix element of the $`BV\mathrm{}^+\mathrm{}^{}`$ decay is needed in the first hand, for which the matrix elements $`V\left\|\overline{f}\gamma _\mu (1\pm \gamma _5)b\right\|B(p_B),V\left\|\overline{f}i\sigma _{\mu \nu }(1+\gamma _5)b\right\|B(p_B)\text{and}V\left\|\overline{f}(1\pm \gamma _5)b\right\|B(p_B),`$ (5) need to be calculated, whose most general forms can be written as $`V(p_V,\epsilon )\left\|\overline{f}\gamma _\mu (1\pm \gamma _5)b\right\|B(p_B)=`$ $`ϵ_{\mu \nu \lambda \sigma }\epsilon ^\nu p_V^\lambda q^\sigma {\displaystyle \frac{2V(q^2)}{m_B+m_V}}\pm i\epsilon _\mu ^{}(m_B+m_V)A_1(q^2)i(p_B+p_V)_\mu (\epsilon ^{}q){\displaystyle \frac{A_2(q^2)}{m_B+m_V}}`$ $`iq_\mu {\displaystyle \frac{2m_V}{q^2}}(\epsilon ^{}q)\left[A_3(q^2)A_0(q^2)\right],`$ $`V(p_V,\epsilon )\left\|\overline{f}i\sigma _{\mu \nu }q^\nu (1\pm \gamma _5)b\right\|B(p_B)=`$ (7) $`4ϵ_{\mu \nu \lambda \sigma }\epsilon ^\nu p_V^\lambda q^\sigma T_1(q^2)\pm 2i\left[\epsilon _\mu ^{}(m_B^2s_M)(p_B+p_V)_\mu (\epsilon ^{}q)\right]T_2(q^2)`$ $`\pm 2i(\epsilon ^{}q)\left[q_\mu (p_B+p_V)_\mu {\displaystyle \frac{q^2}{m_B^2s_M}}\right]T_3(q^2),`$ where $`q=p_Bp_V`$ is the momentum transfer and $`\epsilon `$ is the polarization vector of the vector meson. In order to ensure finiteness of (2) at $`q^2=0`$, we demand that $`A_3(q^2=0)=A_0(q^2=0)`$. The matrix element $`V\left\|\overline{f}(1\pm \gamma _5)b\right\|B`$ can easily be calculated from Eq. (2). For this aim it is enough to contract both sides of Eq. (2) with $`q_\mu `$ and use equation of motion. Taking mass of the strange quark to be zero, it gives $`V(p_V,\epsilon )\left\|\overline{f}(1\pm \gamma _5)b\right\|B(p_B)={\displaystyle \frac{1}{m_b}}\{i(\epsilon ^{}q)(m_B+m_V)A_1(q^2)`$ (8) $`\pm i(m_B^2s_M)(\epsilon ^{}q){\displaystyle \frac{A_2(q^2)}{m_B+m_V}}\pm 2im_V(\epsilon ^{}q)[A_3(q^2)A_0(q^2)]\}.`$ Furthermore, using the equation of motion, the form factor $`A_3`$ can be expressed in terms of the form factors $`A_1`$ and $`A_2`$ (see ) $`A_3(q^2)={\displaystyle \frac{m_B+m_V}{2m_V}}A_1(q^2){\displaystyle \frac{m_Bm_V}{2m_V}}A_2(q^2).`$ (9) Using this relation, the matrix element (8) can be written in the following form $`V(p_V,\epsilon )\left\|\overline{f}(1\pm \gamma _5)b\right\|B(p_B)={\displaystyle \frac{1}{m_b}}\left\{2im_V(\epsilon ^{}q)A_0(q^2)\right\}.`$ (10) As a result of the above considerations, the matrix element of the $`BV\mathrm{}^+\mathrm{}^{}`$ decay can be determined straightforwardly $`(BV\mathrm{}^+\mathrm{}^{})={\displaystyle \frac{G\alpha }{4\sqrt{2}\pi }}V_{tb}V_{tf}^{}`$ (11) $`\times \{\overline{\mathrm{}}\gamma _\mu (1\gamma _5)\mathrm{}[2𝒱_{L_1}ϵ_{\mu \nu \rho \sigma }\epsilon ^\nu p_V^Vq^\sigma i𝒱_{L_2}\epsilon _\mu ^{}+i𝒱_{L_3}(\epsilon ^{}q)(p_B+p_V)_\mu +i𝒱_{L_4}(\epsilon ^{}q)q_\mu ]`$ $`+\overline{\mathrm{}}\gamma _\mu (1+\gamma _5)\mathrm{}\left[2𝒱_{R_1}ϵ_{\mu \nu \rho \sigma }\epsilon ^\nu p_V^\rho q^\sigma i𝒱_{R_2}\epsilon _\mu ^{}+i𝒱_{R_3}(\epsilon ^{}q)(p_B+p_V)_\mu +i𝒱_{R_4}(\epsilon ^{}q)q_\mu \right]`$ $`+\overline{\mathrm{}}(1\gamma _5)\mathrm{}\left[i𝒮_L(\epsilon ^{}q)\right]+\overline{\mathrm{}}(1+\gamma _5)\mathrm{}\left[i𝒮_R(\epsilon ^{}q)\right]\},`$ where $`𝒱_{L_i}`$ and $`𝒱_{R_i}`$ are the coefficients of left– and right–handed leptonic currents with vector structure, and $`𝒮_{L,R}`$ are the weights of scalar leptonic currents with respective chirality, respectively, whose explicit forms are given as $`𝒱_{L_1}`$ $`=`$ $`(C_{LL}^{tot}+C_{RL}){\displaystyle \frac{V(q^2)}{m_B+m_V}}2(C_{BR}+C_{SL}){\displaystyle \frac{T_1}{q^2}},`$ $`𝒱_{L_2}`$ $`=`$ $`(C_{LL}^{tot}C_{RL})(m_B+m_V)A_12(C_{BR}C_{SL}){\displaystyle \frac{T_2}{q^2}}(m_B^2s_M),`$ $`𝒱_{L_3}`$ $`=`$ $`{\displaystyle \frac{C_{LL}^{tot}C_{RL}}{m_B+m_V}}A_22(C_{BR}C_{SL}){\displaystyle \frac{1}{q^2}}\left[T_2+{\displaystyle \frac{q^2}{m_B^2s_M}}T_3\right],`$ $`𝒱_{L_4}`$ $`=`$ $`(C_{LL}^{tot}C_{RL}){\displaystyle \frac{2m_V}{q^2}}(A_3A_0)+2(C_{BR}C_{SL}){\displaystyle \frac{T_3}{q^2}},`$ $`𝒱_{R_1}`$ $`=`$ $`𝒱_{L_1}(C_{LL}^{tot}C_{LR}^{tot},C_{RL}C_{RR}),`$ $`𝒱_{R_2}`$ $`=`$ $`𝒱_{L_2}(C_{LL}^{tot}C_{LR}^{tot},C_{RL}C_{RR}),`$ $`𝒱_{R_3}`$ $`=`$ $`𝒱_{L_3}(C_{LL}^{tot}C_{LR}^{tot},C_{RL}C_{RR}),`$ $`𝒱_{R_4}`$ $`=`$ $`𝒱_{L_4}(C_{LL}^{tot}C_{LR}^{tot},C_{RL}C_{RR}),`$ $`𝒮_L`$ $`=`$ $`(C_{LRRL}C_{RLRL})\left({\displaystyle \frac{2m_V}{m_b}}A_0\right),`$ $`𝒮_R`$ $`=`$ $`(C_{LRLR}C_{RLLR})\left({\displaystyle \frac{2m_V}{m_b}}A_0\right),`$ In order to obtain the full helicity amplitude of the $`BPP^{}\mathrm{}^+\mathrm{}^{}`$ which follows from Eq. (2), the helicity amplitude of the $`BV\mathrm{}^+\mathrm{}^{}`$ decay must be written, which we denote as $`_\lambda ^{\lambda _{\mathrm{}}\overline{\lambda }_{\mathrm{}}}`$ $`_{\lambda _i}^{\lambda _{\mathrm{}}\overline{\lambda }_{\mathrm{}}}={\displaystyle \underset{\lambda _V^{}}{}}\eta _{\lambda _V^{}}L_{\lambda _V^{}}^{\lambda _{\mathrm{}}\overline{\lambda }_{\mathrm{}}}H_{\lambda _V^{}}^{\lambda _i},`$ (12) where (15) where $`\epsilon _V^{}`$ is the polarization vector of the virtual intermediate vector boson satisfying the relation $`g^{\mu \nu }={\displaystyle \underset{\lambda _V^{}}{}}\eta _{\lambda _V^{}}\epsilon _{\lambda _V^{}}^\mu \epsilon _{\lambda _V^{}}^\nu ,`$ where the summation is over the helicities $`\lambda _V^{}=\pm 1,0,s`$ of the virtual intermediate vector meson, with the metric defined as $`\eta _\pm =\eta _0=\eta _s=1`$ (for more detail see ). In Eq. (15), $`J_\mu ^{\mathrm{}}`$ and $`J_\mu ^i`$ are the leptonic and hadronic currents, respectively. Using Eqs. (11)–(15) for the helicity amplitudes $`_{\lambda _i}^{\lambda _{\mathrm{}}\overline{\lambda }_{\mathrm{}}}`$ we get the following expressions (24) where (38) where superscripts denote helicities of the lepton and antilepton and subscripts correspond to the helicity of the vector meson (in our case $`\rho `$ or $`K^{}`$ meson), respectively. (49) where $`\theta `$ is the polar angle of positron in the rest frame of the intermediate boson with respect to its helicity axis. Note that we take $`p_V^2=s_M`$, but not $`m_V^2`$, in order to take into account $`V`$’s being a virtual particle which subsequently decays into $`\pi ^+\pi ^{}`$ or $`K^{}\pi ^+`$ pair. Remembering that the existing CLEO result for the $`BX_s\gamma `$ and $`BK^{}\gamma `$ decays impose strong constraint on the parameter space $`C_{BR}`$ and $`C_{SL}`$. For this reason here we assume that they are equal to each other in the SM. Hence we will take $`C_{LL}^{tot}`$ $`=`$ $`C_9^{eff}C_{10}+C_{LL},`$ $`C_{LR}^{tot}`$ $`=`$ $`C_9^{eff}+C_{10}+C_{LR}.`$ Using the expressions of the helicity amplitudes for the differential decay rate width of the $`BV(PP^{})\mathrm{}^+\mathrm{}^{}`$ decay, we get (56) where (71) where $`\theta _P`$ is the polar angle of the pseudoscalar $`P`$ meson momentum in the rest frame of the vector meson, with respect to the helicity axis, i.e., the outgoing direction of $`V`$ meson, and $`\phi `$ is the azimuthal angle between the planes of the two decays $`VPP^{}`$ and $`V^{}\mathrm{}^+\mathrm{}^{}`$. Kinematically allowed region of the variables are given as (77) We note that, in further analysis the narrow–width approximation for $`V`$ meson will be used, i.e., $`\underset{\mathrm{\Gamma }_V0}{lim}{\displaystyle \frac{\mathrm{\Gamma }_Vm_V/\pi }{(s_Mm_V^2)^2+m_V^2\mathrm{\Gamma }_V^2}}=\delta (s_Mm_V^2),`$ by means of which integration of Eq. (56) over $`s_M`$ can easily be carried and the differential decay rate with respect to dilepton mass $`q^2`$, azimuthal angle $`\phi `$, polar angles $`\theta _P`$ and $`\theta `$ can be written as (84) for which we will use the experimental results for the $`VPP^{}`$ namely, $`(\rho \pi ^+\pi ^{})=(K^{}K\pi )=1`$. It should be noted here that in addition to the variables that exist in $`BV\mathrm{}^+\mathrm{}^{}`$ decay, there appears a new variable $`\theta _P`$, and contrary to the $`BV\mathrm{}^+\mathrm{}^{}`$ case, the dependence of the cascade decay $`BV(PP^{})\mathrm{}^+\mathrm{}^{}`$ on the azimuthal angle $`\phi `$ is not trivial. Incidentally, we would like to remind the reader that, if Eq. (84) is integrated over $`\theta _P`$ and $`\phi `$, the differential decay rate for the $`BV\mathrm{}^+\mathrm{}^{}`$ decay is obtained. The model independent analysis of the $`BK^{}\mathrm{}^+\mathrm{}^{}`$ decay is presented in , in which the dependence of the experimentally measured quantities, such as branching ratio, forward–backward asymmetry and longitudinal polarization of the final lepton and the ratio $`\mathrm{\Gamma }_L/\mathrm{\Gamma }_T`$ of the decay widths when $`K^{}`$ meson is longitudinally and transversally polarized, on the new Wilson coefficients, are systematically studied. As has been noted already, our main goal in this work is to investigate the dependence of such angular distributions on the new Wilson coefficients in the $`BV(VPP^{})\mathrm{}^+\mathrm{}^{}`$ decay, which does not exist in the $`BV\mathrm{}^+\mathrm{}^{}`$ decay and not studied in . It can easily be seen from Eq. (84) that the cascade decay $`BV(PP^{})\mathrm{}^+\mathrm{}^{}`$ has a rich angular structure. Therefore, in the light of this observation, a thorough investigation of different distributions will prove useful in separating various angular coefficients experimentally. Along the lines as suggested by , we adopt two different strategies in further analysis of the problem under consideration, namely investigation of the various individual angular distributions and asymmetries and their relation to the new Wilson coefficients. For the purpose of studying single angle distributions, we integrate Eq. (84) over $`q^2`$, $`\theta `$ and $`\phi `$, which takes the form $`{\displaystyle \frac{d\mathrm{\Gamma }}{d\mathrm{cos}\theta _P}}2\mathrm{cos}^2\theta _P[{\displaystyle \frac{2}{3}}\stackrel{~}{N}_1+{\displaystyle \frac{4}{3}}\stackrel{~}{N}_2+2\stackrel{~}{N}_4]+\mathrm{sin}^2\theta _P[{\displaystyle \frac{4}{3}}\stackrel{~}{N}_5++{\displaystyle \frac{8}{3}}\stackrel{~}{N}_6],`$ (85) where we introduce the notation $`\stackrel{~}{N}_i=N_i𝑑q^2`$. Defining an asymmetry parameter $`\alpha _{\theta _P}`$, from the angular distribution $`W(\mathrm{cos}\theta _P)=1+\alpha _{\theta _P}\mathrm{cos}^2\theta _P`$ we get $`\alpha _{\theta _P}={\displaystyle \frac{\stackrel{~}{N}_1+2\stackrel{~}{N}_2+3\stackrel{~}{N}_4}{\stackrel{~}{N}_5+2\stackrel{~}{N}_6}}1.`$ (86) Integrating Eq. (84) over $`q^2`$, $`\theta _P`$ and $`\phi `$, for the polar $`\mathrm{cos}\theta `$ distribution we get from which one can write the angular distribution $`W(\mathrm{cos}\theta )=1+\alpha _\theta \mathrm{cos}\theta +\beta _\theta \mathrm{cos}^2\theta `$, with the asymmetry parameters $`\alpha _\theta `$ and $`\beta _\theta `$ being defined as (91) Finally we consider the azimuthal angle $`\phi `$ distribution, which is obtained by integrating Eq. (84) over the parameters $`q^2`$, $`\theta _P`$ and $`\theta `$ to yield and the azimuthal angle $`\phi `$ asymmetry parameters $`\alpha _\phi `$ and $`\alpha _\phi `$ can be extracted from $`W(\phi )=1+\alpha _\phi \mathrm{sin}(2\phi )+\beta _\phi \mathrm{cos}(2\phi )`$ to give (96) A second strategy for separating various angular coefficients experimentally, is to define suitable asymmetry ratios that project out the partial rates from Eq. (84). For this purpose we consider the following asymmetries (see also ) $`A_\phi `$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Gamma }(\phi )d\mathrm{\Gamma }(\phi +\pi /2)+d\mathrm{\Gamma }(\phi +\pi )d\mathrm{\Gamma }(\phi +3\pi /2)}{d\mathrm{\Gamma }(\phi )+d\mathrm{\Gamma }(\phi +\pi /2)+d\mathrm{\Gamma }(\phi +\pi )+d\mathrm{\Gamma }(\phi +3\pi /2)}},`$ $`{\displaystyle \frac{\pi }{4}}\phi {\displaystyle \frac{\pi }{4}},`$ $`A_1`$ $`=`$ $`{\displaystyle \frac{N}{D}},`$ (98) where $`N`$ $`=`$ $`d\mathrm{\Gamma }(\theta ,\theta _P,\phi )d\mathrm{\Gamma }(\theta ,\theta _P,\phi +\pi )d\mathrm{\Gamma }(\theta ,\pi \theta _P,\phi )+d\mathrm{\Gamma }(\theta ,\pi \theta _P,\phi +\pi )`$ $``$ $`d\mathrm{\Gamma }(\pi \theta ,\theta _P,\phi )+d\mathrm{\Gamma }(\pi \theta ,\theta _P,\phi +\pi )+d\mathrm{\Gamma }(\pi \theta ,\pi \theta _P,\phi )`$ $``$ $`d\mathrm{\Gamma }(\pi \theta ,\pi \theta _P,\phi +\pi ),`$ $`0\theta _P\pi /2,`$ $`{\displaystyle \frac{\pi }{2}}\theta \pi ,`$ $`{\displaystyle \frac{\pi }{2}}\phi {\displaystyle \frac{\pi }{2}},`$ and the denominator $`D`$ is given by the same expression, with plus signs everywhere, $`A_2`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Gamma }(\theta _P,\phi )d\mathrm{\Gamma }(\theta _P,\phi +\pi )d\mathrm{\Gamma }(\pi \theta _P,\phi )+d\mathrm{\Gamma }(\pi \theta _P,\phi +\pi )}{d\mathrm{\Gamma }(\theta _P,\phi )+d\mathrm{\Gamma }(\theta _P,\phi +\pi )+d\mathrm{\Gamma }(\pi \theta _P,\phi )+d\mathrm{\Gamma }(\pi \theta _P,\phi +\pi )}},`$ $`0\theta _P\pi /2,`$ $`{\displaystyle \frac{\pi }{2}}\phi \pi /2.`$ In expressions (2)–(2) the angles that do not appear in the arguments of the differential rate $`d\mathrm{\Gamma }`$ have been integrated out over their physical ranges $`(0\theta ,\theta _P\pi ,0\phi 2\pi )`$. Integrating over all variables, we are left with the expressions $`\stackrel{~}{A}_\phi ,\stackrel{~}{A}_1,\stackrel{~}{A}_2`$, which depend only on the Wilson coefficients, as follows (here $`\stackrel{~}{}`$ in the notation refers to integration’s being performed over all variables) $`\stackrel{~}{A}_\phi `$ $`=`$ $`{\displaystyle \frac{8\text{Re}[\stackrel{~}{N}_9]}{\pi \left(\stackrel{~}{N}_1+2\stackrel{~}{N}_2+3\stackrel{~}{N}_4+2\stackrel{~}{N}_5+4\stackrel{~}{N}_6\right)}},`$ (100) $`\stackrel{~}{A}_1`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\text{Re}[\stackrel{~}{N}_{10}]}{\pi \left(\stackrel{~}{N}_1+2\stackrel{~}{N}_2+3\stackrel{~}{N}_4+2\stackrel{~}{N}_5+4\stackrel{~}{N}_6\right)}},`$ (101) $`\stackrel{~}{A}_2`$ $`=`$ $`{\displaystyle \frac{3\text{Re}[\stackrel{~}{N}_{11}]}{\sqrt{2}\left(\stackrel{~}{N}_1+2\stackrel{~}{N}_2+3\stackrel{~}{N}_4+2\stackrel{~}{N}_5+4\stackrel{~}{N}_6\right)}}.`$ (102) Before proceeding further, we would like to consider the CP–violating observables that can be constructed by combining the information on $`\overline{B}`$ and $`B`$ decays, namely, $`\overline{B}P\overline{P}^{}\mathrm{}^+\mathrm{}^{}`$ and $`B\overline{P}P^{}\mathrm{}^+\mathrm{}^{}`$, and define CP–odd asymmetry in the following way $`\stackrel{~}{A}_{CP}{\displaystyle \frac{\mathrm{\Gamma }\overline{\mathrm{\Gamma }}}{\mathrm{\Gamma }+\overline{\mathrm{\Gamma }}}},`$ (103) where $`\mathrm{\Gamma }`$ and $`\overline{\mathrm{\Gamma }}`$ are the decay widths of the $`\overline{B}P\overline{P}^{}\mathrm{}^+\mathrm{}^{}`$ and $`B\overline{P}P^{}\mathrm{}^+\mathrm{}^{}`$ processes, respectively. Explicit form of $`\mathrm{\Gamma }`$ can easily be obtained from Eq. (84), by performing integration over the variables $`q^2,\theta ,\theta _P,`$ and $`\phi `$. The decay width for the conjugate process can again be obtained from Eq. (84) by making the replacement $`N_i\overline{N}_i`$, where $`\overline{N}_i`$ are the functions for the conjugate processes. It should be noted that we consider a case in which all new Wilson coefficients are real. Furthermore, form factors which enter into Eq. (49) are computed in framework of light cone QCD sum rules method and this nonperturbative approach predicts that all form factors are real as well. In other words, there is no any new source for CP violation other than that are present in the SM. As is well known, for the $`B\rho \mathrm{}^+\mathrm{}^{}`$ decay in the SM only the coefficient $`C_9^{eff}`$ contains both a weak phase $`\phi _W`$ (associated with the imaginary part of the CKM matrix element) and a strong phase $`\delta _S`$ (attributed to the imaginary parts of the $`c\overline{c}`$ and $`b\overline{b}`$ loops). As a result of these considerations, it follows then that the decay width for the conjugate $`B\overline{P}P^{}\mathrm{}^+\mathrm{}^{}`$ process can be obtained from $`\overline{B}P\overline{P}^{}\mathrm{}^+\mathrm{}^{}`$ channel by the replacements $`\phi _W\phi _W`$ and $`\delta _S\delta _S`$. ## 3 Numerical analysis In this section we present our numerical results for the asymmetries $`A_{CP},\alpha _{\theta _P},\alpha _\theta ,\beta _\theta `$, $`\beta _\phi ,A_\phi ,\alpha _\phi ,\stackrel{~}{A}_1,\stackrel{~}{A}_\phi `$ and $`\stackrel{~}{A}_2`$ for the exclusive rare $`B\pi ^+\pi ^{}\mathrm{}^+\mathrm{}^{}`$ decay only. We take hadronic form factors from Table I and the Wilson coefficients from Table II. The values of the main input parameters used in our analysis are: $`m_b=4.8GeV,m_c=1.35GeV,m_\rho =0.77GeV,m_\tau =1.78GeV,m_\mu =0.105GeV`$ and $`m_B=5.28GeV`$. Here we note that the results for $`BK\pi \mathrm{}^+\mathrm{}^{}`$ can be obtained from $`B\pi \pi \mathrm{}^+\mathrm{}^{}`$ by replacing $`V_{td}`$ with $`V_{ts}`$, $`B\rho `$ transition form factors with $`BK`$ transition form factors, and $`m_\pi `$ with $`m_K`$. In the present work we choose light cone QCD sum rules method predictions for the form factors. In our numerical analysis we will use the results of the work in which the form factors are described by a three–parameter fit where the radiative corrections up to leading twist contribution and SU(3)–breaking effects are taken into account. The $`q^2`$–dependence of the form factors, which appears in our analysis could be parametrized as $`F(s)={\displaystyle \frac{F(0)}{1a_Fs+b_Fs^2}},`$ where $`s=q^2/m_B^2`$ is the dilepton invariant mass in units of $`B`$ meson mass, and the parameters $`F(0)`$, $`a_F`$ and $`b_F`$ are listed in Table 1 for each form factor. In the SM the Wilson coefficients $`C_7^{eff}(m_b)`$ and $`C_{10}(m_b)`$, whose analytical expressions are given in , are strictly real as can be read off from Table 2. In the leading logarithmic approximation, at the scale $`𝒪(\mu =m_b)`$, we have $`C_7^{eff}(m_b)`$ $`=`$ $`0.313,`$ $`C_{10}^{eff}(m_b)`$ $`=`$ $`4.669.`$ Although individual Wilson coefficients at $`\mu m_b`$ level are all real (see Table 2), the effective Wilson coefficient $`C_9^{eff}(m_b,\widehat{s})`$ has a finite phase, and in next–to–leading order $`C_9^{eff}(m_b,\widehat{s})=C_9(m_b)\left[1+{\displaystyle \frac{\alpha _s(\mu )}{\pi }}\omega (\widehat{s})\right]+Y_{SD}(m_b,\widehat{s})+Y_{LD}(m_B,s),`$ (104) where $`C_9(m_b)=4.344`$. Here $`\omega \left(\widehat{s}\right)`$ represents the $`𝒪(\alpha _s)`$ corrections coming from one–gluon exchange in the matrix element of the corresponding operator, whose explicit form can be found in . In (104) $`Y_{SD}`$ and $`Y_{LD}`$ represent, respectively, the short– and long–distance contributions of the four–quark operators $`𝒪_{i=1,\mathrm{},6}`$ . Here $`Y_{SD}`$ can be obtained by a perturbative calculation $`Y_{SD}(m_b,\widehat{s})`$ $`=`$ $`g(\widehat{m}_c,\widehat{s})C^{(0)}{\displaystyle \frac{1}{2}}g(1,\widehat{s})\left[4C_3+4C_4+3C_5+C_6\right]`$ $``$ $`{\displaystyle \frac{1}{2}}g(0,\widehat{s})\left[C_3+3C_4\right]+{\displaystyle \frac{2}{9}}\left[3C_3+C_4+3C_5+C_6\right]`$ $``$ $`\lambda _u\left[3C_1+C_2\right]\left[g(0,\widehat{s})g(\widehat{m}_c,\widehat{s})\right],`$ where $`C^{(0)}`$ $`=`$ $`3C_1+C_2+3C_3+C_4+3C_5+C_6,`$ $`\lambda _u`$ $`=`$ $`{\displaystyle \frac{V_{ub}V_{ud}^{}}{V_{tb}V_{td}^{}}},`$ and the loop function $`g(m_q,s)`$ stands for the loops of quarks with mass $`m_q`$ at the dilepton invariant mass $`s`$. This function develops absorptive parts for dilepton energies $`s=4m_q^2`$: $`g(\widehat{m}_q,\widehat{s})={\displaystyle \frac{8}{9}}\mathrm{ln}\widehat{m}_q+{\displaystyle \frac{8}{27}}+{\displaystyle \frac{4}{9}}y_q{\displaystyle \frac{2}{9}}\left(2+y_q\right)\sqrt{\left\|1y_q\right\|}`$ $`\times \left[\mathrm{\Theta }(1y_q)\left(\mathrm{ln}{\displaystyle \frac{1+\sqrt{1y_q}}{1\sqrt{1y_q}}}i\pi \right)+\mathrm{\Theta }(y_q1)\mathrm{\hspace{0.17em}2}\mathrm{arctan}{\displaystyle \frac{1}{\sqrt{y_q1}}}\right],`$ where $`\widehat{m}_q=m_q/m_b`$ and $`y_q=4\widehat{m}_q^2/\widehat{s}`$. In addition to these perturbative contributions, the $`\overline{c}c`$ loops can excite low–lying charmonium states $`\psi (1s),\mathrm{},\psi (6s)`$ whose contributions are represented by $`Y_{LD}`$ : $`Y_{LD}(m_b,\widehat{s})`$ $`=`$ $`{\displaystyle \frac{3}{\alpha ^2}}\left[{\displaystyle \frac{V_{cf}^{}V_{cb}}{V_{tf}^{}V_{tb}}}C^{(0)}{\displaystyle \frac{V_{uf}^{}V_{ub}}{V_{tf}^{}V_{tb}}}\left(3C_3+C_4+3C_5+C_6\right)\right]`$ $`\times `$ $`{\displaystyle \underset{V_i=\psi \left(1s\right),\mathrm{},\psi \left(6s\right)}{}}{\displaystyle \frac{\pi \kappa _i\mathrm{\Gamma }\left(V_i\mathrm{}^+\mathrm{}^{}\right)M_{V_i}}{\left(M_{V_i}^2\widehat{s}m_b^2iM_{V_i}\mathrm{\Gamma }_{V_i}\right)}},`$ where $`\kappa _i`$ are the Fudge factors (see for example ). Let us first study the dependence of the asymmetry parameter $`\alpha _{\theta _P}`$ on the new Wilson coefficients. Note that in further analysis, only short distance contributions are taken into account and integration over $`q^2`$ is performed in the full physical region $`4m_{\mathrm{}}^2q^2(m_Bm_V^2)^2`$. We assumed that all new Wilson coefficients $`C_X`$ are real, i.e., we do not introduce any new phase in addition to the one present in the SM. In Figs. (1) and (2), we present the dependence of $`\alpha _{\theta _P}`$ on the new Wilson coefficients, for the $`B\pi ^+\pi ^{}e^+e^{}`$ and $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decays, respectively. Here and in all of the following figures, zero value of new Wilson coefficients $`C_X`$ correspond to the SM prediction. In the case of $`B\pi ^+\pi ^{}e^+e^{}`$ decay the asymmetry parameter $`\alpha _{\theta _P}`$ is more sensitive to $`C_{LL}^{tot}`$ and $`C_{RL}`$, while for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay it depends strongly on $`C_{RL}`$. These dependencies can be explained as follows. For the $`B\pi ^+\pi ^{}e^+e^{}`$ decay, if the terms proportional to electron mass are neglected, it easily be seen from Eq. (86) that $`\alpha _{\theta _P}={\displaystyle \frac{\stackrel{~}{N}_2}{\stackrel{~}{N}_6}}1.`$ In the limit $`v1`$ we get (107) In the SM in the large dilepton mass region, say about $`q^25GeV^2`$, $`C_9^{eff}+C_{10}0.4`$ and Re$`\left[C_9^{eff}C_{10}\right]9.5`$. It follows then that the interference terms between $`C_9^{eff}C_{10}`$ and $`C_{LL}(C_{RL})`$ are dominant and hence contributions coming from $`C_{LL}(C_{RL})`$ are large. These figures illustrates that the contributions of $`C_{LL}`$ and $`C_{RL}`$ to $`\alpha _{\theta _P}`$ is positive for $`C_{LL}>0`$ and $`C_{RL}<0`$, and negative for $`C_{LL}<0`$ and $`C_{RL}>0`$. It should be noted that the asymmetry parameter $`\alpha _{\theta _P}`$ can get only positive or negative values for the case $`C_{LL}0`$ and $`C_{RL}0`$, while it is always positive for all other choices of the Wilson coefficients, as is the case in the SM. For this reason determination of the sign of $`\alpha _{\theta _P}`$ can serve as an efficient tool for establishing new physics. In the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ process however, the situation is slightly different compared to that of the $`B\pi ^+\pi ^{}e^+e^{}`$ transition. Largest contribution in this case comes from $`C_{RL}`$ and contributions from all other Wilson coefficients are comparable to one another. This observation can be attributed to mass of the $`\tau `$ lepton, for which $`\alpha _{\theta _P}`$ is positive for the choice of each individual new Wilson coefficients. Depicted in Figs. (3), (4) and Figs. (5), (6) are the dependencies of the asymmetry parameters $`\alpha _\theta `$ and $`\beta _\theta `$ on the new Wilson coefficients $`C_X`$, for the $`B\pi ^+\pi ^{}e^+e^{}`$ and $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decays, respectively. Figs. (3) and (5) depict that the asymmetry parameter $`\alpha _\theta `$ for the $`e^+e^{}`$ channel depends strongly on the new Wilson coefficient $`C_{RL}`$, while it displays similar behavior for all new Wilson coefficients for the $`\tau ^+\tau ^{}`$ case. We observe that the asymmetry parameter $`\beta _\theta `$ depends strongly on $`C_{RL}`$ in the $`e^+e^{}`$ channel and on $`C_{RL}`$ and $`C_{RR}`$ in the $`\tau ^+\tau ^{}`$ channel. It is interesting that $`\beta _\theta `$ changes its sign when $`C_{RL}>2`$ in the $`e^+e^{}`$ channel, while it is negative for all values of $`C_{RL}`$ or $`C_{RR}`$ in the $`\tau ^+\tau ^{}`$ channel. Therefore determination of the sign of $`\beta _\theta `$ is useful in looking for new physics. In Figs. (7) and (8) we present the dependence of the asymmetry parameter $`\beta _\phi `$ for the $`B\pi ^+\pi ^{}e^+e^{}`$ and $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decays, respectively. From these figures one notices that this asymmetry parameter is quite sensitive to the variation in $`C_{RL}`$ and changes its sign at $`C_{RR}>2`$ for both channels. Our investigation of the dependence of the asymmetry parameter $`\alpha _\phi `$ on the new Wilson coefficients shows that for the range $`4<C_X<4`$, $`\alpha _\phi `$ varies between $`6\times 10^3`$ to $`6\times 10^3`$ for the $`B\pi ^+\pi ^{}e^+e^{}`$ and $`5.0\times 10^3`$ to $`2.5\times 10^3`$ for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decays, repectively. Therefore detection of the dependence of the asymmetry parameter $`\alpha _\phi `$ on the new Wilsom coefficients is quite hard from experimental point of view. The asymmetry parameter $`\stackrel{~}{A}_1`$ shows strong dependence on $`C_{RL}`$ and $`C_{LL}`$ for the $`e^+e^{}`$ channel, as depicted in Fig. (9), whose contributions are dominant compared to the other Wilson coefficients. For the $`\tau ^+\tau ^{}`$ channel contributions of $`C_{RL}`$ and $`C_{RR}`$ become dominant, as can be seen in Fig. (10). The asymmetry parameter $`\stackrel{~}{A}_2`$ varies considerably for the $`e^+e^{}`$ channel, in relation to the variations occurring in $`C_{RL}`$, while this behavior is switched to the Wilson coefficient $`C_{LRRL}`$ for the $`\tau ^+\tau ^{}`$ case which are presented in Figs. (11) and (12). Presented in Figs. (13) and (14) are the dependence of the asymmetry parameter $`A_\phi `$ on new Wilson coefficients for the $`e^+e^{}\pi ^+\pi ^{}`$ and $`\tau ^+\tau ^{}\pi ^+\pi ^{}`$ decays, respectively. In both cases the asymmetry parameter $`A_\phi `$ shows strong dependence on $`C_{RL}`$. Finally, in Figs. (15) and (16) we present the dependence of the averaged (i.e., integrated over $`q^2`$ in the full physical region) CP asymmetry on the Wilson coefficients for the $`B\pi ^+\pi ^{}e^+e^{}`$ and $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decays, respectively. We observe that $`A_{CP}`$ is strongly dependent on $`C_{RL}`$ and $`C_{RR}`$ for the $`e^+e^{}`$ channel, while this dependence is switched to $`C_{LL}`$ for the $`\tau ^+\tau ^{}`$ channel. One can easily read from this figures that $`A_{CP}>2.0\times 10^2`$ in the $`e^+e^{}`$ channel, for the values $`C_{RR}2`$ and $`\left\|A_{CP}\right\|>1.7\times 10^2`$ when $`C_{LL}2`$. As the final concluding remark, we presented in this work the model independent analysis of the exclusive $`B\pi ^+\pi ^{}\mathrm{}^+\mathrm{}^{}`$ $`(\mathrm{}=e,\tau )`$ decay is presented. In particular, the sensitivity to the new Wilson coefficients of the experimentally measurable asymmetries and CP violating asymmetry are systematically analyzed. The main result of the present study is that different asymmetry parameters show strong dependence on different new Wilson coefficients. Therefore a combined analysis of the different asymmetries and CP violating asymmetry can give unambiguous information about the existence of new physics beyond the SM and especially about various new Wilson coefficients. ## Figure captions Fig. 1 The dependence of the asymmetry parameter $`\alpha _{\theta _P}`$ on the new Wilson coefficients for the $`B\pi ^+\pi ^{}e^+e^{}`$ decay. Fig. 2 The same as in Fig. (1), but for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay. Fig. 3 The same as in Fig. (1), but for the asymmetry parameter $`\alpha _\theta `$. Fig. 4 The same as in Fig. (1), but for the asymmetry parameter $`\beta _\theta `$. Fig. 5 The same as in Fig. (3), but for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay. Fig. 6 The same as in Fig. (4), but for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay. Fig. 7 The same as in Fig. (1), but for the asymmetry parameter $`\beta _\phi `$. Fig. 8 The same as in Fig. (7), but for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay. Fig. 9 The same as in Fig. (1), but for the asymmetry parameter $`\stackrel{~}{A}_1`$. Fig. 10 The same as in Fig. (9), but for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay. Fig. 11 The same as in Fig. (9), but for the asymmetry parameter $`\stackrel{~}{A}_2`$. Fig. 12 The same as in Fig. (11), but for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay. Fig. 13 The same as in Fig. (9), but for the asymmetry parameter $`\stackrel{~}{A}_\phi `$. Fig. 14 The same as in Fig. (13), but for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay. Fig. 15 The dependence of the averaged CP asymmetry on the new Wilson coefficients for the $`B\pi ^+\pi ^{}e^+e^{}`$ decay. Fig. 16 The same as in Fig. (15), but for the $`B\pi ^+\pi ^{}\tau ^+\tau ^{}`$ decay.
warning/0006/quant-ph0006014.html	ar5iv	text	# A Refutation of Bell’s Theorem1footnote 11footnote 1This work was presented at the international conference ”Foundations of Probability and Physics” held in Växjö, Sweden; Nov. 2000. ## I Introduction Bell’s TheoremJSB2 exhibits a peculiar discrepancy between any local realistic theory and Quantum Mechanics, which leads to empirically distinguishable alternatives. The quandary is that neither local realistic conceptions nor Quantum Mechanics are easy to abandon. Indeed, classical physics and common sense are usually based upon the former, while the latter is rightly presented as the most successful theory of all times. Several experiments have been done, all but a fewFS1 show violations of Bell inequalities. Yet, the ideas brought forth by Bell’s Theorem are so disconcerting that there is still incredulity, not to mention antipathy, evoked by the verdict. The purpose of this article is to provide a refutation of this theorem, within a strictly quantum theoretical framework, without the use of outside assumptions. Although experiments showing violations of Bell’s inequalities are getting increasingly accurate and loophole-freeASP1 , it must be stressed that they, no matter how accurate and close to the ideal, can prove no more than the validity of quantum mechanics, and not the validity of the theorem. Herein, it will be assumed that all tests conducted so far prove conclusively the validity of Quantum Mechanics. In other words, the purpose of this article is not to criticise the numerous experiments, or quantum mechanics for that matter, but Bell’s Theorem itself. ## II The EPRB gedanken experiment ### II.1 Spin observables and the singlet state Bell’s theorem is usually based on a didactic reformulation of the EPR (Einstein, Podolsky and RosenEPR1 ) gedanken experiment, due to D. BohmDB1 . In this EPRB gedanken experiment, a pair of spin-½ particles with total spin zero is produced such that each particle moves away from the source in opposite directions along the y-axis. Two Stern-Gerlach devices are placed at opposite points (left and right) on the y-axis, and are oriented respectively along the directions $`𝐮`$ and $`𝐯`$. The spin observable associated with a measurement given by a Stern-Gerlach device oriented along the unit vector $`𝒖`$ is $`𝝈𝐮`$ (to simplify notation, $`\mathrm{}/2`$ is set to $`1`$ throughout); where the components of $`𝝈`$ are then the Pauli matrices $`\sigma _x`$, $`\sigma _y`$, and $`\sigma _z`$. Let $`_\mathrm{L}`$ and $`_\mathrm{R}`$ be the Hilbert spaces associated with each Stern-Gerlach device respectively. The Hilbert space $``$ associated with the entire EPRB system is the direct product of the Hilbert spaces associated with each Stern-Gerlach device: $$_\mathrm{L}_\mathrm{R}.$$ (1) The spin observables relevant to $`_\mathrm{L}`$ and $`_\mathrm{R}`$ have their respective counterpart in this new product space $``$ as $`𝝈_\mathrm{L}𝐮𝝈𝐮1\mathrm{l}_\mathrm{R},`$ (2a) $`𝝈_\mathrm{R}𝐯1\mathrm{l}_\mathrm{L}𝝈𝐯,`$ (2b) where $`1\mathrm{l}_\mathrm{L}`$ and $`1\mathrm{l}_\mathrm{R}`$ are the identity operators of $`_\mathrm{L}`$ and $`_\mathrm{R}`$. Contrary to the observables $`𝝈𝐮`$ and $`𝝈𝐯`$ which are mutually non commuting when $`𝐮𝐯`$, these new observables $`𝝈_\mathrm{L}𝐮`$ and $`𝝈_\mathrm{R}𝐯`$ do commute, reflecting the fact that the Stern-Gerlach devices are arbitrarily far from each other, and are thus measuring distinct subsystems. The product of these two observables $$(𝝈_\mathrm{L}𝐮)(𝝈_\mathrm{R}𝐯)=𝝈𝐮𝝈𝐯$$ (3) is therefore also an observable and can be understood as a spin correlation observable corresponding to the joint spin measurement of both Stern-Gerlach devices. The product space $``$ is spanned by the product basis formed by the four eigenvectors $`\{\|++,\|+,\|+,\|\}`$ associated with the spin correlation observable $`(𝝈_\mathrm{L}𝐧)(𝝈_\mathrm{R}𝐧)`$ where $`𝐧`$ is a unitary vector. In an EPRB gedanken experiment, the source produces particle pairs with zero total spin, represented by the singlet state $$\|\psi =\frac{1}{\sqrt{2}}[\|+\|+].$$ (4) This singlet state has the important property of being invariant under rotation, which permits one to ignore the explicit form of $`𝐧`$ in expressing the $``$ basis (see, for instance, Ref. GHSZ1 ). ### II.2 Statistical properties of the singlet state As it is, nothing certain can be said either about a single spin measurement, or about a single spin correlation measurement, performed on a system represented by the singlet state. According to the Born interpretation of the state vector, only probabilistic predictions—such as expectation values relevant to numerous measurements in the same context—are allowed. It can be shown (see, for instance, Ref. ctdl1 , chapter IV), that the expectation value of an observable $`\widehat{A}`$ is $`\widehat{A}_\varphi =\varphi \|\widehat{A}\|\varphi `$ and therefore, with the help of Eqs. (2) and (4), that the expectation value of a spin observable for the singlet state $`\|\psi `$ is zero: $`𝝈_\mathrm{L}𝐮_\psi `$ $`=`$ $`\psi \|𝝈𝐮1\mathrm{l}_\mathrm{R}\|\psi =0,`$ (5a) $`𝝈_\mathrm{R}𝐯_\psi `$ $`=`$ $`\psi \|1\mathrm{l}_\mathrm{L}𝝈𝐯\|\psi =0,`$ (5b) whatever $`𝐮`$ and $`𝐯`$, as follows from the rotational invariance of the singlet state. Likewise, the expectation value of the spin correlation observable is $`E^\psi (𝐮,𝐯)=`$ $`\psi \|(𝝈_\mathrm{L}𝐮)(𝝈_\mathrm{R}𝐯)\|\psi `$ (6a) $`=`$ $`𝐮𝐯,`$ (6b) which depends only on the relative angle between $`𝐮`$ and $`𝐯`$ (see, for instance, Refs. FS1 , GHSZ1 , or AB1 ). ### II.3 Perfect correlations and hidden-variables When $`𝐮=𝐯`$, the expectation value of the spin correlation observable (6) is equal to $`1`$, meaning that if both Stern-Gerlach devices are oriented along the same direction, then with certainty the outcomes will be found to be opposite. Since the Stern-Gerlach devices are arbitrarily far from each other, a perfect correlation can be understood from a realistic point of view as invested in the particles at their inception. This, however, would mean that the singlet state is incomplete, and therefore, that it should be possible to give a more precise specification using additional “hidden-variables”. On the other hand, if a more complete description is impossible, then this perfect correlation seems rather mysterious, since a measurement performed on one of the subsystems seems to be capable of influencing the measurement on the other subsystem, whatever the distance between them. In order to facilitate a choice between incompleteness and non-locality of Quantum Mechanics, Bell’s idea was to specify mathematical requirements for a generic local hidden-variables theory, and then to compare its predictions with those from quantum mechanics and the results of experiments. In a local realistic hidden-variables model, a single particle pair is thus supposed to be entirely characterised by means of a set of hidden-variables, which are symbolically represented by a parameter $`\lambda `$, so that the measurement result on the left along $`𝐮`$ can be written as $`A(𝐮,\lambda )`$, and the result on the right along $`𝐯`$ as $`B(𝐯,\lambda )`$. Although the hidden-variables model is supposed to be fully deterministic, it must also be capable of reproducing the stochastic nature of the EPRB gedanken experiment expressed in Eqs. (5) and (6). For that purpose, the complete state specification $`\lambda _i`$ of any particle pair with label $`i`$ must be a random variable: its complete state $`\lambda _i`$ is supposed to be drawn randomly according to a probability distribution $`\rho `$ (see Refs. JSB2 and JSB4 ), meaning that the probability of having $`\lambda _i`$ equal to a particular $`\lambda `$ is $`\rho (\lambda )`$. Consider a set of $`N`$ particle pairs $`\{i=1,\mathrm{},N\}`$, the mean value of joint spin measurements for this set is : $$M^\rho (𝐮,𝐯)=\frac{1}{N}\underset{i=1}{\overset{N}{}}A(𝐮,\lambda _i)B(𝐯,\lambda _i).$$ (7) The probability distribution $`\rho `$ is supposed to assure the equality between this mean value and the expected value, Eq. (6b), given by Quantum Mechanics when $`N`$ goes to infinity. ## III The ‘CHSH’ function In order to establish Bell’s Theorem, a linear combination of correlation functions $`c(𝐚,𝐛)`$ with different argumentsJSB8 is considered, once when these correlation functions are expectation values $`E^\psi (𝐮,𝐯)`$ given by Quantum Mechanics; i.e., Eq. (6), and once when they are mean values $`M^\rho (𝐮,𝐯)`$ given by local hidden-variables theories, Eq. (7); then the results are to be compared. A well known choice of such a linear combination is the CHSH (Clauser, Horne, Shimony and Holt CHSH1 ) function, written with four pairs of arguments: $$S\|c(𝐚,𝐛)c(𝐚,𝐛^{})+c(𝐚^{},𝐛)+c(𝐚^{},𝐛^{})\|.$$ (8) Bell’s Theorem consists in showing that the quantum prediction for the CHSH function violates the maximum possible value given by any local realistic hidden-variables theory. Thus, no such theory will ever be capable of explaining or reproducing these quantum results. Herein, this claim is refuted by demanding the rules of Quantum Mechanics be consistently and meaningfully applied. To begin, the exact meaning of the simultaneous presence of different arguments in a CHSH function must be clarified. Basically, there are two possible interpretations, the strongly objective interpretation and the weakly objective interpretation BDE2 ; BDE1 : implies that all correlation functions are relevant to the same set of $`N`$ particle pairs, that is, all four pairs of directions are considered simultaneously relevant to each particle pair. As such they cannot be relevant to actual experiments but rather with what result would have been obtained if measured on the same set of $`N`$ particle pairs along different directions. implies that each correlation function is actually to be measured on distinct sets of $`N`$ particle pairs. Each set of $`N`$ particle pairs pertains to only one pair of arguments, that is, for each pair only one joint spin measurement is executed. The CHSH function was developed specifically for experimental convenienceCHSH1 . Many experiments have been done (the most famous being Aspect’s ADR1 ) obviously invoking the natural interpretation, namely the weakly objective one. Nevertheless, the strongly objective interpretation must also be considered, since it remains a possible interpretation a priori, and since the choice between strong and weak objectivity is not made at all explicit in many papers, including Bell’s. It must be stressed, moreover, that these interpretations are radically different, not only epistemologically, but also physically. Indeed, the strongly objective interpretation pertains to a single set of $`N`$ particle pairs characterised by the corresponding set of parameters $`\{\lambda _i;i=1,\mathrm{},N\}`$; whereas the weakly objective interpretation pertains to no less than 4 sets of $`N`$ particle pairs. The fact is that a finite set of $`N`$ particle pairs characterised by $`\{\lambda _i\}`$ can’t be identically reproduced, either theoretically (for each complete state $`\lambda _i`$ of any particle pair $`i`$ is a random variable, as defined in Section II.3), or empirically (for the experimenter has no control over the complete state of a particle pair in a singlet state). Of course, when $`N`$ goes to infinity, these four sets of $`N`$ particle pairs necessarily converge to the same ideal set described by the probability distribution $`\rho `$. However, as soon as real experiments are concerned, then $`N\mathrm{}`$ and these four sets are necessarily four different sets of particle pairs (see Ref. AB1 , page 348) respectively characterised by four different sets of hidden-variables parameters $`\{\lambda _{1,i}\}`$, $`\{\lambda _{2,i}\}`$, $`\{\lambda _{3,i}\}`$ and $`\{\lambda _{4,i}\}`$ (for an alternate approach, see Khrennikov1 ). The difference between each interpretation can therefore be embodied in the number of degrees of freedom of the whole system. Let $`f`$ be the degrees of freedom of a single particle pair. In the strongly objective interpretation the degrees of freedom of the whole CHSH system is then $`Nf`$, whereas in the weakly objective interpretation it is 4 times as large, that is, $`4Nf`$. Thus, before initiating Bell’s analysis, one has to choose explicitly one interpretation and stick to it. Unfortunately, this is not what has been done. It will be shown here that the discrepancy exhibited by Bell’s Theorem is due to a meaningless comparison between strongly objective and weakly objective results, which means comparing the numerical value of the CHSH function for two systems, one having $`Nf`$ degrees of freedom, the other $`4Nf`$. ## IV The Strongly objective interpretation: counterfactual properties of $`N`$ particle pairs ### IV.1 A Local realistic inequality within the strongly objective interpretation The local realistic formulation of the CHSH function within strong objectivity is written $`S_{\text{strong}}^\rho =\left\|M^\rho (𝐚,𝐛)M^\rho (𝐚,𝐛^{})+M^\rho (𝐚^{},𝐛)+M^\rho (𝐚^{},𝐛^{})\right\|,`$ (9) or explicitly (using Eq. 7) $$\begin{array}{cc}\hfill S_{\text{strong}}^\rho =\|\frac{1}{N}\underset{i=1}{\overset{N}{}}& A(𝐚,\lambda _i)B(𝐛,\lambda _i)A(𝐚,\lambda _i)B𝐛^{},\lambda _𝐢)\hfill \\ \hfill +& A(𝐚^{},\lambda _i)B(𝐛,\lambda _i)+A(𝐚^{},\lambda _i)B(𝐛^{},\lambda _i)\|,\hfill \end{array}$$ (10) which after factorisation becomes $$\begin{array}{cc}\hfill S_{\text{strong}}^\rho =\|\frac{1}{N}\underset{i=1}{\overset{N}{}}& A(𝐚,\lambda _i)\left[B(𝐛,\lambda _i)B(𝐛^{},\lambda _𝐢)\right]\hfill \\ \hfill +& A(𝐚^{},\lambda _i)[B(𝐛,\lambda _i)+B(𝐛^{},\lambda _i)]\|,\hfill \end{array}$$ (11) where each term can have two values in the summation FS1 ; AB1 $$\begin{array}{cc}\hfill A(𝐚,\lambda _i)& \left[B(𝐛,\lambda _i)B(𝐛^{},\lambda _𝐢)\right]\hfill \\ \hfill +& A(𝐚^{},\lambda _i)\left[B(𝐛,\lambda _i)+B(𝐛^{},\lambda _i)\right]=\pm 2,\hfill \end{array}$$ (12) so that the most restrictive local realistic inequality within the strongly objective interpretation is : $$S_{\text{strong}}^\rho 2.$$ (13) This is the well known generalised formulation of Bell’s inequality due to CHSH CHSH1 . It must be stressed once more, however, that this inequality has been established only within the strongly objective interpretation, which means that each expectation value is relevant to the same set of $`N`$ particle pairs. Hence, this result cannot be compared directly with results from real experimental tests, where in fact mean values from four distinct sets of $`N`$ particle pairs are measured. The question whether the same inequality can be applied to real experiments will be discussed in Section V.2 (weak objectivity). ### IV.2 The Quantum mechanical prediction within the strongly objective interpretation The quantum prediction for the CHSH function within the strongly objective interpretation is written $$S_{\text{strong}}^\psi =\|E^\psi (𝐚,𝐛)E^\psi (𝐚,𝐛^{})+E^\psi (𝐚^{},𝐛)+E^\psi (𝐚^{},𝐛^{})\|.$$ (14) This equation is usually directly evaluated by replacing each expectation value by the scalar product result of Eq. (6b). This, unfortunately, is all too hasty. Indeed, in order to understand better the quantum mechanical meaning of Eq. (14), it is advantageous to take a step backward using Eq. (6a) $$\begin{array}{cc}\hfill S_{\text{strong}}^\psi =\|\psi \|(𝝈_\mathrm{L}𝐚)& (𝝈_\mathrm{R}𝐛)\|\psi \psi \|(𝝈_\mathrm{L}𝐚)(𝝈_\mathrm{R}𝐛^{})\|\psi \hfill \\ \hfill +& \psi \|(𝝈_\mathrm{L}𝐚^{})(𝝈_\mathrm{R}𝐛)\|\psi +\psi \|(𝝈_\mathrm{L}𝐚^{})(𝝈_\mathrm{R}𝐛^{})\|\psi \|,\hfill \end{array}$$ (15) or again $$S_{\text{strong}}^\psi =\left\|\psi \|(𝝈_\mathrm{L}𝐚)(𝝈_\mathrm{R}𝐛)(𝝈_\mathrm{L}𝐚)(𝝈_\mathrm{R}𝐛^{})+(𝝈_\mathrm{L}𝐚^{})(𝝈_\mathrm{R}𝐛)+(𝝈_\mathrm{L}𝐚^{})(𝝈_\mathrm{R}𝐛^{})\|\psi \right\|.$$ (16) Note, however, that the four spin correlation observables in this equation are non commuting observables (this can be shown by calculating the commutator of $`(𝝈_\mathrm{L}𝐮)(𝝈_\mathrm{R}𝐯)`$ and $`(𝝈_\mathrm{L}𝐮)(𝝈_\mathrm{R}𝐯^{})`$ with $`𝐯𝐯^{}`$), so that the meaning of their combination must be questioned. The problem is that it is actually impossible to find an eigenvector for this combination of observables. Indeed, this linear combination of observables is, after factorisation, an operator of the form $`\widehat{A}\widehat{B}+\widehat{C}\widehat{D}`$ with $`[\widehat{A},\widehat{C}]0`$ and $`[\widehat{B},\widehat{D}]0`$. In the Hilbert space $``$, an hypothetical eigenvector $`\|\varphi \|\chi `$ (with $`\alpha `$ being its eigenvalue) of this operator should satisfy $$\left[\widehat{A}\widehat{B}+\widehat{C}\widehat{D}\right]\|\varphi \|\chi =\alpha (\|\varphi \|\chi ),$$ (17) that is, $$\widehat{A}\|\varphi \widehat{B}\|\chi +\widehat{C}\|\varphi \widehat{D}\|\chi =\alpha (\|\varphi \|\chi ).$$ (18) This equation can have solutions only if its left hand side can be factored; that is, either $`\widehat{A}\|\varphi `$ and $`\widehat{C}\|\varphi `$, or $`\widehat{B}\|\chi `$ and $`\widehat{D}\|\chi `$ must be colinear vectors. This, however, can never happen because both $`[\widehat{A},\widehat{C}]0`$ and $`[\widehat{B},\widehat{D}]0`$. Hence, Eq. (17) has no solution, and the linear combination of observables in Eq. (16) has no eigenvector: it is not an observable, and thus it can’t be given physical meaning. Therefore, $`S_{\text{strong}}^\psi `$ is meaningless and is not a proper equation to use in order to make physical predictions. Of course, this does not imply that Quantum Mechanics cannot provide any meaning at all for the CHSH function; it implies only that this meaning cannot be strongly objective. Indeed, according to Von Neumann JVN1 , any linear combination of expectation values of different observables $`\widehat{R}`$, $`\widehat{S},\mathrm{}`$ is meaningful in Quantum Mechanics: $$\widehat{R}+\widehat{S}+\mathrm{}_\varphi =\widehat{R}_\varphi +\widehat{S}_\varphi +\mathrm{}$$ (19) even if $`\widehat{R}`$, $`\widehat{S},\mathrm{}`$ are non commuting observables. The explanation is that Quantum Mechanics is only a weakly objective theory BDE2 ; BDE3 , and that expectation values given by Quantum Mechanics are also weakly objective statements, that is to say, statements relevant to observations. Hence, when $`\widehat{R}`$, $`\widehat{S},\mathrm{}`$ are non commuting observables, the expectation values cannot be simultaneously relevant to the same set of $`N`$ systems: each expectation value is necessarily relevant to a distinct set of $`N`$ systems (all systems being represented by the quantum state $`\|\varphi `$). Likewise, the only possible meaning of Eq. (15) is therefore weakly objective, not strongly objective as desired. Since these expectation values are known with certainty, it is tempting to consider them as counterfactual entities. However, conterfactuality requires at least measurement compatibility, that is, commuting observables. The certainty of a contextual prediction is not sufficient to make it a counterfactual prediction; in other words, weakly objective results known with certainty are not strongly objective results. Incidentally, this is also true in the case of perfect correlations, so that as a general rule, one may not manipulate weakly objective results as if they were strongly objective. The local realistic inequality $`S_{\text{strong}}^\rho `$ cannot be compared with any strongly objective prediction given by Quantum Mechanics, so that Bell’s Theorem cannot be verified with a strongly objective interpretation given to the CHSH function, simply because Quantum Mechanics is not a strongly objective theory. This restriction is the first part of a refutation of Bell’s theorem, though maybe not conclusive, since the strength of Bell’s Theorem is mainly its amenability to experimental test. Still, this was necessary, for now that a strongly objective interpretation is precluded, there is no choice but to rely on the weakly objective interpretation in order to compare hidden-variables theories and Quantum Mechanics. In the next section, a simple method will be provided in order to obtain a unique and meaningful quantum prediction for the CHSH function within weak objectivity. ## V The Weakly objective interpretation : contextual measurements on $`4`$ distinct sets of $`N`$ particle pairs ### V.1 A Quantum mechanical prediction within the weakly objective interpretation It was shown in Section III that strong objectivity and weak objectivity pertain to different physical systems. This difference should therefore appear in the relevant equations. Indeed, the correlation expressed in Eq. (6b) is relevant to spin measurements performed on particles that once constituted a single parent particle. Yet, two particles issued from two distinct parents never have interacted with each other, so that spin measurements performed on such particle pairs can not be correlated. Hence, if left and right spin measurements are performed on two distinct sets of $`N`$ particle pairs, instead of the same set, there should be no correlation, and this property should appear in a generalised spin correlation function (i.e. generalised to the case of spin measurements performed on different sets of particle pairs). This can be easily done within a quantum theoretical framework by means of a distinct EPRB space for each set of $`N`$ particle pairs. Let $`_j`$ be the EPRB Hilbert space associated with the $`j`$th set of particle pairs. In this Hilbert space, the EPRB gedanken experiment is represented by the singlet state $`\|\psi _j`$ (see Section II), $$\|\psi _j=\frac{1}{\sqrt{2}}[\|+_j\|+_j].$$ (20) The whole CHSH experiment with the four sets of particle pairs can be expressed then in terms of a new direct product space $`_{1234}_1_2_3_4`$ in which the state vector is $$\|\psi _{1234}=\|\psi _1\|\psi _2\|\psi _3\|\psi _4.$$ (21) The counterparts of observables in $`_{1234}`$ are obtained as in Section II.1. For instance, the observables pertaining to the right Stern-Gerlach device for the 1st, 2nd, 3rd and 4th set of particle pairs are respectively $`𝝈_{1,\mathrm{R}}𝐮(𝝈_\mathrm{R}𝐮)1\mathrm{l}_21\mathrm{l}_31\mathrm{l}_4,`$ (22a) $`𝝈_{2,\mathrm{R}}𝐮1\mathrm{l}_1(𝝈_\mathrm{R}𝐮)1\mathrm{l}_31\mathrm{l}_4,`$ (22b) $`𝝈_{3,\mathrm{R}}𝐮1\mathrm{l}_11\mathrm{l}_2(𝝈_\mathrm{R}𝐮)1\mathrm{l}_4,`$ (22c) $`𝝈_{4,\mathrm{R}}𝐮1\mathrm{l}_11\mathrm{l}_21\mathrm{l}_3(𝝈_\mathrm{R}𝐮),`$ (22d) where $`1\mathrm{l}_j`$ is the identity operator of the EPRB space $`_j`$. Hence, the expectation value of the product of two spin observables, the first belonging to the $`k`$th set and the second to the $`l`$th set, is $$E_{kl}^\psi (𝐮,𝐯)\psi _{1234}\|(𝝈_{k,L}𝐮)(𝝈_{l,R}𝐯)\|\psi _{1234},$$ (23) and this is the generalised expectation value of spin correlation observables that was sought. The expectation value for measurements performed on the same set ($`k=l`$) of particle pairs is already known, Eq. (6), and $`E_{kk}^\psi (𝐮,𝐯)`$ should provide the same result. Indeed, using Eqs. (21) and (22) leads to $`E_{kk}^\psi (𝐮,𝐯)`$ $`=\psi _k\|(𝝈_L𝐮)(𝝈_R𝐯)\|\psi _k`$ (24) $`=𝐮𝐯,`$ but when $`kl`$, the result is quite different: $`E_{\begin{array}{c}kl\hfill \\ kl\hfill \end{array}}^\psi (𝐮,𝐯)`$ $`=\psi _k\|(𝝈_L𝐮)\|\psi _k\psi _l\|(𝝈_R𝐯)\|\psi _l`$ $`=\psi _k\|𝝈𝐮1\mathrm{l}_\mathrm{R}\|\psi _k\psi _l\|1\mathrm{l}_\mathrm{L}𝝈𝐯\|\psi _l`$ (25) $`=0,`$ in accord with Eq. (5). There are indeed no correlations between two sets of particle pairs, as stipulated in the beginning of this section. Now, contrary to what was done in Section IV.2, it is possible to proceed here in full accord with the quantum mechanical postulates, because the spin correlation observables, Eqs. (22), are mutually commuting, so that a linear combination of these commuting observables is an observable as well. The CHSH experiment can therefore be described by a new observable $$\begin{array}{cc}\hfill \widehat{S}_{\text{weak}}(𝝈_{1,L}𝐚)& (𝝈_{1,R}𝐛)(𝝈_{2,L}𝐚)(𝝈_{2,R}𝐛^{})\hfill \\ \hfill +& (𝝈_{3,L}𝐚^{})(𝝈_{3,R}𝐛)+(𝝈_{4,L}𝐚^{})(𝝈_{4,R}𝐛^{}),\hfill \end{array}$$ (26) and the quantum prediction for the CHSH function within a weakly objective interpretation is therefore obtained by calculating the expectation value of the observable $`\widehat{S}_{\text{weak}}`$ when the system is in the quantum state $`\|\psi _{1234}`$ : $$S_{\text{weak}}^\psi =\left\|\psi _{1234}\|\widehat{S}_{\text{weak}}\|\psi _{1234}\right\|,$$ (27) which using Eqs. (22) and (23) is $$\begin{array}{cc}\hfill S_{\text{weak}}^\psi =\|\psi _1\|(𝝈_L𝐚)& (𝝈_R𝐛)\|\psi _1\psi _2\|(𝝈_L𝐚^{})(𝝈_R𝐛)\|\psi _2\hfill \\ \hfill +& \psi _3\|(𝝈_L𝐚)(𝝈_R𝐛^{})\|\psi _3+\psi _4\|(𝝈_L𝐚^{})(𝝈_R𝐛^{})\|\psi _4\|,\hfill \end{array}$$ (28) that is, using Eq. (24), $$S_{\text{weak}}^\psi =\left\|E_{11}^\psi (𝐚,𝐛)E_{22}^\psi (𝐚,𝐛^{})+E_{33}^\psi (𝐚^{},𝐛)+E_{44}^\psi (𝐚^{},𝐛^{})\right\|.$$ (29) This equation is not ambiguous (as was Eq. 15): it is a linear combination of expectation values, each relevant to a distinct set of $`N`$ particle pairs. This equation is therefore weakly objective, as requested. Finally, using Eq. (24), yields $$S_{\text{weak}}^\psi =\left\|𝐚𝐛𝐚𝐛^{}+𝐚^{}𝐛+𝐚^{}𝐛^{}\right\|,$$ (30) with a well known maximum equal to $$\mathrm{max}(S_{\text{weak}}^\psi )=2\sqrt{2}.$$ (31) This numerical result is indeed the one given in the literature, the only difference here being the fact that the meaning of this result is unambiguously weakly objective. Quantum Mechanics, which is a weakly objective theory BDE2 , provides a clear answer to the CHSH function understood as a weakly objective question. ### V.2 A Local realistic inequality within the weakly objective interpretation The last step consists in comparing the quantum prediction $`S_{\text{weak}}^\psi `$ with its local realistic counterpart $`S_{\text{weak}}^\rho `$. As was stressed in Section III, the $`j`$th set of particle pairs must be characterised by a distinct set of hidden-variables parameters $`\{\lambda _{j,i};i=1,\mathrm{},N\}`$. Hence, to the generalised expectation value of the spin correlation observable Eq. (23) corresponds the generalised mean value of joint spin measurements: $$M_{kl}^\rho (𝐮,𝐯)\frac{1}{N}\underset{i=1}{\overset{N}{}}A(𝐮,\lambda _{k,i})B(𝐯,\lambda _{l,i}),$$ (32) which is a priori capable of reproducing not only the $`k=l`$ prediction, Eq. (24), but also the $`kl`$ prediction, Eq. (V.1). The local realistic CHSH function with a weakly objective interpretation is therefore $$S_{\text{weak}}^\rho =\left\|M_{11}^\rho (𝐚,𝐛)M_{22}^\rho (𝐚,𝐛^{})+M_{33}^\rho (𝐚^{},𝐛)+M_{44}^\rho (𝐚^{},𝐛^{})\right\|,$$ (33) and that is explicitly $$\begin{array}{cc}\hfill S_{\text{weak}}^\rho =\|\frac{1}{N}\underset{i=1}{\overset{N}{}}[A(𝐚,\lambda _{1,i})& B(𝐛,\lambda _{1,i})A(𝐚,\lambda _{2,i})B(𝐛^{},\lambda _{2,i})\hfill \\ \hfill +& A(𝐚^{},\lambda _{3,i})B(𝐛,\lambda _{3,i})+A(𝐚^{},\lambda _{4,i})B(𝐛^{},\lambda _{4,i})\left]\right\|.\hfill \end{array}$$ (34) This expression is to be compared with the one pertaining to the strongly objective interpretation, Eq. (10), which contained terms that could be factored. Here, since each term is different from the others, no factorisation is possible; i.e., there is no way to derive a Bell inequality. This is not the first time this fact has been noticed (see A. Bohm pp. 351, 352 AB1 ), unfortunately, no conclusion was drawn then. Yet, this fact cannot be ignored, for it has been shown in Section IV.2 that Bell’s Theorem cannot be demonstrated within a strongly objective interpretation. Here, the only local realistic inequality that can be derived is obtained by considering (as was done with Eq. 12) the possible numerical values of each term of the summation in Eq. (34), that is, $$\begin{array}{cc}\hfill A(𝐚,\lambda _{1,i})B(𝐛,\lambda _{1,i})A(𝐚,\lambda _{2,i})& B(𝐛^{},\lambda _{2,i})+A(𝐚^{},\lambda _{3,i})B(𝐛,\lambda _{3,i})\hfill \\ \hfill +& A(𝐚^{},\lambda _{4,i})B(𝐛^{},\lambda _{4,i})=+4,+2,0,2,4,\hfill \end{array}$$ (35) for which the extrema are +4 and -4, so that the narrowest local realistic inequality that can be derived from Eq. (34) is nothing but $$S_{\text{weak}}^\rho 4.$$ (36) This most restrictive local realistic inequality (which can also be found in AccardiAccardi ) is not incompatible with the quantum mechanical prediction, as the maximum of $`S_{\text{weak}}^\psi `$ is $`2\sqrt{2}`$. This shows that experiments intended to test Bell’s Theorem were unfortunately not testing the strongly objective inequality (a Bell inequality, Eq. 13), but this weakly objective one, Eq. (36), since all experimental tests necessarily are executed in a weakly objective way, due to the irreducible incompatibility between spin measurements. As was stressed by Sica Sica and Accardi Accardi , a local realistic inequality is nothing but an arithmetic identity, and inequality (36) is definitely too lax to be violated by experimental tests. ## VI Conclusion It was shown that Bell’s Theorem cannot be derived, either within a strongly objective interpretation of the CHSH function, because Quantum Mechanics gives no strongly objective results for the CHSH function (see Section IV.2), or within a weakly objective interpretation, because the only derivable local realistic inequality is never violated, either by Quantum Mechanics or by experiments (see Section V.2). It was demonstrated that the discrepancy in Bell’s Theorem is due only to a meaningless comparison between $`S_{\text{strong}}^\rho 2`$ and $`S_{\text{weak}}^\psi =2\sqrt{2}`$, where the former is relevant to a system with $`Nf`$ degrees of freedom, whereas the latter to one with $`4Nf`$ (see Section III). The only meaningful comparison is between the weakly objective local realistic inequality $`S_{\text{weak}}^\rho 4`$ and the weakly objective quantum prediction $`S_{\text{weak}}^\psi =2\sqrt{2}`$, but these results are not incompatible. Bell’s Theorem, therefore, is refuted.
warning/0006/hep-ex0006040.html	ar5iv	text	# Measurement of the K{_𝐿}→𝑒⁺𝑒⁻𝑒⁺𝑒⁻ Decay Rate ## I Introduction The K$`{}_{L}{}^{}e^+e^{}e^+e^{}`$ decay is expected to proceed mainly via the intermediate state K$`{}_{L}{}^{}\gamma ^{}\gamma ^{}`$ and thus depends on the structure of the K$`{}_{L}{}^{}\gamma ^{}\gamma ^{}`$ vertex. Phenomenological models include vector meson dominance of the photon propagator , QCD inspired models , intermediate pseudoscalar and vector mesons and models based on chiral perturbation theory . The probability for both virtual photons to convert into $`e^+e^{}`$ pairs is calculated to be in the range $`(5.896.50)\times 10^5`$ . The chiral model prediction of corresponds to $`\mathrm{\Gamma }`$(K$`{}_{L}{}^{}e^+e^{}e^+e^{})/\mathrm{\Gamma }(`$K$`{}_{L}{}^{}`$all$`)=3.85\times 10^8`$, including the effect of a form factor, which increases the width by 4%. The interference term due to the identity of particles has been calculated to change the branching ratio by 0.5%. The decay was first observed by the CERN NA31 experiment based on 2 observed events and has been confirmed by later measurements . Here we report the preliminary result obtained from the 1999 data taking of the CERN experiment NA48. ## II Experimental Setup and Data Taking This measurement was carried out as part of the CERN experiment NA48 at the CERN SPS, which has previously reported measurements of the related decays K$`{}_{L}{}^{}e^+e^{}\gamma `$ and K$`{}_{L}{}^{}\mu ^+\mu ^{}\gamma `$ . A detailed and comprehensive description of the detector is in preparation . The NA48 experiment is designed specifically to measure the direct CP violation parameter $`\mathrm{}(ϵ^{}/ϵ)`$ using simultaneous beams of K<sub>L</sub> and K<sub>S</sub>. To produce the K<sub>L</sub> beam, 450 GeV/c protons are extracted from the accelerator during 2.4 s every 14.4 s and $`1.1\times 10^{12}`$ of these are delivered to a beryllium target. Using dipole magnets to sweep away charged particles and collimators to define a narrow beam, a neutral beam of $`2\times 10^7`$ K<sub>L</sub> per burst and divergence $`\pm 0.15`$ mrad enters the decay region. The fiducial volume begins 126 m downstream from the target and is contained in an evacuated cylindrical steel vessel 89 m long and 2.4 m in maximum diameter. The vessel is terminated at the downstream end by a Kevlon-fiber composite window and followed immediately by the main NA48 detector. The sub-detectors which are used in the K$`{}_{L}{}^{}e^+e^{}e^+e^{}`$ analysis are described below in order of succession (see Fig. 1). A magnetic spectrometer consisting of a dipole magnet is preceded and followed by two sets of drift chambers. The drift chambers are each comprised of eight planes of sense wires, two horizontal, two vertical and two along each of the 45 directions. Only the vertical and horizontal planes are instrumented in the third chamber. The volume between the chambers is filled with helium at atmospheric pressure. The momentum resolution is $`\mathrm{\Delta }p/p=0.65\%`$ at 45 GeV/c. Two segmented plastic scintillator hodoscope planes are placed after the helium tank and provide signals for the trigger. A liquid krypton filled calorimeter (LKr) is used for measuring the energy, position and timing of electromagnetic showers. Spacial and timing resolutions of better than 1.3 mm and 300 ps, respectively, have been achieved for energies above 20 GeV. The energy resolution is $`\frac{\sigma (E)}{E}=\frac{0.035}{\sqrt{E}}\frac{0.110}{E}0.006`$, with $`E`$ measured in GeV. A hadron calorimeter composed of 48 steel plates, each 24 mm thick, interleaved with scintillator is used in trigger formation and particle detection studies. ## III Trigger and Data Processing Candidate events were selected by a two-stage trigger. At the first level, a trigger requiring adjacent hits in the hodoscope is put in coincidence with a total energy condition ($`35`$ GeV), defined by adding the energy deposited in the hadronic calorimeter with that seen by the trigger in the LKr calorimeter. The second level trigger uses information from the drift chambers to reconstruct tracks and invariant masses. For the 4-track part of the trigger, the number of clustered hits in each of the first, second, and fourth drift chamber had to be between 3 and 7. All possible 2-track vertices were calculated online. At least two vertices within 6 m of each other in the axial direction had to be found. For the determination of the 4-track trigger efficiency, downscaled events that passed the first level were used. Alternatively, events triggered with the neutral trigger - based on the data of the LKr calorimeter - were selected. The neutral trigger applied the following cuts to the events online: $``$ 5 peaks in each projection, total energy $`>`$ 50 GeV, first moment of cluster energies (‘center of gravity’) $`<15`$ cm, and lifetime uncorrected for deflection in the magnet $`<4.5\tau _{K_S}`$. During the experimental runs, roughly 100 Terabytes of raw data with typically 20 kbytes per event were recorded. Reconstructed output was stored in a compressed data format, 45 times smaller. In addition, several streams of data were formed for events accepted by about 40 filter modules for the analysis of neutral and charged two-pion decays, rare decays, and events for the detector calibration. ## IV Data Analysis The data sample which yields K$`{}_{L}{}^{}e^+e^{}e^+e^{}`$ also has been used to select K$`{}_{L}{}^{}\pi ^+\pi ^{}\pi _{Dalitz}^0`$ and K$`{}_{L}{}^{}\pi ^0\pi _{Dalitz}^0\pi _{Dalitz}^0`$ normalization events, with $`\pi _{Dalitz}^0e^+e^{}\gamma `$. The vertex was reconstructed from the 4 tracks by requiring that the sum of the squared transverse distances from the transverse vertex position, weighted by the inverse track momentum, be minimal. Events were preselected by requiring two positive and two negative tracks with distance of closest approach to the vertex $`<15`$ cm. All clusters in the LKr were required to be in a fiducial area given by an octogon about 5 cm smaller than the outside perimeter of the calorimeter and an inner radius of 15 cm. The distance to dead calorimeter cells had to exceed 2 cm to ensure negligible energy loss. The separation between cluster centers was required to be $`>`$ 3 cm. The energy of each cluster was required to exceed 2 GeV, well above the detector noise of 0.11 GeV per cluster. Electron candidates were identified by requiring that cluster centers in the LKr be within 1.5 cm of the calculated shower maxima based upon the extrapolation of each track. The efficiency of this procedure was measured to be $`(99.7\pm 0.1)\%`$ . To reject pion showers, the ratio of cluster energy to track momentum $`E/p`$ was required to lie between 0.9 and 1.2. The efficiency of this selection was determined to be $`95\%`$ . Those track-associated clusters with $`0<E/p<0.8`$ were classified as pions. From a study of K$`{}_{S}{}^{}\pi ^+\pi ^{}`$ decays, the probability of pions to be wrongly classified as electrons is estimated to be 0.9%; the pion classification is passed by 97.5% of all pions. The fiducial volume was defined by the axial postion 7.50 m $`<z_{vertex}<90`$ m downstream of the K<sub>S</sub> target. Within this volume, 4-track vertices were determined with a typical axial resolution of 0.5 m, as estimated by the MC. The K<sub>L</sub> energy had to be in the range 50 GeV$``$200 GeV. ### A Background Rejection and Selection of K$`{}_{L}{}^{}e^+e^{}e^+e^{}`$ Candidates Candidate events for the decay K$`{}_{L}{}^{}e^+e^{}e^+e^{}`$ with all tracks identified as electrons were selected. In principle, the following four classes of background sources are relevant: * Events with two decays K$`{}_{L}{}^{}\pi e\nu `$ occuring at the same time and for which the pions were misidentified as electrons. Being due to two coincident kaon decays the invariant mass of the system can be around and above the nominal K<sub>L</sub> mass. These events are largely rejected by requiring a good vertex quality: the closest distance of approach of each track w.r.t. the reconstructed vertex had to be $`<`$ 5 cm. Events with separate vertices do not pass the level 2 trigger (see above). Finally it was required that the measured times for clusters associated to the electrons had to be consistent within 3 ns with each other. A study of sidebands in this time distribution shows that the background from this source is negligible. * Events K$`{}_{L}{}^{}\pi ^0\pi ^0\pi ^0`$, where the $`\pi ^0`$’s undergo single or double Dalitz decays or photons convert in the material of the detector, so that 2 positive and 2 negative electrons are detected. Due to the missing photons, the invariant mass of the $`e^+e^{}e^+e^{}`$ system is below the nominal K<sub>L</sub> mass. The loose requirement that the square of the reconstructed transverse momentum $`p_T^2`$ of the reconstructed kaon with respect to the line joining the decay vertex and the K<sub>L</sub> target must be less than 0.0005 (GeV/c)<sup>2</sup>, restricts this background to less than 2.2% . The position of the cut is indicated in Fig. 2. The Monte Carlo simulation indicates that 1.5% of the signal events are lost by the $`p_T^2`$ cut. * Events K$`{}_{L}{}^{}\gamma \gamma `$ and K$`{}_{L}{}^{}e^+e^{}\gamma `$, with conversion of the photons in the material upstream of the spectrometer also yield invariant masses around the nominal K<sub>L</sub> mass. Each pair of oppositely charged tracks was therefore required to be separated by $`2`$ cm in the first drift chamber as indicated in Fig. 3(a). Note that the conversion probability in the material of the NA48 detector is of similar magnitude as that for internal photon conversion to a $`e^+e^{}`$ pair. As the angular opening of oppositely charged tracks peaks mostly at small angles, 60% of the signal events are lost by this cut; according to the MC, there is no remaining background with converted $`\gamma `$’s. * Events K$`{}_{L}{}^{}\pi ^+\pi ^{}e^+e^{}`$ , with the pions misidentified as electrons. Due to the low misidentification probability of 0.9% this background is found to be negligible. The invariant $`e^+e^{}e^+e^{}`$ mass plot resulting from this selection is shown in Fig. 4. Note the slightly asymmetric shape of the K<sub>L</sub> mass peak, which is due to photons radiated off the electrons in the final state. Finally, a mass window of 0.475 GeV/$`c^2<m(e^+e^{}e^+e^{})<`$ 0.520 GeV/$`c^2`$ was set to define the final sample. In total, 132 candidate events were selected in this way. From the number of events observed below the K<sub>L</sub> mass peak and their distribution we estimate that the background contribution to the signal region is negligible. ### B Normalization The four-track decay K$`{}_{L}{}^{}\pi ^+\pi ^{}\pi _{Dalitz}^0`$, with $`\pi _{Dalitz}^0e^+e^{}\gamma `$, was used for normalization. Only events that passed the pretrigger discussed above, downscaled by 60, were selected. This sample was used to determine the efficiency of the 4-track trigger to be $`95.1\pm 0.2\%`$. The 4-track trigger is required for signal and normalization modes. Since both the signal and the normalization modes consist of 4-track events, uncertainties due to tracking tend to cancel in the ratio of acceptances. Selection criteria similar to those used in the signal mode were applied. In addition, at least one extra cluster in the calorimeter not associated with a charged track was required. The invariant mass of the $`e^{}e^{}\gamma `$ system had to be in the range of 0.115 – 0.155 GeV/$`c^2`$. Monte Carlo studies showed that 66 % of the reconstructed K$`{}_{L}{}^{}\pi ^+\pi ^{}\pi _D^0`$ candidate decays are from K$`{}_{L}{}^{}\pi ^+\pi ^{}\pi ^0`$ with one of the external photons converting in the material of the detector (see Fig. 3(b)). All other backgrounds have been estimated to be negligible. Using these cuts, 17123 K$`{}_{L}{}^{}\pi ^0\pi ^0\pi _D^0`$ decays were selected. In a second analysis, K$`{}_{L}{}^{}\pi ^0\pi _{Dalitz}^0\pi _{Dalitz}^0`$ events were selected, yielding 5167 events for normalization. While having a complicated topology of eight clusters, these events have the advantage that all decay products interact electromagnetically in the detector and that the radiative corrections should be similar to those in the signal mode. ### C Acceptance Determination and Kaon Flux For the simulation of the K$`{}_{L}{}^{}e^+e^{}e^+e^{}`$ acceptance, the matrix element was taken from Ref. neglecting the interference of the two virtual photons. The distribution of the angle spanned by the decay planes of the two $`e^+e^{}`$ pairs corresponds to a K<sub>L</sub> which is assumed to be entirely CP$`=1`$. The PHOTOS package has been used to simulate final state radiation both for the signal and normalization channels. The acceptance for K$`{}_{L}{}^{}e^+e^{}e^+e^{}`$ is calculated to be 7.8% for events generated in the range 45 GeV $`<E_{K_L}<215`$ GeV and 5 m $`<z_{vertex}<91`$ m. The normalization to the K<sub>L</sub> signal has been measured from the number of K$`{}_{L}{}^{}\pi ^+\pi ^{}\pi _{Dalitz}^0`$ decays in the same sample. Using the acceptance of 1.34%, calculated by Monte Carlo simulation, a total number of $`5.1\times 10^{10}`$ K<sub>L</sub> decays is obtained. With the cuts described above, the inclusion of radiative corrections decreased the acceptance of signal and normalization channel by 8.8% and 3.4%, respectively. Consistent results for the total number of K<sub>L</sub> were obtained when the alternative normalization channel, K$`{}_{L}{}^{}\pi ^0\pi _{Dalitz}^0\pi _{Dalitz}^0`$, was used instead. ## V Results and Discussion A study of the stability of the branching ratio determination was made as a function of the cuts applied. A systematic error of 3.5% was estimated, mainly being due to the variation in the minimal distance of clusters from the beam pipe and the cut on the vertex quality. A second contribution comes from the 4-track trigger inefficiency. A Monte Carlo simulation of the level 2 algorithm yields a 95.0% efficiency for the normalization mode, in good agreement with the measured value. For the signal, the simulated efficiency is higher (99.8%). We chose not to apply a correction to the branching ratio; instead we introduced a systematic error of $`\pm `$5% for this preliminary result. Finally, the effect of overflows in the drift chambers has been considered. If events with an overflow condition in a window of 312 ns around the event time are removed, 20% of the events are lost and the branching ratio stays constant within 1%. Adding these sources of systematic error in quadrature, we obtain a total systematic error of $`\pm `$ 6.2%. From the numbers given above, a branching ratio of $`{\displaystyle \frac{\mathrm{\Gamma }(\mathrm{K}_Le^+e^{}e^+e^{})}{\mathrm{\Gamma }(\mathrm{K}_L\mathrm{all})}}=(3.67\pm 0.32_{stat}\pm 0.23_{syst}\pm 0.08_{norm})\times 10^8`$is obtained, where the statistical and systematical uncertainties as well as the uncertainty in the branching ratio of the normalization channel are given separately. This result is consistent with the theoretical expectation of and the previous average value of ($`4.1\pm 0.8)\times 10^8`$ , with 5 times the statistics of the single best previous experiment. . ## Acknowledgements It is a pleasure to thank the technical staff of the participating laboratories, universities and affiliated computing centers for their efforts in the construction of the NA48 apparatus, in operation of the experiment, and in the processing of the data.
warning/0006/math0006212.html	ar5iv	text	# Homogeneous symplectic manifolds with Ricci-type curvature ## 1 Let $`(M,\omega )`$ be a symplectic manifold and $``$ be a symplectic connection (a torsion-free connection on $`TM`$ with $`\omega =0`$). The curvature endomorphism $`R`$ of $``$ is defined by $$R(X,Y)Z=\left(_X_Y_Y_X_{[X,Y]}\right)Z$$ for vector fields $`X,Y,Z`$ on $`M`$. The symplectic curvature tensor $$R(X,Y;Z,T)=\omega (R(X,Y)Z,T)$$ is antisymmetric in its first two arguments, symmetric in its last two and satisfies the first Bianchi identity $$_{X,Y,Z}R(X,Y;Z,T)=0$$ where $``$ denotes the sum over the cyclic permutations of the listed set of elements. The Ricci tensor $`r`$ is the symmetric $`2`$-tensor $$r(X,Y)=Trace[ZR(X,Z)Y].$$ $`R`$ also obeys the second Bianchi identity $$_{X,Y,Z}\left(_XR\right)(Y,Z)=0.$$ The Ricci part $`E`$ of the curvature tensor is given by $`E(X,Y;Z,T)`$ $`=`$ $`{\displaystyle \frac{1}{2(n+1)}}[2\omega (X,Y)r(Z,T)+\omega (X,Z)r(Y,T)+\omega (X,T)r(Y,Z)`$ (1.1) $`\omega (Y,Z)r(X,T)\omega (Y,T)r(X,Z)].`$ The curvature is of Ricci type when $`R=E`$. ###### Lemma 1 Let $`(M,\omega )`$ be a symplectic manifold of dimension $`2n4`$. If the curvature of a symplectic connection $``$ on $`M`$ is of Ricci type then there is a $`1`$-form $`u`$ such that $$\left(_Xr\right)(Y,Z)=\frac{1}{2n+1}\left(\omega (X,Y)u(Z)+\omega (X,Z)u(Y)\right).$$ (1.2) Conversely, if there is such a $`1`$-form $`u`$ then the $`W`$ part of the curvature satisfies $$_{X,Y,Z}\left(_XW\right)(Y,Z;T,U)=0.$$ (1.3) Proof When the curvature is of Ricci type, the second Bianchi identity for $`R`$ becomes an identity for $`E`$. Since $`\omega `$ is parallel, covariantly differentiating equation (1.1) and summing cyclically, we get $`0`$ $`=`$ $`{\displaystyle _{X,Y,Z}}2\omega (Y,Z)(_Xr)(T,U)+\omega (Y,T)(_Xr)(Z,U)+\omega (Y,U)(_Xr)(Z,T)`$ (1.4) $`\omega (Z,T)(_Xr)(Y,U)\omega (Z,U)(_Xr)(Y,T).`$ Choose local frames $`\{V_a\}_{a=1}^{2n}`$, $`\{W_a\}_{a=1}^{2n}`$ on $`M`$ such that $`\omega (V_a,W_b)=\delta _{ab}`$. Substitute $`Y=V_a`$ and $`Z=W_a`$ in equation (1.4) and sum over $`a`$ to obtain $`0`$ $`=`$ $`2n(_Xr)(T,U)(_Tr)(X,U)(_Ur)(X,T)`$ (1.5) $`+\omega (X,T){\displaystyle \underset{a}{}}(_{W_a}r)(V_a,U)+\omega (X,U){\displaystyle \underset{a}{}}(_{W_a}r)(V_a,T).`$ If we cyclically permute $`X,T,U`$ in equation (1.5) and sum we get $$(2n2)_{X,T,U}(_Xr)(T,U)=0$$ (1.6) and since $`n2`$ we have $$_{X,T,U}(_Xr)(T,U)=0$$ (1.7) Using equation (1.7) in equation (1.5) gives $$(2n+1)(_Xr)(T,U)+\omega (X,T)\underset{a}{}(_{W_a}r)(V_a,U)+\omega (X,U)\underset{a}{}(_{W_a}r)(V_a,T)=0$$ which is of the desired form if $$u(X)=\underset{a}{}(_{W_a}r)(V_a,X).$$ Conversely, if one substitutes (1.2) into the covariant derivative of (1.1) and cyclically sums then one obtains $$_{X,Y,Z}(_XE)(Y,Z,T,U)=0.$$ Combining this with the second Bianchi identity, gives the second part of the Lemma. ###### Corollary A symplectic manifold with a symplectic connection whose curvature is of Ricci type is locally symmetric if and only if the $`1`$-form $`u`$, defined in the Lemma, vanishes. ###### Remark 1 It will be useful to have an equivalent form of formula (1.2). Denote by $`A`$ the linear endomorphism such that $$r(X,Y)=\omega (X,AY).$$ (1.8) The symmetry of $`r`$ is equivalent to saying that $`A`$ is in the Lie algebra of the symplectic group of $`\omega `$. Denote by $`\overline{u}`$ the vector field such that $$u=i(\overline{u})\omega $$ (1.9) then (1.2) is equivalent to $$_XA=\frac{1}{2n+1}(Xu+\overline{u}i(X)\omega ).$$ (1.10) ###### Lemma 2 Let $`(M,\omega )`$ be a symplectic manifold with a symplectic connection $``$ with Ricci-type curvature. Then, keeping the above notation, the following identities hold: 1. There is a function $`b`$ such that $$u=\frac{1+2n}{2(1+n)}\stackrel{(2)}{r}+b\omega $$ (1.11) where $`\stackrel{(2)}{r}`$ is the $`2`$-form $$\stackrel{(2)}{r}(X,Y)=\omega (X,A^2Y).$$ (1.12) 2. The differential of the function $`b`$ is given by $$db=\frac{1}{1+n}i(\overline{u})r.$$ (1.13) 3. The covariant differential of $`db`$ is given by $$db=\frac{1}{1+n}\left[\frac{1}{1+2n}uu\frac{1+2n}{2(1+n)}\stackrel{(3)}{r}+br\right].$$ (1.14) where $$\stackrel{(3)}{r}(X,Y)=\omega (X,A^3Y).$$ (1.15) Proof We can compute the action of the curvature on endomorphisms in two different ways. On the one hand it is $`R(X,Y)A`$ $`=`$ $`[R(X,Y),A]`$ $`=`$ $`R(X,Y)AAR(X,Y)`$ $`=`$ $`{\displaystyle \frac{1}{2(n+1)}}[X\omega (A^2Y,.)Y\omega (A^2X,.)`$ $`+A^2Y\omega (X,.)A^2X\omega (Y,.)].`$ On the other hand the curvature is of Ricci type so that (1.10) gives $$R(X,Y)A=\frac{1}{2n+1}[X_YuY_Xu+_Y\overline{u}\omega (X,.)_X\overline{u}\omega (Y,.)].$$ If we define an endomorphism $`B`$ of $`TM`$ by $$BY=\frac{2n+1}{2(n+1)}A^2Y+_Y\overline{u}$$ then equality of the two right hand sides yields $$X\omega (BY,.)Y\omega (BX,.)+BY\omega (X,.)BX\omega (Y,.)=0$$ whose only solution is $$B=bId.$$ This gives $$_Yu=\frac{2n+1}{2(n+1)}\omega (A^2Y,.)+b\omega (Y,.)$$ which is equation (1.11). Antisymmetrising (1.11) we get $$du=\frac{2n+1}{n+1}\stackrel{(2)}{r}+2b\omega .$$ Taking the exterior derivative gives $$0=\frac{2n+1}{n+1}d\stackrel{(2)}{r}+2db\omega .$$ But $`d\stackrel{(2)}{r}(X,Y,Z)`$ $`=`$ $`{\displaystyle _{X,Y,Z}}\omega (_XA^2Y,Z)`$ $`=`$ $`{\displaystyle \frac{1}{2n+1}}{\displaystyle _{X,Y,Z}}\omega (u(AY)X+\omega (X,AY)\overline{u},Z)`$ $`+\omega (u(Y)AX+\omega (X,Y)A\overline{u},Z)`$ $`=`$ $`{\displaystyle \frac{2}{2n+1}}{\displaystyle _{X,Y,Z}}\omega (X,Y)r(\overline{u},Z).`$ Substituting, $$[\frac{1}{n+1}r(\overline{u},.)+db]\omega =0.$$ and in dimension $`4`$ or higher this implies $$db=\frac{1}{n+1}r(\overline{u},.)$$ which is (1.13). Covariantly differentiating $`(_Xdb)(Y)`$ $`=`$ $`{\displaystyle \frac{1}{n+1}}\left[(_Xr)(\overline{u},Y)+r(_X\overline{u},Y)\right]`$ $`=`$ $`{\displaystyle \frac{1}{n+1}}\left[{\displaystyle \frac{1}{1+2n}}\omega (X,\overline{u})u(Y)+r({\displaystyle \frac{2n+1}{2(n+1)}}A^2X+bX,Y)\right]`$ which is (1.14). ## 2 Assume $`(M,\omega )`$ is a $`G`$-homogeneous symplectic manifold and $``$ is a $`G`$-invariant symplectic connection with Ricci-type curvature. If $``$ is not locally symmetric the $`G`$-invariant $`1`$-form $`u`$ is everywhere different from zero and the function $`b`$ is also $`G`$-invariant and hence constant. Putting these two facts into (1.13) we see that $`r`$ as a bilinear form is necessarily degenerate $$r(\overline{u},.)=0.$$ (2.16) Also (1.14) implies $$\frac{1}{2n+1}uu+\frac{1+2n}{2(1+n)}\stackrel{(3)}{r}br=0$$ (2.17) or equivalently $$\frac{1}{2n+1}\overline{u}u\frac{1+2n}{2(1+n)}A^3+bA=0.$$ (2.18) Applying $`A`$ to (2.18) and using (2.16) $$\frac{1+2n}{2(1+n)}A^4+bA^2=0.$$ (2.19) It follows that the only possible non-zero eigenvalues of $`A`$ are $`\pm \sqrt{\frac{2(1+n)}{1+2n}b}`$ and so are real or imaginary. ###### Lemma 3 If $`(M,\omega )`$ is a compact homogeneous symplectic manifold admitting a homogeneous symplectic connection $``$ with Ricci-type curvature which is not locally symmetric then $`b=0`$. Proof Recall that for any vector field $`X`$, Cartan’s identity gives $$divX\omega ^n\stackrel{\text{def}}{==}_X\omega ^n=nd\left((i(X)\omega )\omega ^{n1}\right)$$ and $$_X\omega ^n=(_X_X)\omega ^n=n(\omega (_.X,.)+\omega (.,_.X))\omega ^{n1}$$ so that $$divX=Trace[Z_ZX].$$ In particular, by (1.11) $$div\overline{u}=\frac{2n+1}{2(n+1)}TraceA^2+2nb.$$ $`G`$-invariance implies that $`div\overline{u}`$ is constant. But $`M`$ compact with no boundary implies $`_Mdiv\overline{u}\omega ^n=0`$ since the argument is exact; hence the constant is zero. Thus $$b=\frac{2n+1}{4n(n+1)}TraceA^2.$$ On the other hand, (2.19) implies that $`A^2`$ is a multiple of a projection and with $`A`$ symplectic this has even rank $`2p`$ say; using 2.16 we get $`2p<2n`$. Thus $$TraceA^2=\frac{4pb(1+n)}{1+2n}$$ so $$b=\frac{2n+1}{4n(n+1)}.\frac{4pb(1+n)}{1+2n}=\frac{p}{n}b$$ and hence $`b=0`$. It follows that $`A^4=0`$ so $`A`$ is nilpotent; moreover (2.18) tells us that $`A^3`$ has rank 1. ###### Lemma 4 Let $`(M,\omega )`$ be a $`4`$-dimensional homogeneous symplectic manifold admitting a homogeneous symplectic connection $``$ with Ricci-type curvature which is not locally symmetric. Let $`A`$ be the endomorphism associated to the Ricci tensor. Then 1. either $`A`$ is nilpotent, $`b0`$, $`A^2=0`$, and $`A`$ has rank 1 at any point; 2. or $`A`$ is nilpotent, $`b=0`$, and $`A^3`$ has rank 1 at any point; 3. or $`A`$ has a non zero eigenvalue so $`b0`$. Then $`A`$ admits a pair of non zero eigenvalues of opposite sign (real or imaginary) with multiplicity $`1`$ and $`0`$ is an eigenvalue of multiplicity $`2`$ at any point. Furthermore, $`A`$ has necessarily a nilpotent part. Proof The dimension – at any point $`xM`$ – of the generalised $`0`$ eigenspace of $`A`$ is even and non-zero, so is $`2`$ or $`4`$. If it is $`4`$ then $`A`$ is nilpotent and $`A^4=0`$ in dimension $`4`$. Thus, by (2.18), $`bA^2=0`$. If $`b0`$ then $`A^2=0`$ so, by (2.19), $`A`$ has rank 1 at any point. Otherwise $`b=0`$ and $`A^3`$ has rank 1 at any point. When the generalised 0 eigenspace $`V_0`$ is $`2`$-dimensional at any point, then $`\pm \sqrt{\frac{2(1+n)}{1+2n}b}`$ are eigenvalues with multiplicity $`1`$. Choose a globally defined vector field $`vV_0`$ so that $`\omega (v,\overline{u})=1`$. Set $`Av=p\overline{u}`$. Then $$_X(Av)=_X(p\overline{u})=(Xp)\overline{u}+p\left(\frac{1+2n}{2(1+n)}A^2X+bX\right)$$ but it is also equal to $$_X(Av)=(_XA)v+A(_Xv)=\frac{1}{1+2n}\left(Xu(v)+\overline{u}\omega (X,v)\right)+A(_Xv).$$ Observe that $`\omega (A^2X,\overline{u})=\omega (A(_Xv),\overline{u})=0`$, so that $$p\omega (bX,\overline{u})=\omega (\frac{1}{1+2n}u(v)X,\overline{u})=\frac{1}{1+2n}\omega (X,\overline{u}).$$ Hence $`pb=\frac{1}{1+2n}`$ which implies that $`p0`$. Thus $`A`$ has a nilpotent part. ## 3 We first prove Theorem 1 in the simply-connected case. It is standard that a compact simply-connected homogeneous symplectic manifold $`(M,\omega )`$ is symplectomorphic to a coadjoint orbit of a simply-connected compact semisimple Lie group $`G`$. Such a Lie group $`G`$ is a product of simple groups and the orbit is a product of orbits. We may throw away any factors where the orbit is zero dimensional as the remaining group will still act transitively. A $`G`$-invariant symplectic connection $``$ on such an orbit is compatible with the product structure. If the curvature of $``$ is of Ricci type, then it was shown in that the curvature is zero when $`(M,\omega ,)`$ is a product of more than one factor. But a non-trivial compact coadjoint orbit of a simple Lie group does not admit a flat connection since it has a non-zero Euler characteristic. It follows that we can assume $`G`$ is simple and $`(M,\omega )`$ is a coadjoint orbit $`(𝒪,\omega ^𝒪)`$ with its Kirillov–Kostant–Souriau symplectic structure and with an invariant symplectic connection $``$ with curvature of Ricci type. Further, the Euler characteristic of such an orbit is non-zero. If the vector field $`\overline{u}`$ were non-zero, then invariance would imply that it is everywhere non-zero and this cannot happen. Thus $`\overline{u}=0`$ and hence $``$ is locally symmetric ($`R=0`$). Pick a point $`\xi _0𝒪`$ and construct a symmetric symplectic triple $`(𝔩,\sigma ,\mathrm{\Omega })`$ as follows: Let $`𝔞=\{R_{\xi _0}(X,Y)End(T_{\xi _0}𝒪)\|X,YT_{\xi _0}𝒪\}`$ and $`𝔩=T_{\xi _0}𝒪𝔞`$. The bracket is defined by $`[X,Y]`$ $`=`$ $`R_{\xi _0}(X,Y),X,YT_{\xi _0}𝒪;`$ (3.20) $`[B,X]`$ $`=`$ $`BX,B𝔞,XT_{\xi _0}𝒪;`$ (3.21) $`[B,C]`$ $`=`$ $`BCCB,B,C𝔞,`$ (3.22) $`\sigma `$ by $$\sigma =Id_{T_{\xi _0}𝒪}Id_𝔞,$$ and $`\mathrm{\Omega }`$ by $$\mathrm{\Omega }(X+Y,X^{}+Y^{})=\omega _{\xi _0}(X,X^{}),X,X^{}T_{\xi _0},Y,Y^{}𝔞.$$ ###### Lemma 5 $`(𝔩,\sigma ,\mathrm{\Omega })`$ is, indeed, a symmetric symplectic triple. Proof There are two things to check to see that $`𝔩`$ is a Lie algebra. Firstly that the brackets defined above belong to $`𝔩`$. The only ones in doubt are the brackets of two elements of $`𝔞`$. But $`𝔞`$ is in fact the linear infinitesimal holonomy. This follows since the latter is spanned by the values of the curvature endomorphism and its covariant derivatives. The latter vanish by the local symmetry condition. The second thing to check is the Jacobi identity. Obviously this holds if all three elements are in $`𝔞`$ since this is a Lie algebra. If all three are in $`T_{\xi _0}𝒪`$ then $`[X,[Y,Z]]=R_{\xi _0}(Y,Z)X`$ and the Jacobi identity is satisfied for these elements by the first Bianchi identity. When one element is in $`T_{\xi _0}𝒪`$ and two in $`𝔞`$ we have $$[X,[B,C]]+[B,[C,X]]+[C,[X,B]]=[B,C]X+BCXCBX=0.$$ Finally, if two elements are in $`T_{\xi _0}𝒪`$ and one in $`𝔞`$ we have $`[X,[Y,B]]+[Y,[B,X]]+[B,[X,Y]]`$ $`=`$ $`R_{\xi _0}(X,BY)R_{\xi _0}(BX,Y)`$ $`+BR_{\xi _0}(X,Y)R_{\xi _0}(X,Y)B`$ $`=`$ $`(BR_{\xi _0})(X,Y)`$ where $`BR_{\xi _0}`$ denotes the natural action of the holonomy Lie algebra $`𝔞`$ on curvature tensors. But $`R=0`$ if and only if $`BR_{\xi _0}=0,B𝔞`$. The other two properties follow immediately from the definitions. If $`L`$ is the simply connected Lie group associated to $`𝔩`$ and $`K`$ the Lie subgroup associated to the subalgebra $`𝔞`$ then $`K`$ is the connected component of the fixed point set of the automorphism of $`L`$ induced by $`\sigma `$ and $`M_1=L/K`$ is a simply connected symmetric space. $`\mathrm{\Omega }`$ induces a symplectic form $`\omega _1`$ on $`M_1`$ which is parallel for the canonical connection $`_1`$. Consider the point $`\xi _0M=𝒪`$ and the point $`\xi _1=eKM_1`$. There is a linear isomorphism $`\varphi `$ from the tangent space $`T_{\xi _0}M`$ to the tangent space $`T_{\xi _1}M_1`$ so that $$\varphi \left(R_0(X,Y)Z\right)=R_1(\varphi X,\varphi Y)\varphi Z.$$ This implies \[4, p. 259, thm. 7.2\] that there exists an affine symplectic diffeomorphism $`\psi `$ of a neighbourhood $`U_0`$ of $`\xi _0`$ in $`𝒪`$ onto a neighbourhood $`U_1`$ of $`\xi _1`$ in $`M_1`$ such that $`\psi _{_{\xi _0}}=\varphi `$. Both $`(𝒪,\omega ^𝒪,)`$ and $`(M_1,\omega _1,_1)`$ are real analytic, as is $`\psi `$, and $`𝒪`$ is simply connected whilst $`_1`$ is complete. Hence \[4, p. 252, thm. 6.1\] there exists a unique affine map $`\stackrel{~}{\psi }:𝒪M_1`$ such that $`\stackrel{~}{\psi }\|_{U_0}=\psi `$. This map $`\stackrel{~}{\psi }`$ is symplectic since it is an analytic extension of the map $`\psi `$ which is symplectic. Symplectic maps are immersions, and hence local diffeomorphisms when the dimensions are equal as they are in this case. Hence $`\stackrel{~}{\psi }(𝒪)`$ is open. On the other hand, $`𝒪`$ is compact, so $`\stackrel{~}{\psi }(𝒪)`$ is compact and thus closed. Hence $`\stackrel{~}{\psi }`$ is surjective. It follows that $`M_1`$ is compact. From the preceding arguments we see that $`(M_1,\omega _1,_1)`$ is a compact simply connected symmetric symplectic space whose curvature is of Ricci type. The only such space is $`_n()`$ with a multiple of its standard Kähler form $`\omega _0`$ and the Levi-Civita connection $`_0`$ of the Fubini–Study metric. Since $`𝒪`$ and $`M_1`$ are both simply connected they are diffeomorphic and hence we have proved Theorem 1 in the simply connected case. Next we consider the case where $`M`$ has a finite fundamental group, $`(M,\omega )`$ is $`G`$-homogeneous symplectic with a $`G`$-invariant symplectic connection $``$ with curvature of Ricci type. Then the simply connected covering space $`\stackrel{~}{M}`$ is compact and carries such data $`\stackrel{~}{\omega },\stackrel{~}{}`$ for the simply connected covering group $`\stackrel{~}{G}`$. It follows that $`(\stackrel{~}{M},\stackrel{~}{\omega },\stackrel{~}{})`$ is diffeomorphic to $`(_n(),\omega _0,_0)`$ and hence that $`M`$ is diffeomorphic to $`_n()/\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is a discrete subgroup of $`PU(n+1)`$ acting properly discontinuously on $`_n()`$. But non-trivial elements of $`PU(n+1)`$ always have fixed points, so $`\mathrm{\Gamma }`$ must be trivial. This proves Theorem 1. ## 4 We now proceed to give the proof of Theorem 2 indicating along the way why we restrict ourselves to dimension 4 and why we only obtain a local result. Recall that when $`(M,\omega )`$ is homogeneous and admits a non-locally-symmetric invariant symplectic connection with Ricci-type curvature we have the non-zero vector field $`\overline{u}`$ and the Ricci endomorphism satisfies $$A\overline{u}=0,$$ $$\frac{1}{1+2n}\overline{u}u\frac{1+2n}{2(1+n)}A^3+bA=0,$$ (4.23) $$\frac{1+2n}{2(1+n)}A^4bA^2=0.$$ Furthermore, if $`M`$ is compact Lemma 3 tells us that $`b=0`$ so that $`A^4=0`$, and $`A^3`$ has rank 1: $$A^3=\frac{2(1+n)}{(1+2n)^2}\overline{u}u.$$ The $`1`$-form $`u`$ is everywhere non-zero so there is a globally defined vector field $`e_1`$ with $`u(e_1)`$ everywhere $`0`$. The vector fields $`e_1,e_2=Ae_1,e_3=A^2e_1,e_4=A^3e_1`$ form at each point $`xM`$ a basis of a $`4`$-dimensional subspace $`V_x`$ of the tangent space $`T_xM`$. Furthermore, by equation (4.23) $$e_4=\frac{2(1+n)}{(1+2n)^2}u(e_1)\overline{u}.$$ If we choose the vector field $`e_1`$ so that $`\omega (e_1,e_4)=ϵ`$ with $`ϵ^2=1`$, we get $`\frac{2(1+n)}{(1+2n)^2}(u(e_1))^2=ϵ`$ so that $`ϵ=1`$ and $`(u(e_1))^2=\frac{(1+2n)^2}{2(1+n)}`$ so that $`\overline{u}=u(e_1)e_4`$. Remark that we can always assume that $`\omega (e_1,e_2)=0`$ (by adding to $`e_1`$ a multiple of $`e_3`$). So the symplectic form restricted to $`V_x`$ writes in the chosen basis $$\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right).$$ The tangent space at each point $`xM`$ writes $$T_xM=V_xV_x^{}$$ where $`V_x^{}`$ is the $`\omega _x`$-orthogonal to $`V_x`$; it is stable under $`A`$ and, since $`A^3`$ has rank $`1`$, $`A^3\|_{V_x^{}}=0`$ but this is not enough to describe the behaviour of $`A`$ on $`V_x^{}`$. From now on, we restrict ourselves to the $`4`$-dimensional case. We define $`1`$-forms $`\alpha ,\beta ,\gamma ,\delta `$ such that $$_Xe_1=\alpha (X)e_1+\beta (X)e_2+\gamma (X)e_3+\delta (X)e_4.$$ Using formula (1.10) for $`A`$ (i.e. $`_XA=\frac{1}{2n+1}(Xu+u(e_1)e_4i(X)\omega )`$) we obtain $`_Xe_2`$ $`=`$ $`{\displaystyle \frac{u(e_1)}{2n+1}}X^1e_1+\left(\alpha (X){\displaystyle \frac{u(e_1)}{2n+1}}X^2\right)e_2`$ $`+\left(\beta (X){\displaystyle \frac{u(e_1)}{2n+1}}X^3\right)e_3+\left(\gamma (X){\displaystyle \frac{2u(e_1)}{2n+1}}X^4\right)e_4,`$ $`_Xe_3`$ $`=`$ $`{\displaystyle \frac{u(e_1)}{2n+1}}X^1e_2+\left(\alpha (X){\displaystyle \frac{u(e_1)}{2n+1}}X^2\right)e_3+\beta (X)e_4,`$ $`_Xe_4`$ $`=`$ $`{\displaystyle \frac{u(e_1)}{2n+1}}X^1e_3+\left(\alpha (X){\displaystyle \frac{2u(e_1)}{2n+1}}X^2\right)e_4.`$ On the other hand, formula (1.11) gives $$_Xe_4=\frac{u(e_1)}{2n+1}A^2X$$ so that $$\alpha (X)=\frac{u(e_1)}{2n+1}X^2.$$ The fact that $``$ is symplectic gives the additional condition that $$\gamma (X)=\frac{u(e_1)}{2n+1}X^4.$$ The connection is thus determined by the two $`1`$-forms $`\beta `$ and $`\delta `$. The vanishing of the torsion gives the expression of the brackets of the vector fields $`e_j`$. We can now compute the action of the curvature endomorphism on the vector fields $`e_j`$ in two different ways: using the formulas above or using the fact that the curvature is of Ricci type. This yields two identities $$d\beta =\frac{3}{2(n+1)}\omega +\frac{u(e_1)}{2n+1}e_{}^2\beta +\frac{1}{2(n+1)}e_{}^1e_{}^4,$$ $$d\delta =2\gamma \beta \frac{2}{2(n+1)}e_{}^3e_{}^4+2\alpha \delta $$ where the $`e_{}^j`$ are $`1`$-forms so that $`e_{}^j(e_k)=\delta _k^j`$ at each point. Using the formulas for the bracket of vector fields we have $$de_{}^3=\frac{2u(e_1)}{2n+1}e_{}^1e_{}^4\frac{u(e_1)}{2n+1}e_{}^2e_{}^3+e_{}^2\beta ,$$ and substituting $`e_{}^2\beta `$ in $`d\beta `$ yields $$d(\beta \frac{u(e_1)}{2n+1}e_{}^3)=\frac{2}{n+1}\omega $$ which is impossible on a compact manifold. This contradiction tells us that $`u`$ must vanish and hence that $``$ is locally symmetric. Acknowledgement The last author is grateful to the Mathematics Department of the University of Metz for its hospitality during part of this work.
warning/0006/astro-ph0006405.html	ar5iv	text	# Stellar evolution with rotation VI: The Eddington and Ω–limits, the rotational mass loss for OB and LBV stars ## 1 Introduction Recent models of stellar evolution with rotation (Meynet and Maeder MM00 (2000)) have shown that rotation heavily modifies all the model outputs for massive stars. In the course of the above mentioned work, it was realized that several basic points in the stellar physics need to be further clarified, since they may have some important consequences on the evolution. These points concern in particular the correct expression of the break–up velocities and the dependence of the mass loss rates $`\dot{M}`$ on the observed rotation velocities $`v`$. These problems are of great concern for the most luminous stars close to the Eddington limit, like OB stars, supergiants, LBV and WR stars. There is an interesting debate in recent litterature about what is the correct expression for the critical velocity and what is the dependence of the $`\dot{M}`$–rates on the rotation velocities $`v`$. The critical rotation velocity of a star is often written as $`v_{\mathrm{crit}}^2=\frac{GM}{R}(1\mathrm{\Gamma })`$, where $`\mathrm{\Gamma }=L/L_{\mathrm{Edd}}`$ is the ratio of the stellar luminosity to the Eddington luminosity. With this expression, Langer (La97 (1997), La98 (1998)) suggests that “no matter the rotation rate may be, it (the star) will arrive at critical rotation well before $`\mathrm{\Gamma }`$ = 1 is actually reached”. Consequently, Langer introduces the concept that the stars generally reach the break–up limit, i.e. the $`\mathrm{\Omega }`$–limit, earlier in evolution than the $`\mathrm{\Gamma }`$–limit. This point of view was disputed by Glatzel(Gla98 (1998)), who stressed that the $`\mathrm{\Omega }`$–limit is an artefact based on the disregard of gravity darkening and on the assumption of a uniform brightness over the surface of rotating stars. Glatzel concludes that the Eddington factor has no effect on the critical rotation. This problem needs to be further examined and this is what is done here. Another important issue is the dependence of the mass loss rates $`\dot{M}`$ on rotation velocities $`v`$ (cf. Maeder and Meynet MaMe00 (2000)). On one side, Langer (La97 (1997), La98 (1998)), Heger et al. (He00 (2000)), Meynet and Maeder (MM00 (2000)) are using a relation $`\dot{M}`$ vs. $`v`$ from Friend and Abbott (Fr86 (1986)) which formally leads to infinite mass loss rates at break–up velocities. On the other side, Owocki et al. (Ow96 (1996, 1998)), Owocki and Gayley (Ow97 (1997)) and Glatzel (Gla98 (1998)) show that even at extreme rotation the mass loss rates remain finite and do not differ too much from the case of zero–rotation. There also, some further analysis is needed. In a rotating star, the total gravity is the sum of the gravitational, centrifugal and radiative accelerations: $$g_{\mathrm{tot}}=g_{\mathrm{eff}}+g_{\mathrm{rad}}=g_{\mathrm{grav}}+g_{\mathrm{rot}}+g_{\mathrm{rad}}.$$ (1.1) For purpose of clarity, we adopt the following definitions: – We speak of the Eddington or $`\mathrm{\Gamma }`$–limit, when rotation effects can be neglected and $`g_{\mathrm{rad}}+g_{\mathrm{grav}}=0`$, which implies that $$\mathrm{\Gamma }=\frac{\kappa L}{4\pi cGM}\mathrm{\hspace{0.33em}1}.$$ (1.2) In that case $`L=L_{\mathrm{EDD}}=4\pi cGM/\kappa `$. The opacity $`\kappa `$ considered here is the total opacity, unless we specify it differently (cf. Sect. 4.2). –The break–up or $`\mathrm{\Omega }`$–limit is reached, for a star with an angular velocity $`\mathrm{\Omega }`$ at the surface, when the effective gravity $`g_{\mathrm{eff}}=g_{\mathrm{grav}}+g_{\mathrm{rot}}=0`$ and in addition when radiation pressure effects can be neglected. The $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit is reached when the total gravity $`g_{\mathrm{tot}}=0`$, with significant effects of both rotation and radiation. This is the general case, that we study here. It should lead to the two above cases in their respective limits. In Sect. 2, we examine the surface gravity, the Eddington factors and the limiting luminosity. In Sect. 3, the expression of the break–up velocities are considered, while in Sect. 4 we examine the mass loss rates. The equation of the surface with account of rotation and radiative acceleration is discussed in an Appendix. ## 2 Surface gravity, Eddington factors and limiting luminosity ### 2.1 The von Zeipel Theorem close to the $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit The von Zeipel theorem (vZ24 (1924)) expresses that the radiative flux $`F`$ at some colatitude $`\vartheta `$ in a rotating star is proportional to the local effective gravity $`g_{\mathrm{eff}}`$. In a previous work (Maeder Ma99 (1999)), we have generalized this theorem to the case of shellular rotation proposed by Zahn (Za92 (1992)). Shellular rotation results from strong horizontal turbulence which reduces the latitudinal dependence of rotation and makes the angular velocity $`\mathrm{\Omega }`$ constant on an isobar. Here, we shall consider the case of stars with shellular rotation, where the $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit may play a role. As shown by Langer (La97 (1997)), stars close to the Eddington limit tend to develop convection in the outer layers (cf. also Maeder Ma80 (1980)). However, in the outer layers the convective flux is generally negligible and the main transport mechanism is radiative transfer, a point also emphasized by Glatzel (Gla98 (1998)). As a numerical example, in a 60 M model (Meynet and Maeder MM00 (2000)) at the end of the MS phase with log L/L= 5.89 and log T<sub>eff</sub> =4.34 , the convective flux is negligible down to a fractional radius r/R of 0.85. In a 120 M with log L/L = 6.32 and log T<sub>eff</sub> =4.35 at the end of the MS, the convective flux is negligible down to r/R = 0.63. Thus, the basic condition to apply the von Zeipel theorem to stars close to the $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit is fulfilled, since the flux is essentially radiative. The expression of the flux $`F`$ (Maeder Ma99 (1999)) for a star with angular velocity $`\mathrm{\Omega }`$ on the isobaric stellar surface (cf. Appendix) is $$F=\frac{L(P)}{4\pi GM_{}}g_{\mathrm{eff}}[1+\zeta (\vartheta )]\mathrm{with}$$ (2.3) $$M_{}=M\left(1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}\right)\mathrm{and}$$ (2.4) $$\zeta (\vartheta )=\left[\left(1\frac{\chi _T}{\delta }\right)\mathrm{\Theta }+\frac{H_T}{\delta }\frac{d\mathrm{\Theta }}{dr}\right]P_2(\mathrm{cos}\vartheta ).$$ (2.5) There, $`\rho _\mathrm{m}`$ is the internal average density, $`\chi =4acT^3/(3\kappa \rho )`$ and $`\chi _T`$ is the partial derivative with respect to $`T`$. The quantity $`\mathrm{\Theta }`$ is defined by $`\mathrm{\Theta }=\frac{\stackrel{~}{\rho }}{\overline{\rho }}`$, i.e. the ratio of the horizontal density fluctuation to the average density on the isobar, which is given by $`\frac{\stackrel{~}{\rho }}{\overline{\rho }}=\frac{1}{3}\frac{r^2}{\overline{g}}\frac{d\mathrm{\Omega }^2}{dr}`$ where $`\overline{g}`$ is the average gravity on an isobar (cf. Zahn Za92 (1992)). One has the thermodynamic coefficients $`\delta =(\mathrm{ln}\rho /\mathrm{ln}T)_{P,\mu }`$, $`H_T`$ is the temperature scale height. The term $`\zeta (\vartheta )`$, which expresses the deviations of the von Zeipel theorem due to the baroclinicity of the star, is generally very small (cf. Maeder Ma99 (1999)). Let us emphasize that the flux is proportional to $`g_{\mathrm{eff}}`$ and not to $`g_{\mathrm{tot}}`$. This results from the fact that the equation of hydrostatic equilibrium is $`\frac{P}{\rho }=g_{\mathrm{eff}}`$. The effect of radiation pressure is already counted in the expression of $`P`$, which is the total pressure. We may call $`M_{}`$ the effective mass, i.e. the mass reduced by the centrifugal force. This is the complete form of the von Zeipel theorem in a differentially rotating star with shellular rotation, whether or not one is close to the Eddington limit. ### 2.2 Expressions of the gravity and of the local Eddington factor Let us express the total gravity at some colatitude $`\vartheta `$, taking into account the radiative acceleration $$g_{\mathrm{rad}}=\frac{1}{\rho }P_{\mathrm{rad}}=\frac{\kappa (\vartheta )F}{c},$$ (2.6) thus one has with Eq. (1.1), (2.3) and (2.4) $`g_{\mathrm{tot}}=g_{\mathrm{eff}}\left[1{\displaystyle \frac{\kappa (\vartheta )L(P)[1+\zeta (\vartheta )]}{4\pi cGM(1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}})}}\right].`$ (2.7) The rotation effects appear both in $`g_{\mathrm{eff}}`$ and in the term in brackets. When we write $`\kappa (\vartheta )`$, we mean that in a rotating star, the local T<sub>eff</sub> and gravity vary with latitude and so does the opacity.We may also consider the local limiting flux. The condition $`g_{\mathrm{tot}}=0`$ in Eq. (1.1) with Eq. (2.6) for $`g_{\mathrm{rad}}`$ allows us to define a limiting flux, $$F_{\mathrm{lim}}(\vartheta )=\frac{c}{\kappa (\vartheta )}g_{\mathrm{eff}}(\vartheta ).$$ (2.8) From that we may define the ratio $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$ of the actual flux $`F(\vartheta )`$ to the limiting local flux in a rotating star, $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )={\displaystyle \frac{F(\vartheta )}{F_{\mathrm{lim}}(\vartheta )}}={\displaystyle \frac{\kappa (\vartheta )L(P)[1+\zeta (\vartheta )]}{4\pi cGM\left(1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}\right)}}.`$ (2.9) As a matter of fact, $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$ is the local Eddington ratio. For zero rotation $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )=\mathrm{\Gamma }`$ as given by Eq. (1.2). Using relation (2.9), we may write the Eq. (2.7) for the total gravity as $$g_{\mathrm{tot}}=g_{\mathrm{eff}}\left[1\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )\right].$$ (2.10) This shows that the expression for the total acceleration in a rotating star is similar to the usual one, except that $`\mathrm{\Gamma }`$ is replaced by the local value $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$. Indeed, contrarily to expressions such as $`g_{\mathrm{tot}}=g_{\mathrm{eff}}\left(1\mathrm{\Gamma }\right)`$ often found in literature, we see that the appropriate Eddington factor (2.9) also depends on the angular velocity $`\mathrm{\Omega }`$ on the isobaric surface. From (2.9), we note that over the surface of a rotating star, which has a varying gravity and T<sub>eff</sub>, $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$ is the highest at the latitude where $`\kappa (\vartheta )`$ is the largest, (if we neglect the effects of $`\zeta (\vartheta )`$, which is justified in general). If the opacity increases with decreasing T as in hot stars, the opacity is the highest at the equator and there the limit $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )=1`$ may be reached first. Thus, it is to be stressed that if the limit $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )=1`$ happens to be met at the equator, it is not because $`g_{\mathrm{eff}}`$ is the lowest there, but because the opacity is the highest ! Indeed, both dependences in $`g_{\mathrm{eff}}`$ have cancelled each other in the ratio given by Eq. (2.9). ### 2.3 The luminosity at the $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit The $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit is reached, when the local Eddington ratio $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )=1`$ at some colatitude $`\vartheta `$. The condition $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )=1`$ allows us to define a limiting luminosity $`L_{\mathrm{\Omega }\mathrm{\Gamma }}`$ at the $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit, i.e. when both the effects of radiative acceleration and rotation are important. From (2.9) we have $$L_{\mathrm{\Omega }\mathrm{\Gamma }}=\frac{4\pi cGM}{\kappa (\vartheta )\left[1+\zeta (\vartheta )\right]}\left(1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}\right).$$ (2.11) It means that for a certain angular velocity $`\mathrm{\Omega }`$ on the isobaric surface, the maximum permitted luminosity of a star is reduced by rotation, with respect to the usual Eddington limit (cf. Sect. 1). This conclusion was also reached by Glatzel (Gla98 (1998)). In the above relation, $`\kappa (\vartheta )`$ is the largest value of the opacity on the surface of the rotating star. For O–type stars with photospheric opacities dominated by electron scattering, the opacity $`\kappa `$ is the same everywhere on the star. For the equation of the surface discussed in the Appendix, the maximum value of $`\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}=0.361`$, (with more digits it is 0.360747). ## 3 The break–up velocities We have seen above that rotation may be considered as reducing the maximum possible luminosity for a star. An alternative way to consider the problem is to ask the question: what happens to the break–up velocity for a star close to the Eddington limit ? Most authors (Langer La97 (1997, 1998, 1999); Lamers et al. Lam99 (1999); Heger et al. He00 (2000)) write this critical velocity $`v_{\mathrm{crit}}`$ like $$v_{\mathrm{crit}}^2=\frac{GM}{R}(1\mathrm{\Gamma }).$$ (3.12) This relation is true if we assume that the brightness of the rotating star is uniform over its surface, which is in contradiction with von Zeipel’s theorem. Surprisingly, some authors use this relation simultaneously with the von Zeipel theorem. Eq. (3.12), which we do not support, in agreement with Glatzel (Gla98 (1998)), implies that the break–up velocity is reduced by the proximity to the Eddington limit. The problem needs to be further studied carefully. The critical velocity is reached when somewhere on the star one has $`g_{\mathrm{tot}}=0`$, i.e. according to (2.10) $$g_{\mathrm{eff}}\left[1\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )\right]=0.$$ (3.13) This equation has two roots. The first one $`v_{\mathrm{crit},1}`$ is given by the usual condition $`g_{\mathrm{eff}}=0`$, which implies the equality $`\mathrm{\Omega }^2R_{\mathrm{eb}}^3/(GM)=1`$ at the equator (cf. Eq. A2 in Appendix). This corresponds to an equatorial critical velocity $$v_{\mathrm{crit},1}=\mathrm{\Omega }R_{\mathrm{eb}}=\left(\frac{2}{3}\frac{GM}{R_{\mathrm{pb}}}\right)^{\frac{1}{2}}.$$ (3.14) $`R_{\mathrm{eb}}`$ and $`R_{\mathrm{pb}}`$ are respectively the equatorial and polar radius at the break–up velocity and they obey to the surface equation. We notice that the critical velocity $`v_{\mathrm{crit},1}`$ is independent on the Eddington factor. To this extent, this is in agreement with Glatzel (Gla98 (1998)). The basic physical reason for this independence is quite clear: the radiative flux decreases at the equator, when the effective gravity decreases. Equation (3.13) has a second root, which is given by the condition $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$ = 1. As seen above, this condition will in general be met at the equator first. We thus have to search for the corresponding critical velocity $`v_{\mathrm{crit},2}`$ for a given value of the stellar luminosity. This second root has to be compared to the first one. For given values of M and L, the lowest of the two roots $`v_{\mathrm{crit},1}`$ and $`v_{\mathrm{crit},2}`$ is the significant one, since as soon as it will be reached the matter at the surface of the star is no longer bound. The condition $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$ = 1 gives, if we neglect $`\zeta (\vartheta )`$, $$\frac{\kappa (\vartheta )L(P)}{4\pi cGM}=1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}.$$ (3.15) Let us write $$\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}=\frac{16}{81}\omega ^2V^{}(\omega )\mathrm{with}$$ (3.16) $$V^{}(\omega )=\frac{V(\omega )}{\frac{4}{3}\pi R_{\mathrm{pb}}^3}\mathrm{and}\omega ^2=\frac{\mathrm{\Omega }^2R_{\mathrm{eb}}^3}{GM}.$$ (3.17) The quantity $`\omega `$ is the fraction of the angular velocity at the classical break–up given by Eq. (3.14). The density $`\rho _\mathrm{m}(\omega )=M/V(\omega )`$, where the stellar volume $`V(\omega )`$ depends on rotation. The quantity $`V^{}(\omega )`$ is the ratio of the actual volume of a star with rotation $`\omega `$ to the volume of a sphere of radius $`R_{\mathrm{pb}}`$. $`V^{}(\omega )`$ is obtained by the integration of the solutions of the surface equation (A6) for a given value of the parameter $`\omega `$. At break–up velocity $`v_{\mathrm{crit},1}`$ , the value of $`V^{}(\omega )`$ = 1.829, which gives the maximum value of $`\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}`$ = 0.361. If we call $`\mathrm{\Gamma }_{\mathrm{max}}`$ the maximum Eddington ratio $`\kappa (\vartheta )L(P)/(4\pi cGM)`$ over the surface (in general at the equator), Eq. (3.15) can thus be written $$\frac{16}{81}\omega ^2V^{}(\omega )=\mathrm{\hspace{0.33em}1}\mathrm{\Gamma }_{\mathrm{max}}.$$ (3.18) For a given value of $`\mathrm{\Gamma }_{\mathrm{max}}`$, one must search the value of $`\omega `$ which satisfies this equation. This is easily obtained by solving numerically the surface equation. For a given large enough $`\mathrm{\Gamma }_{\mathrm{max}}`$ (i.e. larger than 0.639), the obtained $`\omega `$–value is lower than 1, and this implies a corresponding critical velocity $`v_{\mathrm{crit},2}`$ given by $`v_{\mathrm{crit},2}^2=\mathrm{\Omega }^2R_\mathrm{e}^2(\omega )={\displaystyle \frac{81}{16}}{\displaystyle \frac{1\mathrm{\Gamma }_{\mathrm{max}}}{V^{}(\omega )}}{\displaystyle \frac{GM}{R_{\mathrm{eb}}^3}}R_\mathrm{e}^2(\omega )=`$ $`{\displaystyle \frac{9}{4}}v_{\mathrm{crit},1}^2{\displaystyle \frac{1\mathrm{\Gamma }_{\mathrm{max}}}{V^{}(\omega )}}{\displaystyle \frac{R_\mathrm{e}^2(\omega )}{R_{\mathrm{pb}}^2}},`$ (3.19) where R$`{}_{\mathrm{e}}{}^{}(\omega )`$ is the equatorial radius for a given value of the rotation parameter $`\omega `$. Fig. 1 illustrates the results. We notice that the second root $`v_{\mathrm{crit},2}`$, expressed as a fraction of the first root $`v_{\mathrm{crit},1}`$ given by (3.14), tends towards zero when the Eddington factor $`\mathrm{\Gamma }_{\mathrm{max}}`$ tends towards 1. The first root defined by (3.14) can still be written, but the second root defined by the condition $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )=1`$ is met first. This reduction of the critical velocity with respect to the classical expression only occurs for Eddington factors larger than 0.639, since the maximum value of $`\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}`$ is 0.361 (cf. Eq. 3.15). This second root results physically from both effects of rotation and radiation: for $`\mathrm{\Gamma }_{\mathrm{max}}>0.639`$, a zero value of $`g_{\mathrm{tot}}`$ can be achieved for non–extreme rotations. This enters through the reduction due to rotation of the effective mass M, which is the significant mass in the local Eddington factor $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$. For $`\mathrm{\Gamma }_{\mathrm{max}}<0.639`$, Eq. (3.18) has no solution. Physically this means that when the star is sufficiently far from the Eddington limit, the reduction of the effective mass M by rotation is not sufficient to bring $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$ to 1. In that case, Eq. (3.13) has only one root given by the classical Eq. (3.14). We hope that these results clarify the debate between Glatzel (Gla98 (1998)) and Langer(La97 (1997)). On one hand, we see that the claim by Langer that stars close to the Eddington limit have a lower rotation limit is correct, even if the Eq. (3.12) by Langer is not the right one. On the other side, Glatzel has claimed that the Eddington factor does not affect the break–up velocity, we see that this is true in general for most stars for the reasons given above, however for $`\mathrm{\Gamma }_{\mathrm{max}}0.639`$ this statement does not apply. The moment when stars reach their critical velocities is far from being an academic one, since when this occurs large mass loss enhancements may result, a point which is examined below. ## 4 The mass loss rates as a function of $`\mathrm{\Omega }`$ and $`\mathrm{\Gamma }`$ ### 4.1 Present context The effects of rotation on the mass loss rates have been studied both observationally and theoretically. Observationally, very large changes of the $`\dot{M}`$–rates, i.e. up to 2–3 orders of a magnitude, were suggested by Vardya (Var85 (1985)). However, Nieuwenhuijzen and de Jager (Nieu88 (1988)) claimed with reason that the correlation found by Vardya was largely the reflect of the distribution of the mass loss rates and rotational velocities over the HR diagram. When disentangling the various effects, Nieuwenhuijzen and de Jager found much smaller effects of rotation. However, they noticed that the $`\dot{M}`$–rates of the Be–stars are larger by about a factor 10<sup>2</sup>. Since Be–stars are fast rotating stars, we may wonder whether the effects of rotation on the mass loss rates are really so negligible, as these last authors considered them. Certainly further observational studies are also required. On the theoretical side, Pauldrach et al. (Pau86 (1986)), Friend and Abbott (Fr86 (1986)) find only a moderate increase of the $`\dot{M}`$–rates, of about 30 % for $`v=350`$ km/s. Friend and Abbott find an increase of the $`\dot{M}`$–rates which can be fitted by the relation (Langer La98 (1998); Heger and Langer He98 (1998)) $$\dot{M}(v)=\dot{M}(v=0)\left(\frac{1}{1\frac{v}{v_{\mathrm{crit}}}}\right)^\xi $$ (4.20) with $`\xi `$ = 0.43. This expression, often used in evolutionary models, is based on wind models which do not account for the von Zeipel theorem. We notice that Eq. (4.20) diverges at break–up, while as shown by Glatzel (Gla98 (1998)) and Owocki et al. (Ow96 (1996, 1998)), the stellar mass loss rates should not diverge at the $`\mathrm{\Omega }`$–limit, (however see Sect. 4.3). When a proper account of the gravity darkening is made, there are two main terms contributing to the anisotropic mass loss rates from a rotating star (cf. Maeder Ma99 (1999)). 1) The “g<sub>eff</sub>–effect” which favours polar ejection, since the polar caps of a rotating star are hotter. 2) The “opacity or $`\kappa `$–effect”, which may favour an equatorial ejection, when the opacity is large enough at the equator due to the lower T<sub>eff</sub>. In O–type stars, since opacity is due mainly to the T–independent electron scattering, the g<sub>eff</sub>–effect is likely to dominate, raising a fast highly ionized polar wind. In B– and later type stars, the opacity effect may favour a dense equatorial wind and ring formation, with low terminal velocities and low ionization. Recently Petrenz and Puls (PP00 (2000)) have constructed 2–D models of line driven winds for rotating O–type stars. They found that the mass loss from hot star is essentially polar due to the $`g_{\mathrm{eff}}`$–effect. Quantitatively their results differ very little from previous works by Pauldrach et al. (Pau86 (1986)). ### 4.2 Expression of the $`\dot{M}`$–rates as a function of $`\mathrm{\Omega }`$ and $`\mathrm{\Gamma }`$ It is worth to further examine the consequences of the above results on the dependence of the mass loss rates on rotation. According to the radiative wind theory (cf. Castor et al. Ca75 (1975); Pauldrach et al. 1986; Kudritzki et al. Ku89 (1989); Puls et al. Pu96 (1996)), we may write the mass loss fluxes $`\mathrm{\Delta }\dot{M}/\mathrm{\Delta }\sigma `$ by surface elements $`\mathrm{\Delta }\sigma `$ $$\frac{\mathrm{\Delta }\dot{M}(\vartheta )}{\mathrm{\Delta }\sigma }\left(k\alpha \right)^{1/\alpha }\left(\frac{1\alpha }{\alpha }\right)^{\frac{1\alpha }{\alpha }}F(\vartheta )^{1/\alpha }g_{tot}^{1\frac{1}{\alpha }}(\vartheta )$$ (4.21) where $`k`$ and $`\alpha `$ are the force multiplier parameters. At some temperatures, the ionisation equilibrium of the stellar wind is changing abruptly and so does the opacity of the plasma. Consequently, the values of the force multipliers undergo rapid transitions for certain values of $`T_{\mathrm{eff}}`$, particularly at 21 000 K and maybe also at 10 000 K (cf. Lamers et al. Lam95 (1995); Lamers Lam97 (1997)). Such fast transitions of the wind properties are called by Lamers a bi–stability of the stellar winds, since near the transition limit the wind can exist in two states. There are both empirical and theoretical determinations of $`\alpha `$, however they lead to rather different values (cf. Lamers et al. Lam95 (1995)). The empirical ones, based on the values of the observed terminal velocities, are in general smaller than the theoretical estimates. As empirical values, Lamers et al. (Lam95 (1995)) obtain, for example, $`\alpha `$ = 0.52 for $`4.70\mathrm{log}T_{\mathrm{eff}}4.35`$, (type B1.5 or earlier); $`\alpha `$ = 0.24, 0.21, 0.17, 0.15 for $`\mathrm{log}T_{\mathrm{eff}}`$ = 4.30 (type B2.5), 4.20 (B5), 4.00 (B9.5), 3.90 (A7) respectively. These transitions may produce jumps in the mass loss rates, with the high rates on the low side of the transition. As $`T_{\mathrm{eff}}`$ is decreasing from the pole to the equator, one may thus expect, on the surface of a fast rotating star of type B or later, the occurence of some bi–stability limits and the corresponding variations of the force multipliers and of the mass loss rates. In a star, where a bi–stability limit is crossed at some latitude, a steep increase of the mass flux will happen between this latitude and the equator, possibly leading to a huge equatorial ring. With the expressions of the flux (2.3) and of $`g_{\mathrm{tot}}`$ (2.10), we get for the mass flux $`{\displaystyle \frac{\mathrm{\Delta }\dot{M}(\vartheta )}{\mathrm{\Delta }\sigma }}A\left[{\displaystyle \frac{L(P)}{4\pi GM_{}(P)}}\right]^{\frac{1}{\alpha }}{\displaystyle \frac{g_{\mathrm{eff}}[1+\zeta (\vartheta ]^{\frac{1}{\alpha }}}{(1\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta ))^{\frac{1}{\alpha }1}}}`$ $`\mathrm{with}A=\left(k\alpha \right)^{\frac{1}{\alpha }}\left({\displaystyle \frac{1\alpha }{\alpha }}\right)^{\frac{1\alpha }{\alpha }},`$ (4.22) with M given by Eq. (2.4). We notice the gravity–effect, which favours mass loss at the pole, where the total gravity is higher, and the $`\kappa `$–effect which favours high mass loss where $`\alpha `$ is small. The proximity to the Eddington limit will enhance the mass flux due to the term $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$, while rotation enhances the mass flux through both the terms M and $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$. In the theory of radiatively driven winds, the total opacity at a given optical depth is expressed with the force multipliers in terms of the electron scattering opacity $`\kappa _{\mathrm{es}}`$. This means that in Eq. (4.22), $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$ is just $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )={\displaystyle \frac{\kappa _{\mathrm{es}}L(P)[1+\zeta (\vartheta )]}{4\pi cGM\left(1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}\right)}}.`$ (4.23) The dependence on latitude of $`\mathrm{\Gamma }_\mathrm{\Omega }(\vartheta )`$ would only come through the term $`\zeta (\vartheta )`$. ### 4.3 Dependence of the global mass loss rates on rotation Let us estimate how the global mass loss rates depend on the rotation velocities and on the proximity of the $`\mathrm{\Omega }\mathrm{\Gamma }`$ limit. For that we henceforth neglect the small corrective term $`\zeta (\vartheta )`$ in the expression of the flux. We have, if $`\mathrm{\Sigma }(\omega )`$ is the total surface $`{\displaystyle \frac{\dot{M}}{\mathrm{\Sigma }(\omega )}}A\left[{\displaystyle \frac{L(P)}{4\pi GM_{}}}\right]^{\frac{1}{\alpha }}\overline{\left({\displaystyle \frac{g_{\mathrm{eff}}}{(1\mathrm{\Gamma }_\mathrm{\Omega })^{\frac{1}{\alpha }1}}}\right)}`$ $`A\left[{\displaystyle \frac{L(P)}{4\pi GM_{}}}\right]^{\frac{1}{\alpha }}{\displaystyle \frac{\overline{g_{\mathrm{eff}}}}{(1\mathrm{\Gamma }_\mathrm{\Omega })^{\frac{1}{\alpha }1}}},`$ (4.24) since $`\mathrm{\Gamma }_\mathrm{\Omega }`$ is independent on $`\vartheta `$, and with appropriate $`\alpha `$– and $`A`$–values. For the average effective gravity, we have $$\overline{g_{\mathrm{eff}}}=\frac{g_{\mathrm{eff}}d\sigma }{\mathrm{\Sigma }(\omega )}=\frac{4\pi GM_{}}{\mathrm{\Sigma }(\omega )},$$ (4.25) after integration over the stellar surface which is an isobar (cf. Appendix). This leads to the following expression for the total mass loss rate from the star $$\dot{M}\frac{AL(P)^{\frac{1}{\alpha }}}{\left(4\pi GM\right)^{\frac{1}{\alpha }1}\left[1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}\right]^{\frac{1}{\alpha }1}(1\mathrm{\Gamma }_\mathrm{\Omega })^{\frac{1}{\alpha }1}}.$$ (4.26) This relation expresses how the total mass loss rate from a star depends on mass, luminosity, Eddington factor and rotation, (see Fig. A1 for a simple expression of the rotation parameter). If we omit rotation, Eq. (4.26) is identical to the typical relations used in literature (cf. Pauldrach et al. Pau86 (1986); Lamers Lam97 (1997)). The amplitude of the effects very much depends on the opacity and in particular on the value of the force multiplier $`\alpha `$. Let us consider a rotating star with angular velocity $`\mathrm{\Omega }`$ and a non–rotating star of the same mass $`M`$ at about the same location in the HR diagram. The ratio of their mass loss rates can be written, $$\frac{\dot{M}(\mathrm{\Omega })}{\dot{M}(0)}=\frac{\left(1\mathrm{\Gamma }\right)^{\frac{1}{\alpha }1}}{\left[1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}\right]^{\frac{1}{\alpha }1}(1\mathrm{\Gamma }_\mathrm{\Omega })^{\frac{1}{\alpha }1}},$$ (4.27) where $`\mathrm{\Gamma }`$ is the Eddington ratio corresponding to electron scattering opacity for the non–rotating star. From Eq. (4.23), we have the relation $$\mathrm{\Gamma }_\mathrm{\Omega }=\frac{\mathrm{\Gamma }}{1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}}.$$ (4.28) We get finally $$\frac{\dot{M}(\mathrm{\Omega })}{\dot{M}(0)}=\frac{\left(1\mathrm{\Gamma }\right)^{\frac{1}{\alpha }1}}{\left[1\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}\mathrm{\Gamma }\right]^{\frac{1}{\alpha }1}}.$$ (4.29) If $`\mathrm{\Omega }=0`$, this ratio is of course equal to 1. This ratio, which is the main result of this work, can also be expressed with the ratio $`v/v_{\mathrm{crit},1}`$ of the rotational velocity $`v`$ to the critical velocity given by the usual Eq. (3.14). From Eq. (A7) in the Appendix, we have $`\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}\frac{4}{9}\frac{v^2}{v_{\mathrm{crit},1}^2}`$ over a large range of values (cf. Fig. A1), thus one can write $$\frac{\dot{M}(\mathrm{\Omega })}{\dot{M}(0)}\frac{\left(1\mathrm{\Gamma }\right)^{\frac{1}{\alpha }1}}{\left[1\frac{4}{9}(\frac{v}{v_{\mathrm{crit},1}})^2\mathrm{\Gamma }\right]^{\frac{1}{\alpha }1}}.$$ (4.30) For a star with a small Eddington factor, it simplifies to $$\frac{\dot{M}(\mathrm{\Omega })}{\dot{M}(0)}\frac{1}{\left[1\frac{4}{9}(\frac{v}{v_{\mathrm{crit},1}})^2\right]^{\frac{1}{\alpha }1}}.$$ (4.31) Equation (4.31) shows that the effects of rotation on the $`\dot{M}`$–rates remain moderate in general. This is in agreement with the results by Owocki et al. (Ow96 (1996)), by Owocki and Gayley (Ow97 (1997)) and by Glatzel (Gla98 (1998)), and also with more elaborate non–LTE 1–D and 2–D models by Pauldrach et al. (Pau86 (1986)), Petrenz and Puls (PP00 (2000)). However, this is only true for stars far enough from the Eddington limit. When $`\mathrm{\Gamma }`$ is significant, rotation may drastically increase the mass loss rates as shown by (4.29) or (4.30). This is particularly the case for low values of $`\alpha `$, i.e. for stars with log T$`{}_{\mathrm{eff}}{}^{}`$ 4.30. In the extreme cases where $`\mathrm{\Gamma }>0.639`$, a moderate rotation may even make the denominator of (4.29) or (4.30) to vanish, thus leading to extreme mass loss. Table 1 shows some numerical results based on Eq. (4.29). For different initial stellar masses in the Geneva models at Z =0.02 (Schaller et al. Scha92 (1992)), the $`\mathrm{\Gamma }`$ factors at the end of the Main Sequence (MS) phase are given, as well as the predicted ratios $`\dot{M}(\mathrm{\Omega })/\dot{M}(0)`$ of the mass loss rates for a star at break–up rotation and for a non–rotating star of the same mass and luminosity. These ratios are given at $`\mathrm{log}T_{\mathrm{eff}}4.35`$ ($`\alpha =0.52`$), at $`\mathrm{log}T_{\mathrm{eff}}`$ = 4.30 ($`\alpha =0.24`$), at $`\mathrm{log}T_{\mathrm{eff}}`$ = 4.00 ($`\alpha =0.17`$) and at $`\mathrm{log}T_{\mathrm{eff}}`$ = 3.90 ($`\alpha =0.15`$) for the same value of $`\mathrm{\Gamma }`$. This covers the range of the typical $`T_{\mathrm{eff}}`$ of OB and LBV stars, the differences with T<sub>eff</sub> result from the differences in the $`\alpha `$–parameter. The indication $`\mathrm{}`$ in Table 1 means that the bracket term in (4.29) or (4.30) may vanish at maximum rotation, which leads to extreme mass outflows. The ratios $`\dot{M}(\mathrm{\Omega })/\dot{M}(0)`$ keep quite moderate even at extreme rotation for MS stars up to 40 M, while for MS stars above 60 M they can become very large. These ratios may also be very large for B–type supergiants and LBV stars. In particular, we notice that for stars close to the Humphreys–Davidson limit the ratios $`\dot{M}(\mathrm{\Omega })/\dot{M}(0)`$ may diverge. Such stars are typically at the $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit. For log T$`{}_{\mathrm{eff}}{}^{}4.30`$, the force multiplier $`\alpha `$ is also very small, which favours extreme mass loss. On the whole, it is striking that the domain where $`\dot{M}(\mathrm{\Omega })/\dot{M}(0)`$ has the possibility to diverge so closely corresponds to the observed domain of LBV stars. The present results will enable us to better specify the changes of the $`\dot{M}`$ rates in massive star models. ## 5 Conclusion We conclude that the concept of an $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit reached during the evolution of the most massive stars is not an artefact, but the existence of this limit is confirmed by consistent developments based on the von Zeipel theorem. However, we emphasize that the expression currently used for the critical velocity is not correct. We have also clarified the dependence of the mass loss rates on the rotation velocities in the general case. We can make the following remarks on the various limits: –1. The $`\mathrm{\Gamma }`$–limit : The mass loss rates grow steeply as the Eddington limit is approached, even in absence of rotation. This is a well known result of the classical wind theory. –2. The $`\mathrm{\Omega }`$–limit : We see that the case of only rotational effects does not apply for O–type stars and even for the early B–type stars, since they always have a significant $`\mathrm{\Gamma }`$–value. Only for spectral types later than B3 on the MS, the $`\mathrm{\Gamma }`$ term can be ignored. In the framework of the radiative wind theory, the growth of the mass loss rates remains limited. –3. The $`\mathrm{\Omega }\mathrm{\Gamma }`$–limit : This general case is met for rotating OB stars, LBV stars, supergiants and Wolf–Rayet stars, because both $`\mathrm{\Gamma }`$ and rotation are important. The bracket in (4.29) is reduced by rotation and by the proximity to the Eddington limit. As shown by Table 1, both effects produce steep enhancements of the mass loss rates, especially for lower $`T_{\mathrm{eff}}`$ since $`\alpha `$ is lower. This may explain the very large mass loss rates for LBV stars, blue and yellow supergiants (cf. de Jager et al. deJ88 (1988)). If the ratio $`\mathrm{\Gamma }=\frac{\kappa _{\mathrm{es}}L}{4\pi cGM}`$ is bigger than 0.639, the break–up limit is reached for reduced rotation velocities, as illustrated by Fig. 1. Then, extremely high mass loss rates may occur, a situation likely corresponding to the case of the LBV stars and maybe also to some WR stars. Some words of caution are necessary. The $`\dot{M}`$–rates given here are the values predicted in the framework of the radiative wind theory. It is probable that close to break–up several other effects not included here may intervene, such as important horizontal fluxes, formally vanishing $`T`$– and $`P`$– gradients, instabilities, etc… Also, we may point out that if the flux vanishes, the radiative wind theory should not apply. Thus, for the detailed physics of the break–up, more complex analyses are certainly needed. Finally, we note that it was generally believed that in addition to $`L`$ and $`T_{\mathrm{eff}}`$, the mass loss rates only depend on metallicity Z. We see here another dependence which is quite significant and may introduce some scatter in the values of the $`\dot{M}`$–rates. Thus, we may expect that for a given initial mass the evolution is very different according to rotation, due to both rotational mixing, meridional circulation (Maeder and Zahn MZ98 (1998)) and to the induced differences in the mass loss rates. ## Appendix A The equation of the stellar surface in a rotating star with high radiation pressure Shellular rotation, with an angular velocity $`\mathrm{\Omega }`$ constant on horizontal surfaces, was proposed by Zahn (Za92 (1992)). This rotation law results from strong horizontal geostrophic–like turbulence which homogeneizes rotation on the horizontal surfaces. As noted by Meynet and Maeder (MM97 (1997)), the isobars for shellular rotation are identical to the equipotentials of the conservative case, which are $$\mathrm{\Psi }=\frac{GM}{r(\vartheta )}+\frac{1}{2}\mathrm{\Omega }^2r^2(\vartheta )\mathrm{sin}^2\vartheta =\mathrm{const}$$ (A.32) The components of the effective gravity are $`g_{\mathrm{eff},r}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }}{r}}+\mathrm{\Omega }^2r\mathrm{sin}\vartheta `$ $`g_{\mathrm{eff},\vartheta }`$ $`=`$ $`{\displaystyle \frac{1}{r}}{\displaystyle \frac{\mathrm{\Phi }}{\vartheta }}+\mathrm{\Omega }^2r\mathrm{sin}\vartheta \mathrm{cos}\vartheta `$ (A.33) where $`\mathrm{\Phi }=GM/r`$. In vectorial form, one can write $$P=\rho g_{\mathrm{eff}}=\rho (\mathrm{\Psi }r^2\mathrm{sin}^2\vartheta \mathrm{\Omega }\mathrm{\Omega }).$$ (A.34) This equation is interesting: it shows that if $`\mathrm{\Omega }`$ is constant on isobars, $`\mathrm{\Psi }`$ is also constant on isobars. Moreover, for a motion $`ds`$, the equation of the surface must satisfy $$g_{\mathrm{tot}}ds=0.$$ (A.35) Since we have $$g_{\mathrm{tot}}=g_{\mathrm{eff}}\left[1\mathrm{\Gamma }_\mathrm{\Omega }\right],$$ (A.36) this implies that the surface is perpendicular to $`g_{\mathrm{tot}}`$, thus to $`g_{\mathrm{eff}}`$ and to the P—gradient. The surface is an isobar, also if the radiation pressure is important. The equation of the surface can therefore be represented by Eq. (A1) and the procedure to calculate is quite conventional, i.e. $$\frac{1}{x}+\frac{4}{27}\omega ^2x^2\mathrm{sin}^2\vartheta =1$$ (A.37) with $`x=\frac{R}{R_\mathrm{p}}`$. At break–up, the equatorial radius $`R_{\mathrm{eb}}`$ equals 1.5 times the polar radius $`R_{\mathrm{pb}}`$. In the above demonstration, there is no need of the assumption $`\mathrm{\Omega }=\mathrm{\Omega }(r)`$, we just need the assumption that $`\mathrm{\Omega }`$ is constant on horizontal surface, which is less restrictive. This is true, whether the radiative acceleration is important or not. The critical rotation parameter naturally appearing in this work was the ratio $`\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}`$. Account must be given to the change of the average density of the star with rotation. We can express this ratio by (3.16) or in term of the actual rotational velocity $`v`$ and of the critical velocity $`v_{\mathrm{crit},1}`$ (3.14), $$\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}=\frac{4}{9}\frac{v^2}{v_{\mathrm{crit},1}^2}V^{}(\omega )\frac{R_{\mathrm{pb}}^2}{R_\mathrm{e}^2(\omega )}.$$ (A.38) The relation between $`\frac{\mathrm{\Omega }^2}{2\pi G\rho _\mathrm{m}}`$ and the ratio $`\frac{v^2}{v_{\mathrm{crit},1}^2}`$ is illustrated in Fig. A1. The product $`V^{}(\omega )\frac{R_{\mathrm{pb}}^2}{R_\mathrm{e}^2(\omega )}`$ has a limited range of variation, being equal to 1 for zero rotation and to 0.813 at break–up velocity. This means that for a crude estimate at low or moderate velocities, one may just ignore this product in Eq. (A7). ###### Acknowledgements. We express our thanks to Dr. Joachim Puls for very valuable remarks during this work.
warning/0006/hep-ph0006335.html	ar5iv	text	# Institute of Experimental Physics Warsaw University hep-ph/0006335 Leptoquark Searches at Future Colliders ## 1 Introduction New result on atomic parity violation (APV) in Cesium and unitarity of the CKM matrix, as well as recent LEP2 hadronic cross-section measurements indicate possible deviations from the Standard Model predictions. Exchange of leptoquark type objects with masses above 250GeV has been proposed as a possible explanation for these effects . If the observed signal is real it should become clearly visible in future colliders. The aim of the present analysis is to compare leptoquark search limits expected from different experiments. The Buchmüller-Rückl-Wyler model used in this analysis is described in section 2. Results from the global analysis of available experimental data and the possible leptoquark signal are briefly summarized in section 3. Parameters of existing and future colliders considered in this analysis are presented in section 4. In section 5 expected limits on the leptoquark mass to the coupling ratio are calculated in the limit of very high leptoquark masses, using the contact interaction approximation. Limits expected from precise measurements of the Drell-Yan lepton pair production at the Tevatron and LHC, hadronic cross-section measurements at TESLA $`e^+e^{}`$ and high-$`Q^2`$ NC DIS at HERA, THERA and TESLA ($`e^+e^{}`$ and $`e\gamma `$) are compared with the existing limits. In sections 6 and 7 expected limits from direct leptoquark production are considered. Leptoquark pair production is studied for the Tevatron, LHC and TESLA ($`e^+e^{}`$ and $`\gamma \gamma `$ scattering) whereas single leptoquark production is considered for HERA, THERA and TESLA ($`e^+e^{}`$ and $`e\gamma `$ scattering). Results expected from different experiments are compared in section 8. ## 2 Leptoquark models In this paper a general classification of leptoquark states proposed by Buchmüller, Rückl and Wyler will be used. The Buchmüller-Rückl-Wyler (BRW) model is based on the assumption that new interactions should respect the $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ symmetry of the Standard Model. In addition leptoquark couplings are assumed to be family diagonal (to avoid FCNC processes) and to conserve lepton and baryon numbers (to avoid rapid proton decay). Taking into account very strong bounds from rare decays it is also assumed that leptoquarks couple either to left- or to right-handed leptons. With all these assumptions there are 14 possible states (isospin singlets or multiplets) of scalar and vector leptoquarks. Table 1 lists these states according to the so-called Aachen notation . An S(V) denotes a scalar(vector) leptoquark and the subscript denotes the weak isospin. When the leptoquark can couple to both right- and left-handed leptons, an additional superscript indicates the lepton chirality. A tilde is introduced to differentiate between leptoquarks with different hypercharge. Listed in Table 1 are the leptoquark fermion number F, electric charge Q, and the branching ratio to an electron-quark pair (or electron-antiquark pair), $`\beta `$. The leptoquark branching fractions are predicted by the BRW model and are either 1, $`\frac{1}{2}`$ or 0. For a given electron-quark branching ratio $`\beta `$, the branching ratio to the neutrino-quark is by definition $`(1\beta )`$. Also included in Table 1 are the flavours and chiralities of the lepton-quark pairs coupling to a given leptoquark type. In three cases the squark flavours (in supersymmetric theories with broken R-parity) with corresponding couplings are also indicated. Present analysis takes into account only leptoquarks which couple to the first-generation leptons ($`e`$, $`\nu _e`$) and first-generation quarks ($`u`$, $`d`$), as most of the existing experimental data constrain this type of couplings. Second- and third-generation leptoquarks as well as generation-mixing leptoquarks will not be considered in this paper. It is also assumed that one of the leptoquark types gives the dominant contribution, as compared with other leptoquark states and that the interference between different leptoquark states can be neglected. Using this simplifying assumption, different leptoquark types can be considered separately. Finally, it is assumed that different leptoquark states within isospin doublets and triplets have the same mass. ## 3 Current limits from global analysis In a recent paper available data from HERA, LEP and the Tevatron, as well as from low energy experiments are used to constrain the Yukawa couplings $`\lambda _{LQ}`$ and masses $`M_{LQ}`$ for scalar and vector leptoquarks. To compare the data with predictions of the BRW model the global probability function $`𝒫(\lambda _{LQ},M_{LQ})`$ is introduced, describing the probability that the data come from the model described by parameters $`\lambda _{LQ}`$ and $`M_{LQ}`$. The probability function is defined in such a way that the Standard Model probability $`𝒫_{SM}1`$. Constraints on the leptoquark couplings and masses were studied in the limit of very high leptoquark masses (using the contact interaction approximation ) as well as for finite leptoquark masses, with mass effects correctly taken into account. Excluded on 95% confidence level are all models (parameter values) which result in the global probability less than 5% of the Standard Model probability: $`𝒫(\lambda _{LQ},M_{LQ})<0.05`$. For models which describe the data much better than the Standard Model ($`𝒫_{max}\underset{\lambda ,M}{\mathrm{max}}𝒫(\lambda ,M)1`$) the 95% CL signal limit is defined by the condition: $`𝒫(\lambda _{LQ},M_{LQ})>0.05𝒫_{max}`$. Four leptoquark models are found to describe the existing experimental data much better than the Standard Model ($`𝒫(\lambda _{LQ},M_{LQ})>20`$). The signal limits for these models, at 68% and 95% CL are compared with exclusion limits in the $`(\lambda _{LQ},M_{LQ})`$ space in Figure 1. For $`S_1`$ and $`\stackrel{~}{V}_{}`$ leptoquarks the observed increase in the global probability by factor 367 and 142 respectively corresponds to more than a 3$`\sigma `$ effect. The leptoquark “signal” is mostly resulting from the new data on the atomic parity violation (APV) in cesium. After the theoretical uncertainties have been significantly reduced, the measured value of the cesium weak charge is now 2.5$`\sigma `$ away from the Standard Model prediction. Also the new hadronic cross-section measurements at LEP2, for $`\sqrt{s}`$=192–202 GeV, are on average about 2.5% above the predictions. The effect is furthermore supported by the slight violation of the CKM matrix unitarity and HERA high-$`Q^2`$ results. ## 4 Future experiments In the presented analysis experiments at the following existing and future colliders are considered: * HERA Since 1998 HERA collides 920 GeV protons with 27.5 GeV electrons or positrons, resulting in $`\sqrt{s}`$ 318 GeV. After the accelerator upgrade in 2000/2001 HERA is expected to deliver about 200 $`pb^1`$ of data every year. In this analysis the integrated luminosity of 400 $`pb^1`$ for each beam ($`e^{}p`$ and $`e^+p`$) is assumed. * The Tevatron The Tevatron is expected to start collecting data again in 2001. After the accelerator upgrade the $`p\overline{p}`$ center-of-mass energy of $`\sqrt{s}`$= 2 TeV is expected and the luminosity of about 1-2 $`fb^1`$ per year. Up to 10 $`fb^1`$ of the data can be collected before LHC turns on. Results given in this paper were calculated for integrated luminosities of 1 and 10 $`fb^1`$. * LHC The Large Hadron Collider (LHC), currently under construction at CERN, will collide proton beams at the center-of-mass energy of $`\sqrt{s}`$= 14 TeV. The luminosity expected at the very beginning is about 10 $`fb^1`$ per year and should increase up to about 100 $`fb^1`$ per year in the next years. Results given in this paper were calculated for integrated luminosities of 10 and 100 $`fb^1`$. * TESLA TESLA is one of the existing proposals for the next-generation $`e^+e^{}`$ linear collider. It would collide 250 GeV electron and positron beams ($`\sqrt{s}`$= 0.5 TeV), delivering the integrated luminosity of up to 500 $`fb^1`$. After the accelerator upgrade, beam energies of up to 500 GeV ($`\sqrt{s}`$= 1 TeV) should be reachable. High energy $`e^+e^{}`$ collisions can be also used to study $`e\gamma `$ and $`\gamma \gamma `$ interaction, with the effective photon flux described by the Waizsäcker-Williams Approximation (WWA). However, high quality electron beams of TESLA could be also used to produce high energy and high intensity photon beams from the Compton backscattering of laser light . In that case scattered photons usually take most of the electron energy. The energy spectrum is much harder than for WWA and peaked at the maximum photon energy $`E_\gamma ^{max}0.83E_e`$. Taking into account possible production of high intensity photon beams, three different scenarios are considered for TESLA, for $`\sqrt{s_{ee}}`$= 0.5 and 1 TeV, and the integrated luminosity of 100 $`fb^1`$: + $`e^+e^{}`$ scattering; $`e\gamma `$ and $`\gamma \gamma `$ collisions are also considered using the WWA effective photon flux, + $`e\gamma `$ scattering, with photon beam produced by Compton backscattering; the maximum $`e\gamma `$ center-of-mass energy $`\sqrt{s_{e\gamma }}0.91\sqrt{s_{ee}}`$, + $`\gamma \gamma `$ scattering, with both photon beams produced by Compton backscattering; the maximum $`\gamma \gamma `$ center-of-mass energy $`\sqrt{s_{\gamma \gamma }}0.83\sqrt{s_{ee}}`$, * THERA If TESLA project is approved, it is also possible to consider scattering of 250 or 500 GeV electron beam from TESLA with 1 TeV proton beam from HERA, resulting in the center-of-mass energy of $`\sqrt{s}`$= 1 and 1.4 TeV respectively. Expected integrated luminosity is of the order of 100 $`pb^1`$. ## 5 Contact interaction limit In the limit $`M_{LQ}\sqrt{s}`$ the effect of leptoquark production or exchange is equivalent to a vector type $`eeqq`$ contact interaction . Limits on the effective contact interaction mass scale $`\mathrm{\Lambda }`$ (related to the leptoquark mass to the coupling ratio $`M_{LQ}/\lambda _{LQ}`$) can be extracted from precise measurements of different Standard Model processes. Expected limits from future experiments are calculated assuming that no deviations from the Standard Model predictions will be observed. The method used has been described in details in . Limits from the following future measurements are considered: * Drell-Yan electron pair production at the Tevatron and LHC; from comparison with event selection efficiency is assumed to be 50% at the Tevatron and 80% at LHC, * hadronic cross-section measurements at TESLA $`e^+e^{}`$, assuming that $`\sigma (e^+e^{}q\overline{q})`$ is measured with 1% precision, * high-$`Q^2`$ NC DIS cross-section for $`e^\pm p`$ scattering at HERA and $`e^{}p`$ scattering at THERA, * high-$`Q^2`$ NC DIS cross-section for $`e\gamma `$ scattering at TESLA, for $`e^+e^{}`$ ($`\gamma `$ from WWA) and $`e\gamma `$ ($`\gamma `$ from Compton backscattering) scenarios. Parton densities are described by the MRST distribution functions for the proton and GRV LO distribution functions for the photon. 95% CL exclusion limits on $`M_{LQ}/\lambda _{LQ}`$, expected from different future experiments, are presented in Tables 2 and 3. Current limits from global analysis are included for comparison. Strongest limits on the leptoquark mass to the coupling ratio, in the contact interaction approximation are expected from the Drell-Yan electron pair production at LHC and hadronic cross-section measurements at TESLA. Results presented for Drell-Yan electron pair production and high-$`Q^2`$ NC DIS measurements are the mean values from about 2000 MC experiments. Poisson fluctuations in the observed numbers of events can result in the statistical fluctuation of the actual limit of 10-20%. ## 6 Limits from leptoquark pair production Leptoquark pair production has been considered for $`p\overline{p}`$ collisions at the Tevatron, $`pp`$ collisions at LHC and $`\gamma \gamma `$ scattering at TESLA (in $`e^+e^{}`$ and $`\gamma \gamma `$ scenarios). The advantage of the leptoquark pair production is that the cross-section depends only on the the strong or electromagnetic coupling constant and does not depend on the leptoquark Yukawa coupling. The leptoquark mass limits derived from the search for the leptoquark pair production are therefor valid for arbitrary values of $`\lambda _{LQ}`$. Presented results are based on the cross-sections given in . For vector leptoquark production at the Tevatron and LHC the anomalous coupling values resulting in the minimal pair production cross-section are assumed. For vector leptoquark production at TESLA minimal couplings $`\kappa `$=1, $`\lambda `$=0 are used. For leptoquark pair production at LHC the Standard Model background estimate is taken from . For leptoquark pair production searches at the Tevatron and at TESLA the Standard Model background is assumed to be negligible. Event selection efficiency is 25% at the Tevatron and 30% at LHC (from comparison with ). For TESLA it is assumed that (due to much “cleaner” environment of $`e^+e^{}`$ collisions) high event selection efficiency (close to 100%) is possible. 95% CL exclusion limits on the leptoquark mass $`M_{LQ}`$, expected from the negative search results at different future experiments, are presented in Table 4. In all cases strongest limits on the leptoquark mass are expected from the pair production search at LHC. ## 7 Limits from single leptoquark production Production of single leptoquarks has been considered for $`e^\pm p`$ scattering at HERA and THERA and for $`e\gamma `$ scattering at TESLA (in $`e^+e^{}`$ and $`e\gamma `$ scenarios). In the narrow-width approximation, the cross-section for single $`F=2`$ leptoquark production in electron-proton scattering (via the electron-quark fusion) is given by: $`\sigma ^{epLQX}(M_{LQ},\lambda _{LQ})`$ $`=`$ $`(J+1){\displaystyle \frac{\pi \lambda _{LQ}^2}{4M_{LQ}^2}}x_{LQ}q(x_{LQ},M_{LQ}^2)`$ (1) where $`J`$ is the leptoquark spin, $`q(x,Q^2)`$ is the quark momentum distribution function in the proton and $`x_{LQ}=\frac{M_{LQ}^2}{s}`$. For the single leptoquark production in $`e\gamma `$ collisions two approaches are possible: * leptoquark is produced in the electron fusion with a quark inside the photon. The production cross-section is given by the formula (1), with $`q(x,Q^2)`$ describing the quark momentum distribution in the photon. As before, GRV LO parton densities for the photon are used. * photon directly participates in the process $`e\gamma LQq`$. The cross-section for this process is taken from . Both approaches give very similar numerical results for the single leptoquark production cross-section. This is due to the fact that the dominant contribution to the direct photon process comes from the diagram in which photon splits into the $`q\overline{q}`$ pair, which is also described in the resolved photon approach, mentioned above. Only the leptoquark signal in the electron-jet decay channel is considered in this study and the resolution of the mass reconstruction is assumed to be 5%. Expected signal from single leptoquark production, for given leptoquark mass $`M_{LQ}`$ and Yukawa coupling $`\lambda _{LQ}`$, is compared with the observed number of events from the Standard Model background ($`ep`$ or $`e\gamma `$ NC DIS). The background can be suppressed by applying a cut on the Bjorken variable $`y`$, which is optimized for every leptoquark type as a function of the leptoquark mass. The 95%CL exclusion limit on the leptoquark Yukawa coupling $`\lambda _{LQ}`$ corresponds to the decrease of the Poisson probability for the number of events observed in the leptoquark mass window to 5% of the Standard Model probability. Average $`\lambda _{LQ}`$ exclusion limits expected from future experiments are calculated based on about 2000 MC experiments generated according to the Standard Model expectations. In Figures 2, 3 and 4 combined 95% CL exclusion limits in $`(\lambda _{LQ},M_{LQ})`$ space, expected from single leptoquark production, leptoquark pair production and contact interaction analysis, are presented for future $`e^+e^{}`$ scattering at TESLA, $`e\gamma `$ and $`\gamma \gamma `$ scattering at TESLA and $`ep`$ scattering at THERA, respectively. Results obtained for TESLA electron beam energy of 250 and 500 GeV are compared with existing limits. In all cases search for single leptoquark production significantly improves the limits coming from leptoquark pair production and indirect searches. For coupling values $`\lambda _{LQ}`$ 0.1 the leptoquark mass limits can be extended up to the kinematic limit $`M_{LQ})\sqrt{s}`$. For lower leptoquark masses (but above the limit from the leptoquark pair production), limits on the leptoquark Yukawa coupling $`\lambda _{LQ}`$ can be improved by an order of magnitude. ## 8 Summary of results In Figures 5, 6, 7 and 8, 95% CL exclusion limits in $`(\lambda _{LQ},M_{LQ})`$ space, expected from different experiments are compared with existing limits and possible leptoquark signal, for $`S_{1/2}^R`$, $`S_1`$, $`\stackrel{~}{V}_{}`$ and $`V_{1/2}^L`$ leptoquark models respectively. New HERA data are not expected to improve the existing leptoquark limits significantly. Also the effect of future measurements at the Tevatron is moderate: leptoquark mass limits will increase by 100-150 GeV (see also Table 4). Sizable improvement of the leptoquark mass and coupling limits is expected from high luminosity LHC data. Exclusion limits expected from the leptoquark pair production at LHC will extend up to the leptoquark masses of about 2.2 TeV. The possible $`S_1`$ and $`\stackrel{~}{V}_{}`$ leptoquark “signal”, resulting from the new data on the atomic parity violation, can be confirmed or excluded even for very high leptoquark masses. If the leptoquark type particle is discovered at LHC, with mass below 1 TeV, TESLA and THERA will be an ideal place to study its properties, provided that the Yukawa coupling is not too small. For $`e\gamma `$ scattering at TESLA leptoquarks with Yukawa couplings down to $`\lambda _{LQ}`$ 0.01 can be studied. ## Acknowledgements This work has been partially supported by the Polish State Committee for Scientific Research (grant No. 2 P03B 035 17).
warning/0006/cond-mat0006199.html	ar5iv	text	# Can Polymer Coils be modeled as “Soft Colloids”? ## Abstract We map dilute or semi-dilute solutions of non-intersecting polymer chains onto a fluid of “soft” particles interacting via a concentration dependent effective pair potential, by inverting the pair distribution function of the centers of mass of the initial polymer chains. A similar inversion is used to derive an effective wall-polymer potential; these potentials are combined to successfully reproduce the calculated exact depletion interaction induced by non-intersecting polymers between two walls. The mapping opens up the possibility of large-scale simulations of polymer solutions in complex geometries. A statistical description of polymer solutions in complex geometries, such as the colloid-polymer mixtures which have recently received much experimental attention, generally relies on a nanometer scale segment representation of the polymer coils, a computationally very demanding task except in the special case of ideal (non-intersecting) polymers obeying Gaussian statistics. This obviously follows from the fact that, although the colloidal particles may reasonably be modeled by hard impenetrable spheres or other complex shapes lacking internal structure, each polymer coil involves $`L`$ segments which must satisfy a non-intersection constraint. It thus appears natural to attempt a mesoscale coarse-graining, whereby polymer coils interact via effective pair potentials acting between their centers of mass (CM). Since polymers can interpenetrate, the effective potential $`\beta v(r)`$ is expected to be soft, with a range of the order of the radius of gyration $`R_g`$ of individual coils. Such a coarse-grained description has been a long-time goal in the statistical mechanics of polymer solutions, dating back to the first attempts by Flory and Krigbaum who employed mean-field theory to find an interaction for which the strength at overlap scales as: $`\beta v(r`$$`=`$$`0)L^{0.2}`$. Later, scaling arguments, field-theoretical renormalization group calculations, and simulations confirmed that the range of the interaction between two isolated polymer coils is of order $`R_g`$, but found that in the scaling limit the strength $`\beta v(r`$$`=`$$`0)`$ is independent of $`L`$ and of order $`k_BT`$. In this letter, we show that a meaningful “soft colloid” picture of polymer coils may be built on a coherent “first principles” statistical mechanical foundation. We derive both the effective wall-polymer CM interaction $`\beta \varphi (z)`$, and the “best” local effective pair-potential $`\beta v(r)`$ between polymer CM’s for finite polymer concentrations. These potentials are then applied to simulate bulk polymer solutions, as well inhomogeneous polymers near a hard wall and polymers confined between two parallel walls to extract the effective depletion potential between plates. The “soft colloid” approach turns out to be successful not only in the dilute regime but also, perhaps more surprisingly, well into the semi-dilute regime. A related “soft particle” picture has been applied to polymer melts and blends, but the corresponding phenomenological implementation differs substantially from the present “first principles” approach. We consider a popular model for polymers in a good solvent, namely $`N`$ excluded volume polymer chains of $`L`$ segments undergoing non-intersecting self avoiding walks (SAW) on a simple cubic lattice of $`M`$ sites, with periodic boundary conditions. The packing fraction is equal to the fraction of lattice sites occupied by polymer segments, $`c`$$`=`$$`N\times L/M`$, while the concentration of polymer chains is $`\rho `$$`=`$$`c/L`$$`=`$$`N/M`$. For a single SAW chain, the radius of gyration $`R_gL^\nu `$, where $`\nu 0.6`$ is the Flory exponent. The overlap concentration $`\rho ^{}`$, signaling the onset of the semi-dilute regime, is such that $`4\pi \rho ^{}R_g^3/31`$, and hence $`\rho ^{}L^{3\nu }`$. We have carried out MC simulations for chains of length $`L`$$`=`$$`100`$ and $`L`$$`=`$$`500`$, and covered a range of concentrations up to $`\rho /\rho ^{}5`$. The pair distribution function $`g(r)`$ of the centers of mass was computed for several concentrations; $`g(r`$$`=`$$`0)`$ is always non-zero, thus confirming the “softness” of the effective pair potential $`\beta v(r)`$. The latter was then derived from $`g(r)`$ by an inversion procedure based on the hyperneted-chain approximation (HNC) closure relation: $$g(r)=\mathrm{exp}\{\beta v(r)+g(r)c(r)1\},$$ (1) where $`\beta `$$`=`$$`1/k_BT`$, while $`c(r)`$ is the direct pair correlation function, related to $`g(r)`$ by the Ornstein-Zernike (OZ) relation. To any given $`g(r)`$ and density there corresponds a unique effective pair potential $`\beta v(r)`$, capable of reproducing the input $`g(r)`$, irrespective of the underlying many-body interactions in the system; in a variational sense this $`\beta v(r)`$ provides the “best” pair representation of the true interactions, and leads back to the true thermodynamics via the compressibility relation. While the simple HNC inversion procedure would be inadequate for dense fluids of hard core particles, where more sophisticated closures or iterative procedures are required, we are able to demonstrate the consistency of the HNC inversion in the present case. If the resulting effective $`\beta v(r)`$, examples of which are shown in Fig. 1, are used directly in MC simulations, the calculated “exact” $`g(r)`$ for this effective representation coincides within statistical errors with the $`g(r)`$ derived from the simulation of the full initial polymer segment model. In fact, the HNC closure turns out to be quasi-exact when applied to the simple Gaussian model whereby particles interact via the potential $`\beta v(r)`$$`=`$$`ϵ\mathrm{exp}[\alpha \left(r/R_g\right)^2]`$, which yields a reasonable fit to the effective pair potentials shown in Fig. 1. Even the much cruder random phase approximation closure, $`c(r)`$$`=`$$`\beta v(r)`$, yields semi-quantitatively accurate results in the regime of interest. Careful inspection of Fig. 1 reveals that the effective pair potential is not very sensitive to the polymer concentration. The value at $`r`$$`=`$$`0`$ first increases slightly with $`\rho `$, before decreasing again at the highest concentration. More strikingly, and perhaps not surprisingly, the range of $`\beta v(r)`$ increases with $`\rho `$. The effective potential becomes slightly negative $`(𝒪(10^3k_BT))`$ for $`r/R_g3`$ at the higher concentrations. The properties of soft-core fluids are significantly different from their hard-core counterparts. For example, for potentials of the type shown in Fig. (1), the pressure is very well described by $`\beta P`$$`=`$$`\rho +1/2\beta \widehat{V}(0)\rho ^2`$ over the entire density range. Here $`\widehat{V}(0)`$ is the Fourier transform of the potential, at $`k`$$`=`$$`0`$. Our observation that potentials become slightly longer ranged at higher densities implies that the pressure scales with an exponent slightly higher than $`2`$, so that the equation of state (e.o.s.) is at consistent with the well-known $`\rho ^{2.25}`$ law. At first sight it may seem surprising that a two-body potential could reproduce the full e.o.s. without explicit many body terms. However, the effective potential we use is constructed to reproduce the true thermodynamics through the compressibility relation (ignoring small volume terms); the relative insensitivity of $`\beta v(r)`$ to concentration implies that many-body interactions are not very important. This insensitivity to concentration makes it possible to apply the effective potential appropriate for a given mean concentration to inhomogeneous cases, where the local polymer concentration deviates from the mean. Such a situation occurs when a polymer solution is confined by a hard wall. Using the same explicit SAW polymer model in MC simulations, we have computed the “exact” profiles $`h(z)`$$`=`$$`\rho (z)/\rho 1`$, where $`z`$ denotes the perpendicular distance of the polymer CM from the wall. Examples of $`h(z)`$ for several bulk concentrations are shown in the inset of Fig. 2. The corresponding adsorptions $`\mathrm{\Gamma }`$ are defined by: $$\mathrm{\Gamma }=\frac{(\mathrm{\Omega }^{ex}/A)}{\mu }=\rho _0^{\mathrm{}}h(z)𝑑z,$$ (2) where $`\mathrm{\Omega }^{ex}/A`$ is the excess grand potential per unit area, $`\rho `$ the bulk concentration of the polymers and $`\mu `$ their chemical potential. From a knowledge of the concentration profile $`\rho (z)`$, and the bulk direct correlation function between polymers CM’s $`c(r)`$, one may extract an effective wall-polymer potential $`\beta \varphi (z)`$ by combining the wall-polymer OZ relations with the HNC closure, resulting in: $$\beta \varphi (z)=\beta \varphi ^{MF}(z)+\rho 𝑑𝐫^{}h(z^{})c(\|𝐫𝐫^{}\|).$$ (3) The first term is the usual potential of mean force $`\beta \varphi ^{MF}(z)`$$`=`$$`\mathrm{ln}\left[\rho (z)/\rho \right]`$, to which $`\beta \varphi (z)`$ would reduce in the $`\rho 0`$ limit, while the second term arises from correlations between the polymer coils next to the wall. Using the $`c(r)`$ extracted from the earlier bulk simulations of $`g(r)`$, together with Eq. (3), we are able to extract $`\beta \varphi (z)`$ from the density profiles. Results for various bulk concentrations are plotted in the Fig. (2). The range of the effective wall-polymer repulsion increases with increasing concentration, while the density profiles actually move in closer to the wall. It is important to stress that the correlation term considerably enhances the repulsion compared to the potential of mean force. We have tested the consistency of the inversion procedure (which, to the best of our knowledge, has not been attempted before for any wall/fluid interface) by using $`\beta \varphi (z)`$, and the pair potential $`\beta v(r)`$ for the appropriate bulk concentration, in MC simulations based on these effective interactions (such simulations are at least an order of magnitude faster than simulations of the initial segment model). The resulting concentration profile of the effective “soft colloids” agrees to within statistical accuracy with the initial $`\rho (z)`$ obtained from the detailed segment simulations, and the corresponding adsorption $`\mathrm{\Gamma }`$ differs by less than $`1\%`$ from the exact value, thus demonstrating the adequacy of the “soft colloid” representation of the interacting polymer coils. An even more severe test of this representation is provided by a calculation of the depletion interaction between two hard walls confining the polymers within a slit of width $`d`$. Using direct grand-canonical simulations of the full SAW polymer model, we computed the osmotic pressure exerted by the polymer-coils on the walls; the interaction free energy per unit area A, $`\beta \mathrm{\Delta }F/A`$, is then obtained by integrating the osmotic pressure calculated for different values of the spacing $`d`$ between the walls. These simulations are extremely computer intensive, and were only be carried out for $`L`$$`=`$$`100`$. In the “soft colloid” picture, the interactions of the polymer CM’s with each other, $`\beta v(r)`$, and with a wall, $`\beta \varphi (z)`$, are calculated once with the HNC inversion procedures from the $`g(r)`$ and $`\rho (z)`$ of a full SAW polymer simulation at the bulk density. These are then used in grand-canonical MC simulations of soft particles between two walls, and in Fig. 3 they are compared to the ‘exact” grand-canonical MC simulations of $`L`$$`=`$$`100`$ SAW polymers (for $`\rho /\rho ^{}=0.95`$) . The results are in good agreement, but the “soft colloid” calculations are at least two orders of magnitude faster. Contrary to the more widely studied case of colloid-colloid mixtures, the “exact” interaction exhibits no significant repulsive barrier, whilst the “soft colloid” model leads to a flat maximum; the corresponding barrier height is, however, very small compared to the attractive minimum at contact, which agrees well with the “exact” data, as does the slope of the attraction. In fact, the repulsive barrier does not increase significantly with density, and its origin can be traced to our use of the “potential overlap approximation”, namely that the interaction of the soft particles with two parallel walls a distance $`d`$ apart can be written as the sum of the two individual wall-particle interactions. This is exact for simple liquids with true intermolecular interactions, but not for polymers described by effective potentials, even if the polymers are ideal. For the sake of consistency, the MC simulations for the “soft colloid” model were carried out with effective wall-polymer and polymer-polymer potentials appropriate for $`L`$$`=`$$`100`$. However, we checked that the data obtained with effective interactions appropriate for longer polymers ($`L`$$`=`$$`500`$), which cannot be easily handled within the full segment model, are very close to the $`L`$$`=`$$`100`$ results, so we are confident that we are close to the scaling regime. In Fig. 3 we also compare two results derived in the spirit of the Asakura-Oosawa (AO) approximation. The free-energy difference $`\beta \mathrm{\Delta }F(z)`$ is modeled by the density times the exact depletion volume, $`\mathrm{\Delta }V_{id}(z)`$, excluding one ideal Gaussian polymer of size $`R_g`$, or by a popular phenomenological improvement: $`\beta \mathrm{\Delta }F(z)`$$`=`$$`\beta \mathrm{\Pi }_b\mathrm{\Delta }V_{id}(z)`$, where $`\mathrm{\Pi }_b`$ is the bulk osmotic pressure of the interacting polymers. Note that for the density under consideration, these approximations are seen to be very poor, both as regards the depth and the range of the depletion attraction. In fact, the range of the depletion interaction for interacting polymer coils is significantly reduced compared to the AO predictions, valid for ideal polymers. For low densities we find, as expected, that all the above approaches converge. These observations can be understood within the “soft colloid” representation and the HNC approximation, where the interacting free energy per unit area is given by: $`{\displaystyle \frac{\beta \mathrm{\Delta }F(z)}{A}}`$ $`=`$ $`\rho {\displaystyle _{\mathrm{}}^{\mathrm{}}}h(s)h(zs)𝑑s`$ (4) $`+`$ $`\rho {\displaystyle _{\mathrm{}}^{\mathrm{}}}h(zs)\left[\beta \varphi (s)\beta \varphi ^{MF}(s)\right]𝑑s.`$ (5) Here $`h(z)`$$`=`$$`\rho (z)/\rho 1`$ is the single wall density profile, $`\beta \varphi (z)`$ is the corresponding effective wall-polymer potential, and $`\beta \varphi ^{MF}(z)`$ is the corresponding potential of mean force. The first term on the r.h.s. of Eq. (4) is the density overlap approximation and would be the only-contribution in the case of ideal (Gaussian) polymer coils. The second term arises from the correlations between coils; this dominates the first term in the semi-dilute regime $`(\rho /\rho ^{}1)`$. The standard AO approximation may be derived from Eq. (4) by replacing the density profile by a step-function of width $`R_g`$ in the first term of Eq. (4) and neglecting the correlation term. In Fig. 3 we compare the HNC approach of Eq. (4) for the wall-wall interaction to the “exact” results and the MC simulations of the “soft colloids”. As was found for the homogeneous case and for the single wall, HNC works very well here, demonstrating that knowledge of $`\beta v(r)`$ and $`\beta \varphi (z)`$ quickly leads to accurate predictions for the slit geometry, paving the way for the use of integral equation techniques in other, more complex, geometries. To summarize, the coarse-grained representation of polymer coils as “soft” colloids has been shown to be very reliable, yielding pair distribution functions and concentration profiles which agree closely with the results for the full SAW segment model, while being much more efficient from a computational point of view. Much of the success of the coarse-graining lies in our finding that the “best” effective pair potential between CM’s of neighboring coils does not depend strongly on polymer concentration, and is reasonably close to its $`\rho 0`$ limit. Similar conclusions were reached in recent work on the phase-behavior of star-polymers, where the $`\rho 0`$ limit of the pair potential was used to calculate the phase-behavior at finite concentration. Our results for the linear polymer case suggest that the full pair-potential for star polymers may not be strongly concentration dependent, and that our approach could be used for star polymers in confined geometries. Finally, we note that the “soft colloid” description is expected to work best in complex geometries where the curvature is not too large on the scale of $`R_g`$, such as colloid-polymer mixtures where the colloid radius $`RR_g`$. For such systems, the “soft colloid” model may now be used in large scale simulations or fluid integral equations of polymers in complex geometries, such as the structure, phase behavior, interactions, and metastability of colloid polymer mixtures, which cannot be achieved with the detailed model of non-intersecting polymer chains. AAL acknowledges support from the Isaac Newton Trust, Cambridge, PB acknowledges support from the EPSRC under grant number GR$`/`$M88839, EJM acknowledges support from the Royal Netherlands Academy of Arts and Sciences. We thank David Chandler, Daan Frenkel, Christos Likos, Hartmut Löwen, and Patrick Warren for helpful discussions.
warning/0006/quant-ph0006043.html	ar5iv	text	# Hidden-variable theorems for real experiments ## Abstract It has recently been questioned whether the Kochen-Specker theorem is relevant to real experiments, which by necessity only have finite precision. We give an affirmative answer to this question by showing how to derive hidden-variable theorems that apply to real experiments, so that non-contextual hidden variables can indeed be experimentally disproved. The essential point is that for the derivation of hidden-variable theorems one does not have to know which observables are really measured by the apparatus. Predictions can be derived for observables that are defined in an entirely operational way. In general quantum mechanics only makes probabilistic predictions for individual events. The question whether one can go beyond quantum mechanics in this respect has been a subject of debate and research since the early days of the theory. There are famous theorems placing restrictions on possible hidden-variable theories reproducing the results of quantum mechanics. Bell’s theorem excludes local hidden-variables. The Kochen-Specker theorem excludes non-contextual hidden variables. In local hidden-variable theories the pre-determined results for a given measurement are independent of which measurements are performed at space-like separation. In non-contextual hidden-variable theories the pre-determined results are independent of any measurements that are performed jointly. The present work is concerned with the derivation of hidden-variable theorems that apply to real experiments. It was motivated by the thought–provoking work of Meyer , who claimed that the Kochen-Specker theorem was “nullified” for real experiments because of the unavoidably finite measurement precision, i.e. because the observer does not have full control over his experimental setup. However, our approach is very general and applies in a much broader context. It allows one to deal with all possible experimental imperfections, and it is not restricted to the specific class of non-contextual hidden variables. In particular, theorems on local hidden variables can be derived in the same way. Let us begin our discussion with some remarks on experimental tests of hidden variables in general. Hidden-variable theorems provide us with predictions made by certain classes of hidden-variable theories which can be tested experimentally and which differ from the quantum mechanical predictions. For local hidden variables, Bell’s inequalities can be derived without making any reference to quantum mechanics. Then one can check whether they are fulfilled by the quantum mechanical predictions, and by experiment. For non-contextual hidden variables, one can use the Kochen-Specker argument to derive specific quantitative predictions for the results of experiments, which can then be compared to quantum mechanics and to nature. For experimental proposals and Kochen-Specker type experiments that were performed in this spirit see . It is sometimes argued that the Kochen-Specker theorem just makes a conceptual point about the formal structure of quantum mechanics and that actually performing experimental tests for non-contextual hidden variables is unnecessary, since quantum mechanics is a very well tested theory. Irrespective of this discussion, we feel that the question whether such tests could be performed in principle is interesting because it concerns the falsifiability of a rather simple and fundamental concept. In the derivation of hidden-variable theorems frequently no direct reference to experiments is made, or the experimental situation is treated in an idealized way. For example, in the original derivation of the Kochen-Specker theorem it was just shown that non-contextual hidden variables are incompatible with quantum mechanics, without making the experimental predictions of non-contextual hidden variables explicit, while in the original derivation of Bell’s theorem perfect correlations and perfect detection efficiency were assumed. If one wants to consider an experimental test for a certain class of hidden-variable theories, one has to give up these idealized assumptions and derive hidden-variable theorems that apply to the true experimental situation. For local hidden variables, the case of non-perfect correlations and detection efficiency was analyzed soon after Bell’s original derivation, starting with the work of Clauser and Horne . The same kind of analysis is possible for non-contextual hidden variables. However, there is one more idealization made in the usual derivations of hidden-variable theorems which, to our knowledge, has never been discussed in an explicit way, namely the accuracy with which the experimental setup can be maintained (e.g. how well the Stern-Gerlach magnets in a spin measurement can be aligned and kept stable). This question has become particularly important for the Kochen-Specker theorem in view of recent claims by Meyer that this theorem is no longer applicable when the measurements have only finite precision. In its original form, the Kochen-Specker theorem states that is is impossible to assign values to observables corresponding to all directions on the sphere subject to a constraint for triads of orthogonal directions. The observables under consideration are the squares of the spin components of a spin-1 particle along the respective direction, and the constraint is given by the total spin. Meyer’s claim was based on the fact that it is possible to assign values compatible with the constraint to all rational directions on the sphere, which constitute a dense subset of all directions . He argued that, since measurements with finite precision cannot discriminate a dense subset from its closure, this implies that non-contextual hidden variables cannot be excluded by any real Kochen-Specker type experiment. However, he did not construct an explicit non-contextual hidden-variable model for real experiments with finite precision . At first sight, it does not seem possible to refute Meyer’s claim on the basis of existing theorems, for the following reason. An essential feature of all hidden-variable theorems known to us is that some observables have to appear in different experimental contexts: an observable $`A`$ has to be measured together with another observable $`B`$ on a sample of systems, and also together with a third observable $`C`$ on another sample of systems from the same source, where both $`B`$ and $`C`$ commute with $`A`$, but they do not commute with each other and thus cannot be measured jointly. Note that even for tests of non-contextuality it is not necessary to perform several incompatible measurements on the same system, which would of course be impossible . For example, in the original Kochen-Specker situation, one can only arrive at a contradiction by considering several triads of directions that have at least some directions in common. For Kochen-Specker experiments this implies that, at least for some directions, the observables corresponding to these directions have to appear in different triads. When the finite precision of real experiments is taken into account, it seems impossible to ascertain that the same observable is really measured more than once in different experimental contexts. Thus the usual derivations of hidden-variable theorems seem to run into problems. Does this mean that it is impossible in principle to rule out non-contextuality experimentally? Could even the experimental results on local hidden variables be in danger for the same reasons? We are going to show that the answer is no. There are predictions of non-contextual hidden variables, which can be derived within a framework that is sufficiently general to apply to real experiments, and which can therefore also be tested by real experiments with sufficiently high, but finite, precision. In order to achieve this, it is important to realize that predictions for classes of hidden-variable theories can be derived without reference to quantum mechanics. In such an approach, the observables playing a role in the hidden-variable predictions are not defined via the quantum-mechanical formalism, but in an operational way. For concreteness, imagine that an observer wants to perform a measurement of the spin square of a spin-1 particle along a certain direction $`\stackrel{}{n}`$. There will be an experimental procedure for trying to do this as accurately as possible. We will refer to this procedure by saying that he sets the ”control switch” of his apparatus to the position $`\stackrel{}{n}`$. In all experiments that we will discuss only a finite number of different switch positions is required. By definition different switch positions are clearly distinguishable for the observer, and the switch position is all he knows about. Therefore, in an operational sense the measured physical observable is entirely defined by the switch position. From the above definition it is clear that the same switch position can be chosen again and again in the course of an experiment (while of course the system measured will always be different, cf. ). In general one has to allow for the possibility that the switch position $`\stackrel{}{n}`$ does not uniquely determine the physical state of the measuring apparatus, i.e. there may be degrees of freedom of the apparatus over which the observer does not have full control but which may influence the result of any given measurement. In the context of deriving hidden-variable theorems, this possibility can be accomodated in a very simple and general way. Following the philosophy of deterministic hidden variable theories , there must also be some (in general hidden, i.e. unknown) variables determining the behaviour of the apparatus, and one has to assume that the result of any measurement will be determined by the hidden variables of the system and by those of the apparatus together. Notice that in such an approach as described in the two preceding paragraphs it does not matter which observable is “really” measured by the apparatus and to what precision. One just derives general predictions for the behaviour of system and apparatus together, provided that certain switch positions are chosen. These predictions only depend on the properties of the class of hidden-variable theories considered. The question of the correct quantum-mechanical description of the non-ideal measurement considered arises only when, as a next step, one wants to obtain the quantum-mechanical predictions for the given situation. Following the method described above, we will now show how non-contextual hidden variables can be tested and thus potentially excluded by real experiments. We will consider the context of the original Kochen-Specker argument, i.e. exactly the case considered by Meyer. In the original Kochen-Specker situation one considers a spin-1 particle. In the ideal case of perfect precision, the relevant quantum-mechanical observables are the squares of the spin components, denoted by $`\widehat{S}_\stackrel{}{n}^2`$ for arbitrary directions $`\stackrel{}{n}`$. For a spin-1 particle one has $$\widehat{S}_{\stackrel{}{n}_1}^2+\widehat{S}_{\stackrel{}{n}_2}^2+\widehat{S}_{\stackrel{}{n}_3}^2=2$$ (1) for every triad of orthogonal directions $`\{\stackrel{}{n}_1,\stackrel{}{n}_2,\stackrel{}{n}_3\}`$. As the possible results for every $`\widehat{S}_{\stackrel{}{n}_i}^2`$ are 0 or 1, this implies that in the ideal case for every measurement of three orthogonal spin squares two of the results will be equal to one, and one of them will be equal to zero. Let us emphasize that in our approach, in the derivation of the hidden-variable predictions, the observables are defined operationally by the switch positions, i.e. by the best effort and knowledge of the experimenter, and cannot be identified with the quantum-mechanical observables. Of course, for a specific experiment, there should be some approximate correspondence in order to ensure that the quantum-mechanical predictions will be sufficiently close to the ideal case so that they are still in conflict with the relevant hidden-variable predictions. In the following the symbol $`S_\stackrel{}{n}^2`$ (without the hat) will denote the operational observable defined by the switch position $`\stackrel{}{n}`$, and the term direction will be used as a synonym for switch position. In a deterministic hidden variable theory (cf. ) one assumes that for every individual particle the result of the measurement of any observable $`S_\stackrel{}{n}^2`$ is predetermined by hidden variables. In non-contextual hidden variable theories it is furthermore assumed that this predetermined result does not depend on the ”context” of the measurement, i.e. on which other observables are measured together with $`S_\stackrel{}{n}^2`$, but only on the switch position $`\stackrel{}{n}`$ and the hidden variables . In general the result may depend both on the hidden variables of the system and of the apparatus. Let us denote the hidden variables of the system by $`\lambda `$ and those of the apparatus by $`\mu `$. As explained above, the philosophy of non-contextual hidden variables implies the existence of a function $`S_\stackrel{}{n}^2(\lambda ,\mu )`$ taking values 0 and 1 which describes the result of a measurement with switch position $`\stackrel{}{n}`$ on a system characterized by $`\lambda `$ with an apparatus characterized by $`\mu `$. For fixed $`\lambda `$ and $`\mu `$ this function therefore assigns a value 0 or 1 to the switch position $`\stackrel{}{n}`$ . A Kochen-Specker experiment can now be performed by testing the validity of Eq. (1) (without hats) for a judiciously chosen set of triads of directions. Therefore the apparatus is required to have three switches where the three directions of a given triad can be chosen. Because the switch positions do not correspond to the ideal quantum mechanical observables the sum of the three results will not always be equal to 2. Nevertheless a contradiction between non-contextuality and quantum mechanics can be obtained in the following way. From the Kochen-Specker theorem it follows that there are finite sets of triads for which no value assignment consistent with Eq. (1) (again without the hats) is possible . Let us choose such a Kochen-Specker set of triads $$\{\{\stackrel{}{n}_1,\stackrel{}{n}_2,\stackrel{}{n}_3\},\{\stackrel{}{n}_1,\stackrel{}{n}_4,\stackrel{}{n}_5\},\mathrm{},\}.$$ (2) Let us emphasize that at least some of the switch positions $`\stackrel{}{n}_i`$ have to appear in several of the triads; clearly otherwise there could be no inconsistency. Let us denote the number of triads in the Kochen-Specker set (2) by $`N`$. The set is constructed in such a way that for any fixed values of $`\lambda `$ and $`\mu `$ the equation $$S_{\stackrel{}{n}_i}^2(\lambda ,\mu )+S_{\stackrel{}{n}_j}^2(\lambda ,\mu )+S_{\stackrel{}{n}_k}^2(\lambda ,\mu )=2$$ (3) must be violated for at least one out of the $`N`$ triads $`\{\stackrel{}{n}_i,\stackrel{}{n}_j,\stackrel{}{n}_k\}`$. Suppose that one can establish experimentally that for all triads in the Kochen-Specker set the sum of the results is equal to 2 in a fraction greater than $`1ϵ`$ of all cases. For the hidden variables this implies that Eq. (3) must hold for a fraction $`1ϵ`$ of all $`(\lambda ,\mu )`$, for all triads in the set. But for sufficiently small $`ϵ`$ this implies that there would have to be pairs $`(\lambda ,\mu )`$, for which Eq. (3) holds for all triads in the set, which is impossible according to the Kochen-Specker argument. To determine the required value for $`ϵ`$, it is convenient to use a set-theoretic language. Let us denote the set of all pairs of hidden variables $`(\lambda ,\mu )`$ by $`\mathrm{\Lambda }`$. Furthermore let us denote the subset of hidden variables for which the sum of spin squares is equal to $`2`$ for the $`k`$-th triad by $`\mathrm{\Lambda }_k`$. The value of $`ϵ`$ has to be sufficiently small such that the intersection of all the $`\mathrm{\Lambda }_k`$ cannot be empty. If we define the measure (i.e. the size in terms of probability) of $`\mathrm{\Lambda }`$ to be $`1`$, then according to our assumptions all the $`\mathrm{\Lambda }_k`$ have measures larger than $`1ϵ`$, which implies that the measure of the intersection of all the $`\mathrm{\Lambda }_k`$ is larger than $`1Nϵ`$. This follows from the fact that the complement of each $`\mathrm{\Lambda }_k`$ is smaller than $`ϵ`$, such that the size of the union of these complements, which is the complement of the intersection of all the $`\mathrm{\Lambda }_k`$, cannot be larger than $`Nϵ`$. Thus non-contextuality is experimentally disproved as soon as $`ϵ`$ is smaller than $`1/N`$, because then there would have to be hidden variables which lead to a sum of spin squares equal to $`2`$ for all the triads, which is impossible because of the structure of the Kochen-Specker set. Note that $`ϵ`$ describes all the imperfections of a real experiment including finite precision but also e.g. imperfect state preparation and non-unit detection efficiency. The value of $`N`$ and therefore the bound on $`ϵ`$ depends on the particular Kochen-Specker set used . As noted above, an inevitable requirement for the contradiction to be obtained is the fact that the function $`S_{\stackrel{}{n_1}}^2(\lambda ,\mu )`$, or in general functions corresponding to at least some switch positions, appear in more than one out of the $`N`$ triads. This appearance of the same function in different lines of the mathematical proof, corresponding to different experimental contexts, is possible in spite of finite experimental precision only because we defined our observables operationally via the switch positions. We have shown how non-contextual hidden-variable theories can be disproved by real experiments. This shows that Meyer’s coloring of a dense subset of the sphere cannot be used to construct non-contextual hidden-variable theories according to the above definition. The values assigned to specific switch positions would have to depend on the context, i.e. on the other switch positions chosen simultaneously by the experimenter. This form of contextuality is a feature of existing explicit models based on Meyer’s idea . In view of our results, we would assert that the Kochen-Specker theorem is not ”nullified” by finite measurement precision. Let us note that arguments in favor of this conclusion were given in . Our suggestion how to perform a Kochen-Specker experiment was inspired by some of Mermin’s remarks in . Let us emphasize that using the method of the present paper one can also show that local hidden variables can be disproved in real experiments, e.g. using the GHZ form of Bell’s theorem which is also based on sets of propositions that cannot be consistently satisfied by a class of hidden-variable theories (local, in this case). Inequalities corresponding to our bound on $`ϵ`$ can be derived and compared to the experimental data . Thus our work confirms the fact that fundamental concepts about the world can indeed be put to experimental test. When this work was completed, we learned from J.-A. Larsson that he has come to similar conclusions using a somewhat related approach . C. S. would like to thank L. Hardy for a useful discussion. This work has been supported by the Austrian Science Foundation (FWF), projects S6504 and F1506, and by the QIPC Program of the European Union.
warning/0006/hep-lat0006005.html	ar5iv	text	# References Three-Quark Potential in SU(3) Lattice QCD T.T. Takahashi<sup>1</sup>, H. Matsufuru<sup>1</sup>, Y. Nemoto<sup>2</sup> and H. Suganuma<sup>3,1</sup> <sup>1</sup> RCNP, Osaka University, Mihogaoka 10-1, Osaka 567-0047, Japan <sup>2</sup> YITP, Kyoto University, Kitashirakawa, Sakyo, Kyoto 606-8502, Japan <sup>3</sup> Faculty of Science, Tokyo Institute of Technology, Tokyo 152-8551, Japan ## Abstract The static three-quark (3Q) potential is studied in SU(3) lattice QCD with $`12^3\times 24`$ and $`\beta =5.7`$ at the quenched level. From the 3Q Wilson loop, 3Q ground-state potential $`V_{3\mathrm{Q}}`$ is extracted using the smearing technique for ground-state enhancement. With accuracy better than a few %, $`V_{3\mathrm{Q}}`$ is well described by a sum of a constant, the two-body Coulomb term and the three-body linear confinement term $`\sigma _{3\mathrm{Q}}L_{\mathrm{min}}`$, with $`L_{\mathrm{min}}`$ the minimal value of total length of color flux tubes linking the three quarks. Comparing with the Q-$`\overline{\mathrm{Q}}`$ potential, we find a universal feature of the string tension, $`\sigma _{3\mathrm{Q}}\sigma _{\mathrm{Q}\overline{\mathrm{Q}}}`$, and the OGE result for Coulomb coefficients, $`A_{3\mathrm{Q}}\frac{1}{2}A_{\mathrm{Q}\overline{\mathrm{Q}}}`$. In usual, the three-body force is regarded as a residual interaction in most fields of physics. In QCD, however, the three-body force among three quarks is expected to be a “primary” force reflecting the SU(3)<sub>c</sub> gauge symmetry. Indeed, the three quark (3Q) potential is directly responsible to the structure and properties of baryons , similar to the relevant role of the Q-$`\overline{\mathrm{Q}}`$ potential upon meson properties . In contrast with a number of studies on the Q-$`\overline{\mathrm{Q}}`$ potential using lattice QCD , there were only a few lattice QCD studies for the 3Q potential done mainly more than 13 years ago . (In Ref., the author only showed a preliminary result on the equilateral-triangle case without enough analyses.) In Refs., the 3Q potential seemed to be expressed by a sum of two-body potentials, which supports the $`\mathrm{\Delta }`$-type flux tube picture . On the other hand, Ref. seemed to support the Y-type flux-tube picture rather than the $`\mathrm{\Delta }`$-type one. These controversial results may be due to the difficulty of the accurate measurement of the 3Q ground-state potential in lattice QCD. For instance, in Refs., the authors did not use the smearing for ground-state enhancement, and therefore their results may include serious contamination from the excited-state component. The 3Q static potential can be measured with the 3Q Wilson loop, where the 3Q gauge-invariant state is generated at $`t=0`$ and is annihilated at $`t=T`$, as shown in Fig.1. Here, the three quarks are spatially fixed in $`𝐑^3`$ for $`0<t<T`$. The 3Q Wilson loop $`W_{3\mathrm{Q}}`$ is defined in a gauge-invariant manner as $$W_{3\mathrm{Q}}\frac{1}{3!}\epsilon _{abc}\epsilon _{a^{}b^{}c^{}}U_1^{aa^{}}U_2^{bb^{}}U_3^{cc^{}}$$ (1) with $`U_k\mathrm{P}\mathrm{exp}\{ig_{\mathrm{\Gamma }_k}𝑑x^\mu A_\mu (x)\}`$ ($`k=1,2,3`$). Here, $`P`$ denotes the path-ordered product along the path denoted by $`\mathrm{\Gamma }_k`$ in Fig.1. Similar to the derivation of the Q-$`\overline{\mathrm{Q}}`$ potential from the Wilson loop, the 3Q potential $`V_{3\mathrm{Q}}`$ is obtained as $`V_{3\mathrm{Q}}=lim_T\mathrm{}\frac{1}{T}\mathrm{ln}W_{3\mathrm{Q}}.`$ Physically, the true ground state of the 3Q system, which is of interest here, is expected to be expressed by the flux tubes instead of the strings, and the 3Q state which is expressed by the three strings generally includes many excited-state components such as flux-tube vibrational modes. Of course, if the large $`T`$ limit can be taken, the ground-state potential would be obtained. However, $`W_{3\mathrm{Q}}`$ decreases exponentially with $`T`$, and then the practical measurement of $`W_{3\mathrm{Q}}`$ becomes quite severe for large $`T`$ in lattice QCD simulations. Therefore, for the accurate measurement of the 3Q ground-state potential $`V_{3\mathrm{Q}}`$, the smearing technique for ground-state enhancement is practically indispensable. However, this smearing technique was not applied to the past lattice QCD studies for $`V_{3\mathrm{Q}}`$ in Refs. , since the smearing technique were mainly developed after their works. In this paper, we study the 3Q ground-state potential $`V_{3\mathrm{Q}}`$ using the ground-state enhancement by the gauge-covariant smearing technique for the link-variable in SU(3)<sub>c</sub> lattice QCD with the standard action with $`\beta `$=5.7 and $`12^3\times `$ 24 at the quenched level . We consider 16 patterns of the 3Q configuration where the three quarks are put on $`(i,0,0)`$, $`(0,j,0)`$ and $`(0,0,k)`$ in $`𝐑^3`$ with $`0i,j,k3`$ in the lattice unit. Here, the junction point $`O`$ is set at the origin $`(0,0,0)`$ in $`𝐑^3`$, although the final result of the ground-state potential $`V_{3\mathrm{Q}}`$ should not depend on the artificial selection of $`O`$. For actual lattice QCD calculations of the 3Q Wilson loop, we use the translational, the rotational and the reflection symmetries on lattices on the choice of the origin $`O`$ and the direction of $`\widehat{x},\widehat{y},\widehat{z}`$ The standard smearing for link-variables is expressed as the iteration of the replacement of the spatial link-variable $`U_i(s)`$ ($`i=1,2,3`$) by the obscured link-variable $`\overline{U}_i(s)\mathrm{SU}(3)_c`$ which maximizes $$\mathrm{Re}\mathrm{tr}\left\{\overline{U}_i^{}(s)\left[\alpha U_i(s)+\underset{ji}{}\{U_j(s)U_i(s+\widehat{j})U_j^{}(s+\widehat{i})+U_j^{}(s\widehat{j})U_i(s\widehat{j})U_j(s+\widehat{i}\widehat{j})\}\right]\right\}$$ (2) with the smearing parameter $`\alpha `$ being a real number. The $`n`$-th smeared link-variables $`U_\mu ^{(n)}(s)`$ $`(n=1,2,..,N_{\mathrm{smear}})`$ are iteratively defined starting from $`U_\mu ^{(0)}(s)U_\mu (s)`$ as $$U_i^{(n)}(s)\overline{U}_i^{(n1)}(s)(i=1,2,3),U_4^{(n)}(s)U_4(s).$$ (3) As an important feature, this smearing procedure keeps the gauge covariance of the “fat” link-variable $`U_\mu ^{(n)}(s)`$ properly. In fact, the gauge-transformation property of $`U_\mu ^{(n)}(s)`$ is just the same as that of the original link-variable $`U_\mu (s)`$, and therefore the gauge invariance of $`F(U_\mu ^{(n)}(s))`$ is ensured whenever $`F(U_\mu (s))`$ is a gauge-invariant function. Since the fat link-variable $`U_\mu ^{(n)}(s)`$ includes a spatial extension, the “line” expressed with $`U_\mu ^{(n)}(s)`$ physically corresponds to a “flux tube” with a spatial extension. Therefore, if a suitable smearing is done, the “line” of the fat link-variable is expected to be close to the ground-state flux tube. (This smearing technique is actually successful for the extraction of the Q-$`\overline{\mathrm{Q}}`$ potential in lattice QCD .) Now, we investigate the magnitude of the ground-state component in the 3Q-state operator at $`t=0,T`$ in the 3Q Wilson loop $`W_{3\mathrm{Q}}`$. From the similar argument in the Q-$`\overline{\mathrm{Q}}`$ system , the overlap of the 3Q-state operator with the ground state is estimated with $$C_0W_{3\mathrm{Q}}(T)^{T+1}/W_{3\mathrm{Q}}(T+1)^T,$$ (4) which we call the ground-state overlap. If the 3Q state at $`t=0,T`$ in Fig.1 is the perfect ground state, $`W_{3\mathrm{Q}}(T)=e^{V_{3\mathrm{Q}}T}`$ and then $`C_0`$=1 hold. (Here, $`W_{3\mathrm{Q}}(T)`$ in Eq.(1) is normalized as $`W_{3\mathrm{Q}}(T=0)=1`$.) The ground-state potential $`V_{3\mathrm{Q}}`$ can be measured accurately if $`C_0`$ is large enough. Then, we check the ground-state overlap $`C_0`$ of the 3Q Wilson loop $`W_{3\mathrm{Q}}(U_\mu ^{(n)}(s),T)`$ composed of the fat link-variable $`U_\mu ^{(n)}(s)`$ in lattice QCD simulations, and search reasonable values of the smearing parameter $`\alpha `$ and the iteration number $`N_{\mathrm{smear}}`$ for this purpose. In lattice QCD simulations, the ground-state overlap $`C_0`$ is largely enhanced as $`0.8<C_0<1`$ for $`T3`$ by the smearing with $`\alpha =2.3`$ and $`N_{\mathrm{smear}}=12`$ for all of the 3Q configurations in consideration, as shown in Fig.2. Here, we make a theoretical consideration on the potential form in the Q-$`\overline{\mathrm{Q}}`$ and 3Q systems with respect to QCD. In the short-distance limit, perturbative QCD is applicable and the Coulomb-type potential appears as the one-gluon-exchange (OGE) result. In the long-distance limit at the quenched level, the flux-tube picture would be applicable from the argument of the strong-coupling limit of QCD , and hence a linear-type confinement potential proportional to the total flux-tube length is expected to appear. Of course, it is nontrivial that these simple arguments on the ultraviolet and infrared limits of QCD hold for the intermediate region as $`0.2\mathrm{fm}<r<1\mathrm{fm}`$. Nevertheless, lattice QCD results for the Q-$`\overline{\mathrm{Q}}`$ ground-state potential is well described by $$V_{\mathrm{Q}\overline{\mathrm{Q}}}(r)=\frac{A_{\mathrm{Q}\overline{\mathrm{Q}}}}{r}+\sigma _{\mathrm{Q}\overline{\mathrm{Q}}}r+C_{\mathrm{Q}\overline{\mathrm{Q}}}$$ (5) at the quenched level . Actually, we measure $`V_{\mathrm{Q}\overline{\mathrm{Q}}}`$ from the on-axis Wilson loop with the smearing with $`\alpha =2.3`$ and $`N_{\mathrm{smear}}=20`$, and find a good fitting of Eq.(5) with the parameters ($`A_{\mathrm{Q}\overline{\mathrm{Q}}}`$, $`\sigma _{\mathrm{Q}\overline{\mathrm{Q}}}`$, $`C_{\mathrm{Q}\overline{\mathrm{Q}}}`$) listed in Table 1. (In general, the adequate values of $`\alpha `$ and $`N_{\mathrm{smear}}`$ depend on the operator.) In fact, $`V_{\mathrm{Q}\overline{\mathrm{Q}}}`$ is described by a sum of the short-distance OGE result and the long-distance flux-tube result. Also for the 3Q ground-state potential $`V_{3\mathrm{Q}}`$, we try to apply this simple picture of the short-distance OGE result plus the long-distance flux-tube result. Then, the 3Q potential $`V_{3\mathrm{Q}}`$ is expected to take a form of $$V_{3\mathrm{Q}}=A_{3\mathrm{Q}}\underset{i<j}{}\frac{1}{\|𝐫_i𝐫_j\|}+\sigma _{3\mathrm{Q}}L_{\mathrm{min}}+C_{3\mathrm{Q}},$$ (6) where $`L_{\mathrm{min}}`$ denotes the minimal value of total length of color flux tubes linking the three quarks. Denoting three sides of the 3Q triangle by $`a`$,$`b`$ and $`c`$, $`L_{\mathrm{min}}`$ is expressed as $$L_{\mathrm{min}}=\left[\frac{1}{2}(a^2+b^2+c^2)+\frac{\sqrt{3}}{2}\sqrt{(a+b+c)(a+b+c)(ab+c)(a+bc)}\right]^{\frac{1}{2}},$$ (7) when all angles of the 3Q triangle do not exceed $`2\pi /3`$. In this case, there appears the physical junction which connects the three flux tubes originating from the three quarks, and the shape of the 3Q system is expressed as a Y-type flux tube , where the angle between two flux tubes is found to be $`2\pi /3`$ . When an angle of the 3Q triangle exceeds $`2\pi /3`$, one finds $`L_{\mathrm{min}}=a+b+c\mathrm{max}(a,b,c)`$. Now, we show lattice QCD results for the 3Q ground-state potential. We generate 210 gauge configurations using SU(3)<sub>c</sub> lattice QCD Monte-Carlo simulation with the standard action with $`\beta =5.7`$ and $`12^3\times 24`$ at the quenched level. The lattice spacing $`a0.19\mathrm{fm}`$ is determined so as to reproduce the string tension as $`\sigma `$=0.89 GeV/fm in the Q-$`\overline{\mathrm{Q}}`$ potential $`V_{\mathrm{Q}\overline{\mathrm{Q}}}`$. Here, the gauge configurations are taken every 500 sweeps after a thermalization of 5000 sweeps using the pseudo-heat-bath algorithm. In this study, lattice QCD calculations have been performed on NEC-SX4 at Osaka University. We measure the 3Q ground-state potential $`V_{3\mathrm{Q}}`$ using the smearing technique, and compare the lattice data with the theoretical form of Eq.(6). Owing to the smearing, the ground-state component is largely enhanced, and therefore the 3Q Wilson loop $`W_{3\mathrm{Q}}`$ composed with the smeared link-variable exhibits a single-exponential behavior as $`W_{3\mathrm{Q}}e^{V_{3\mathrm{Q}}T}`$ even for a small value of $`T`$. Then, for each 3Q configuration, we extract $`V_{3\mathrm{Q}}^{\mathrm{latt}}`$ from the least squares fit with the single-exponential form $$W_{3\mathrm{Q}}=\overline{C}e^{V_{3\mathrm{Q}}T}$$ (8) in the range of $`T_{\mathrm{min}}TT_{\mathrm{max}}`$ listed in Table 2. Here, we choose the fit range of $`T`$ such that the stability of the “effective mass” $`V(T)\mathrm{ln}\{W_{3\mathrm{Q}}(T)/W_{3\mathrm{Q}}(T+1)\}`$ is observed in the range of $`T_{\mathrm{min}}TT_{\mathrm{max}}1`$. For each 3Q configuration, we summarize the lattice data $`V_{3\mathrm{Q}}^{\mathrm{latt}}`$ as well as the prefactor $`\overline{C}`$ in Eq.(8), the fit range of $`T`$ and $`\chi ^2/N_{\mathrm{DF}}`$ in Table 2. The statistical error of $`V_{3\mathrm{Q}}^{\mathrm{latt}}`$ is estimated with the jackknife method. We find again a large ground-state overlap as $`\overline{C}>0.8`$ for all 3Q configurations. Now, we consider the potential form. In Fig.3, we plot the 3Q ground-state potential $`V_{3\mathrm{Q}}`$ as the function of $`L_{\mathrm{min}}`$. Apart from a constant, $`V_{3\mathrm{Q}}`$ is almost proportional to $`L_{\mathrm{min}}`$ in the infrared region. We show in Table 1 the best fitting parameters in Eq.(6) for $`V_{3\mathrm{Q}}`$. On the goodness of this fitting, $`\chi ^2`$ seems relatively large as $`\chi ^2/N_{\mathrm{DF}}=3.76`$, which may reflect a systematic error on the finite lattice spacing. We add in Table 2 the comparison of the lattice data $`V_{3\mathrm{Q}}^{\mathrm{latt}}`$ with the fitting function $`V_{3\mathrm{Q}}^{\mathrm{fit}}`$ in Eq.(6) with the parameters listed in Table 1. The three-quark ground-state potential $`V_{3\mathrm{Q}}`$ is well described by Eq.(6) with accuracy better than a few %. Next, we compare the coefficients $`(\sigma _{3\mathrm{Q}},A_{3\mathrm{Q}},C_{3\mathrm{Q}})`$ in the 3Q potential $`V_{3\mathrm{Q}}`$ in Eq.(6) with $`(\sigma _{\mathrm{Q}\overline{\mathrm{Q}}},A_{\mathrm{Q}\overline{\mathrm{Q}}},C_{\mathrm{Q}\overline{\mathrm{Q}}})`$ in the Q-$`\overline{\mathrm{Q}}`$ potential $`V_{\mathrm{Q}\overline{\mathrm{Q}}}`$ in Eq.(5) as listed in Table 1. As a remarkable fact, we find a universal feature of the string tension, $`\sigma _{3\mathrm{Q}}\sigma _{\mathrm{Q}\overline{\mathrm{Q}}}`$, as well as the OGE result for the Coulomb coefficient, $`A_{3\mathrm{Q}}\frac{1}{2}A_{\mathrm{Q}\overline{\mathrm{Q}}}`$. As a check, we consider the diquark limit, where two quark locations coincide in the 3Q system. In the diquark limit, the static 3Q system becomes equivalent to the Q-$`\overline{\mathrm{Q}}`$ system, which leads to a physical requirement on the relation between $`V_{3\mathrm{Q}}`$ and $`V_{\mathrm{Q}\overline{\mathrm{Q}}}`$. Our results, $`\sigma _{3\mathrm{Q}}\sigma _{\mathrm{Q}\overline{\mathrm{Q}}}`$ and $`A_{3\mathrm{Q}}\frac{1}{2}A_{\mathrm{Q}\overline{\mathrm{Q}}}`$, are consistent with the physical requirement in the diquark limit. Here, the constant term is to be considered carefully in the diquark limit, because there appears a singularity or a divergence from the Coulomb term in $`V_{3\mathrm{Q}}`$ in the continuum diquark limit, $`lim_{𝐫_j𝐫_i}\frac{A_{3\mathrm{Q}}}{\|𝐫_i𝐫_j\|}=\mathrm{}`$. In the lattice regularization, this ultraviolet divergence is regularized to be a finite constant with the lattice spacing $`a`$ as $`\frac{A_{3\mathrm{Q}}}{\|𝐫_i𝐫_j\|}\frac{A_{3\mathrm{Q}}}{\omega a}`$, where $`\omega `$ is a dimensionless constant satisfying $`0<\omega <1`$ and $`\omega 1`$. Then, we find $`C_{3\mathrm{Q}}+\frac{A_{3\mathrm{Q}}}{\omega a}=C_{\mathrm{Q}\overline{\mathrm{Q}}}`$, i.e., $`C_{3\mathrm{Q}}C_{\mathrm{Q}\overline{\mathrm{Q}}}=\frac{A_{3\mathrm{Q}}}{\omega a}(>0)`$ in the diquark limit. This is the requirement for the constant term in the diquark limit on the lattice. Our lattice QCD results for $`C_{3\mathrm{Q}}`$, $`C_{\mathrm{Q}\overline{\mathrm{Q}}}`$ and $`A_{3\mathrm{Q}}`$ are thus consistent with this requirement, and one finds $`\omega 0.46`$. Finally, we also try to fit $`V_{3\mathrm{Q}}^{\mathrm{latt}}`$ with the $`\mathrm{\Delta }`$-type flux-tube ansatz of $`V_{3\mathrm{Q}}=A_\mathrm{\Delta }_{i<j}\frac{1}{\|𝐫_i𝐫_j\|}+\sigma _\mathrm{\Delta }_{i<j}\|𝐫_i𝐫_j\|+C_\mathrm{\Delta }`$, which was suggested in Refs.. However, this fitting seems rather worse, because $`\chi ^2`$ is unacceptably large as $`\chi ^2/N_{\mathrm{DF}}=10.1`$ even for the best fit with $`A_\mathrm{\Delta }=0.1410(64)`$, $`\sigma _\mathrm{\Delta }=0.0858(16)`$ and $`C_\mathrm{\Delta }=0.9344(210)`$ in the lattice unit. As an approximation, $`V_{3\mathrm{Q}}`$ seems described by a simple sum of the effective two-body Q-Q potential with a reduced string tension as $`\sigma _{\mathrm{QQ}}0.53\sigma `$. This reduction factor can be naturally understood as a geometrical factor rather than the color factor, since the ratio between $`L_{\mathrm{min}}`$ and the perimeter length $`L_P`$ satisfies $`\frac{1}{2}\frac{L_{\mathrm{min}}}{L_P}\frac{1}{\sqrt{3}}`$, which leads to $`L_{\mathrm{min}}\sigma =L_P\sigma _{\mathrm{QQ}}`$ with $`\sigma _{\mathrm{QQ}}=(0.50.58)\sigma `$. We have studied the static 3Q ground-state potential in SU(3) lattice QCD at the quenched level, using the smearing technique for ground-state enhancement. We have found that $`V_{3\mathrm{Q}}`$ is well described by a sum of a constant, the two-body Coulomb term and the three-body linear confinement term $`\sigma _{3\mathrm{Q}}L_{\mathrm{min}}`$, where $`L_{\mathrm{min}}`$ denotes the minimal length of the color flux tube linking the three quarks. We have also found a universal feature of the string tension as $`\sigma _{3\mathrm{Q}}\sigma _{\mathrm{Q}\overline{\mathrm{Q}}}`$, and the OGE result for Coulomb coefficients as $`A_{3\mathrm{Q}}\frac{1}{2}A_{\mathrm{Q}\overline{\mathrm{Q}}}`$. In lattice QCD studies, however, there appear systematic errors relating to the finite lattice-spacing effect and so on. To obtain the conclusive result on the 3Q potential, we are investigating with finer lattices with large $`\beta `$’s.
warning/0006/hep-ph0006113.html	ar5iv	text	# 1 Introduction ## 1 Introduction Undoubtedly the axial anomaly represents a fundamental issue for understanding the basic aspects of quantum field theory. This issue has been analysed deeply over the years. The anomaly problem has been treated by means of renormalization procedure, giving the interpretation of its origin in terms of ultraviolet divergences . A more formal analysis of the axial anomaly can be made by using the path integral formalism . Dolgov and Zakharov have shown an alternative approach to the axial anomaly, by studying the $`VVA`$ triangular diagram through dispersion relations. From this approach follows the interpretation of the axial anomaly as an infrared phenomenon. It appears as due to a singularity present in the chiral limit in the absorbitive part of the triangular diagram. The infrared aspect of the axial anomaly, rised in this paper, is complementary to the more familiar ultraviolet one, which emerges from the renormalization procedure. A particularly interesting feature of this approach is that it allows to shed light upon the physical meaning of the anomalous chiral symmetry breaking, which is connected to a non conservation of helicity. The connection between the anomaly and the breaking of a given symmetry has received a lot of attention in the literature and this subject has been discussed and developed in several papers. Gribov , in a seminal work, has described the source of the anomalies as a collective motion of particles with arbitrarily large momenta in the vacuum. Related to this work, in ref. , Mueller has discussed the manifestation of the axial anomaly as a flow of Landau levels. In the papers of refs. the origin of the axial anomaly has been studied in two dimensions and, again, as a level crossing phenomenon. The infrared interpretation of the axial anomaly, according to the Dolgov and Zakharov approach, has been possibly advanced in . In a series of papers the leading terms in the chiral limit have been correctly evaluated. Furthermore the dispersive analysis of the triangular $`VVA`$ diagram is fundamental in the formulation of the ’t Hooft consistency condition . The axial anomaly plays an essential role also in the interpretation of spin dependent parton distribution (see ). In this work we attempt to relate the Dolgov and Zakharov approach to the axial anomaly to some effects in the dynamics of physical reactions, as the radiative decays of $`\pi ^+`$ and $`Z^0`$. We will try to show that the axial anomaly can be related to polarized radiative decays, as in the usual ultraviolet interpretation it is connected to the $`\pi ^0\gamma \gamma `$ decay. We calculate the corresponding decay rates for the cases where the outgoing leptons are in a definite helicity state and we examine in some detail the structure of the mass singularities and their cancellation. We study how the Kinoshita-Lee-Nauenberg (KLN) theorem applies and we consider the analogies and the differences with respect to the corresponding unpolarized decay rates. The paper is organized as follows. In the section 2 we briefly reconsider the dispersive approach to the axial anomaly. In particular we concentrate on its physical origin. In section 3 we extend the Dolgov and Zakharov approach to the study of the radiative pion decay. We calculate the differential decay rate for the process, where the outgoing lepton undergoes an helicity flip and we interpret its behaviour in the chiral limit, as a manifestation of the axial anomaly. In section 4 we study the behaviour of the Inner Bremsstrahlung contribution to the pion decay rate in the collinear and infrared limits. We consider separately the unpolarized decay rate and both the cases of right-handed and of left-handed outgoing lepton. We find that the mass singularities cancellation mechanism occurs in different ways, according to the polarization of the outgoing lepton. We discuss the various realizations of the KLN theorem. We also consider a more general process, i.e. the radiative $`Z^0`$ decay in a lepton-antilepton pair, with a right-handed polarized lepton. We examine how the mass singularities cancel in this case and discuss the differences with respect to the pion decay. Finally, in the Conclusions, we summarize our arguments. A short discussion of the main results obtained has been already given in . ## 2 The dispersive approach to axial anomaly The Dolgov and Zakharov approach to derive the axial anomaly is based on the dispersion relation method. In this framework, the triangular diagram with two vector and one axial vertices is seen as the lowest order contribution to the process: $$axialvectorsource\mathrm{\hspace{0.33em}2}realphotons,$$ as described by the diagrams of fig. 1. In the physical region, the triangular diagram possesses a branch cut along the real axis, from $`4m^2`$ to infinity. $`T_{\alpha \beta \mu }^5`$ represents the corresponding amplitude. We express it by means of a dispersion relation in the variable $`s=q^2=(k_1+k_2)^2`$, where $`k_1`$ and $`k_2`$ are the outgoing photon momenta. By requiring parity, Lorentz invariance, Bose symmetry and that $`T_{\alpha \beta \mu }^5`$ satisfies the vector Ward-Takahashi identity, $$k_1^\alpha T_{\alpha \beta \mu }^5(k_1,k_2)=k_1^\beta T_{\alpha \beta \mu }^5(k_1,k_2)=0,$$ we can write it in terms of an invariant scalar function $`g_1(q^2)`$ as: $$\mathrm{T}_{\alpha \beta \mu }^5(k_1,k_2)=\frac{2\alpha }{\pi }g_1(q^2)ϵ_{\alpha \beta \sigma \rho }k_1^\sigma k_2^\rho q_\mu .$$ (1) Similarly, we can express the contribution of the triangular diagram with the vertex $`\gamma _\mu \gamma _5`$ substituted by $`\gamma _5`$ as: $$\mathrm{T}_{\alpha \beta }^5(k_1,k_2)=\frac{\alpha m}{\pi }g_2(q^2)ϵ_{\alpha \beta \sigma \rho }k_1^\sigma k_2^\rho ,$$ (2) where $`g_2(q^2)`$ is another invariant scalar function. In terms of dispersion relations, the functions $`g_i(q^2)`$, $`i=1,2`$ can be expressed as follows: $$g_i(q^2)=\frac{1}{\pi }_{4m^2}^+\mathrm{}ds\frac{\mathrm{Im}g_i(s)}{sq^2}i=1,2.$$ (3) We can use unsubtracted dispersion relations, because the integrals contained in $`g_i(q^2)`$, $`i=1,2`$ are convergent, since these functions are multiplied by three and two powers of momentum, respectively. The imaginary part of the invariant scalar functions can be derived from the absorbitive part of the triangular diagram, calculated by cutting the diagram as shown in fig. 2 and by using the Cutkosky rules or the perturbative unitarity relation. We obtain : $`\mathrm{Im}g_1(q^2)`$ $`=`$ $`{\displaystyle \frac{2\pi m^2}{q^4}}\theta (q^24m^2)\mathrm{ln}\left({\displaystyle \frac{1+\sqrt{14m^2/q^2}}{1\sqrt{14m^2/q^2}}}\right).`$ (4) $`\mathrm{Im}g_2(q^2)`$ $`=`$ $`2\pi \theta (q^24m^2){\displaystyle \frac{1}{q^2}}\mathrm{ln}\left({\displaystyle \frac{1+\sqrt{14m^2/q^2}}{1\sqrt{14m^2/q^2}}}\right)`$ (5) Using eqs. (1), (2), (3) and the above expressions, we derive the complete triangular diagram contribution: $$\mathrm{T}_{\alpha \beta \mu }^5(k_1,k_2)=\frac{2m}{q^2+iϵ}\mathrm{T}_{\alpha \beta }^5(k_1,k_2)q_\mu +\frac{2\alpha }{\pi }\frac{1}{q^2+iϵ}ϵ_{\alpha \beta \sigma \rho }k_1^\sigma k_2^\rho q_\mu $$ (6) and the anomalous axial-vector Ward-Takahashi identity: $$q^\mu \mathrm{T}_{\alpha \beta \mu }^5(k_1,k_2)=\frac{2m}{q^2}\mathrm{T}_{\alpha \beta }^5(k_1,k_2)+\frac{2\alpha }{\pi }ϵ_{\alpha \beta \sigma \rho }k_1^\sigma k_2^\rho .$$ (7) Thus, in the dispersive approach we find the same result obtained by using the renormalization procedure . The method above allows a more direct treatment, since we avoid evaluating divergent integrals and introducing regularization schemes. The fact that the axial anomaly can be derived without using the renormalization procedure, suggests that this should not be considered as the only origin of the anomalous breaking of the chiral symmetry. Moreover, by studying the axial anomaly with the renormalization procedure, that is by considering its ultraviolet interpretation, an important aspect of this phenomenon remains obscure and we are bound by a formal derivation only. As stressed in and , the dispersive approach shows that the anomaly is related to the chiral limit and therefore, it can be interpreted as an infrared effect. In this work we are close to this infrared interpretation, which, as we will see, can help to shed light on some aspects of the physics connected to the axial anomaly. Let us briefly discuss the physics involved in the amplitudes contributing to the absorbitive part of the triangular diagram. The two Born diagrams, obtained from the cut of the triangular diagram (see fig. 1), describe the following two processes: 1. the production of a fermion-antifermion pair (for example $`e^+e^{}`$) by an axial-vector source; 2. the subsequent annihilation of the pair into two real photons. In both these processes there occur helicity flips, thus the chirality is not conserved in the zero mass limit. Let us go to the center of mass frame of the two final photons. In the first process the axial-vector source produces an $`e^+e^{}`$ pair of total spin zero, since a spin 1 state cannot annihilate into a two real photons state. Thus $`e^+`$ and $`e^{}`$ must have the same helicity and hence opposite chirality in the massless limit<sup>1</sup><sup>1</sup>1 Equivalently, considering the outgoing antifermion as an incoming fermion, one can say that the latter makes an helicity flip interacting with the axial-vector source.. In the process b) the $`e^+e^{}`$ pair annihilates into two real photons by going through an intermediate virtual state. There are four possible virtual states : one is drawn in fig. 3 ($`𝐩_1`$ and $`𝐩_2`$ are the $`e^{}`$ and $`e^+`$ linear momenta), a second one is obtained by reversing the virtual state helicity; the remaining two are obtained replacing the outgoing virtual fermion with an incoming virtual antifermion. Let us study the case shown in fig. 3 and assume that $`e^+`$ and $`e^{}`$ are both right-handed. In the vertex B the chirality is conserved in the massless limit, while in the vertex A there is an helicity flip, thus the chiral symmetry is broken. As can be easily checked, for all the remaining virtual states we always have an allowed vertex and a forbidden one. At the Born level, these reactions are described by the classical QED Lagrangian, which, in the $`m0`$ limit, is invariant under chiral transformations. Thus it seems at first sight that the absorbitive part of the triangular diagram, being proportional to the product of amplitudes relative to processes forbidden by chiral invariance, vanishes. On the contrary, one sees that taking the limit $`m0`$ in (5) gives a finite result : $$\mathrm{Im}g_1(q^2)\pi \delta (q^2)\mathrm{as}m0.$$ (8) Therefore, by studying the absorbitive part of the triangular diagram, one establishes that a non-conservation of helicity occurs, becoming, in the massless limit, a non-conservation of chirality. We interpret this as related to the presence of the anomalous term in the divergence of the axial-vector current. Thus the axial anomaly can be derived by studying the properties of the amplitude in the infrared region. Even if in this work we will analyse the cancellation of mass singularities in polarized processes, we will not discuss the physical implication of the zero fermion mass limit. As stated in refs. , the result in eq. (8) indicates that there occurs a cancellation of the suppression factor $`m^2`$, due to the terms coming from the vertices with helicity flip. In (4) the logarithmic factor conspires to give a finite result. This logarithm is a collinear one; we shall discuss about this kind of logarithms in section 4. Its presence is a manifestation of the singularity occurring in the fermion propagator as $`m0`$, which exactly cancels the suppression factor $`m^2`$ in the numerator. We observe that the behaviour of the absorbitive part of the triangular diagram given in eq. (8) shows that the $`m0`$ limit is not smooth; if we evaluate the amplitudes with the massless theory, they identically vanish. If, on the contrary, we take the chiral limit after summing their product over the intermediate states, we obtain a result different from zero. ## 3 Helicity changing processes and polarized pion decay ### 3.1 Helicity changing processes In the process b) contributing to the absorbitive part of the triangular diagram a charged polarized fermion, changes helicity by emitting a photon. Thanks to the collinear singularity within the propagator relative to this fermion, the chiral suppression factor $`m^2`$ is cancelled and the absorbitive part of the triangular diagram doesn’t vanish in the chiral limit. Let us now extend these remarks to physical processes having the same features as the ones characterizing the absorbitive part of the triangular diagram. This means that we want to investigate about a manifestation of the axial anomaly in reactions where a fermion changes helicity by emitting a photon. We call these processes helicity changing processes. They have been considered by Lee and Nauenberg , who observed that states with opposite helicity don’t decouple in the massless limit, provided that this limit is taken after having summed the transition probability over the final phase space. In the process b) the incoming $`e^+`$ and $`e^{}`$ are in a definite helicity state, as discussed above. In the reactions we want to study, we calculate the probability that an outgoing fermion assumes an helicity opposite to the one required by the interaction before the emission of the photon. In other words, we evaluate the probability that in the fermion-photon vertex occurs an helicity flip. According to the Dolgov and Zakharov analysis, we interpret the presence of a term independent of the fermion mass in the corresponding cross sections as a manifestation of the axial anomaly. Due to this term, the probability for a process with helicity flip doesn’t vanish in the chiral limit. ### 3.2 Polarized radiative pion decay We first examine the non radiative pion decay $$\pi ^+l^++\nu _l,$$ where $`l^+`$ is an antilepton ($`e^+`$ or $`\mu ^+`$) and $`\nu _l`$ is the associated neutrino. At the Born level the total decay rate is given by: $$\mathrm{\Gamma }_0(\pi ^+l^+\nu _l)=\frac{G^2f_\pi ^2}{8\pi }V_{ud}^2\frac{m_l^2}{m_\pi ^3}(m_\pi ^2m_l^2)^2,$$ (9) where $`G`$ is the Fermi coupling constant, $`f_\pi `$ is the pion decay constant, $`V_{ud}`$ is the CKM matrix element, $`m_l`$ and $`m_\pi `$ are the lepton and pion masses, respectively. The decay rate (9) is proportional to $`m_l^2`$, since, due to angular momentum conservation, the pion produces a left-handed lepton, while the structure of the weak coupling requires the $`l^+`$ to be right-handed for $`m_l=0`$. This situation is confirmed by the expression of the total decay rate for the process in which the antilepton is polarized: $`\mathrm{\Gamma }_0^{pol}(\pi ^+l^+\nu _l)={\displaystyle \frac{G^2f_\pi ^2}{16\pi }}V_{ud}^2{\displaystyle \frac{m_l^2}{m_\pi ^3}}(m_\pi ^2m_l^2)^2\left(1{\displaystyle \frac{𝐩_l𝐬_𝐥}{\|𝐩_l\|}}\right),`$ (10) where $`𝐩_l`$ and $`𝐬_𝐥`$ are the linear momentum and the spin vector of the lepton, respectively, giving $`𝐩_l𝐬_l/\|𝐩_l\|=\pm 1`$, for right-handed and left-handed lepton, respectively. The eq. (10) indicates that the lepton is mandatory left-handed, as requested by angular momentum conservation. We now consider the radiative correction, due to the emission of a real photon, to the pion decay, $$\pi ^+l^++\nu _l+\gamma ,$$ when the outgoing antilepton is polarized, i.e. when it is in a definite helicity state. As we have seen in the non-radiative case, due to angular momentum conservation, the $`\pi ^+`$ is coupled to a left-handed lepton. We calculate the probability that the lepton flips its helicity and becomes right-handed, by emitting a real photon. The amplitude describing the radiative pion decay can be divided into two parts, the Inner Bremsstrahlung and the Structure Dependent amplitudes : $$M(\pi ^+l^+\nu _l\gamma )=M_{IB}+M_{SD}.$$ (11) The Inner Bremsstrahlung amplitude, where the photon is radiated from the external charged particles, can be calculated using the rules of QED, with a point like pion coupling; the Structure-Dependent amplitude is governed by the strong interactions. Clearly, the relevant part for the problem we are considering is the Inner Bremsstrahlung contribution described by the diagrams of fig. 4. The term associated to the $`IB_3`$ diagram is the so called contact term and it is introduced to ensure gauge invariance (see, for example, ). We consider only tree level diagrams, since, to reveal the effects of the axial anomaly, it is sufficient to take into account the contribution of the vertex with the helicity flip. For the moment, we may neglect the corrections due to the emission of virtual photons; these will be discussed in section 4. Finally, at this order in perturbation theory, we retain terms of all powers in the lepton mass. We will argue that the manifestation of the axial anomaly is strictly connected to these terms. The Inner Bremsstrahlung amplitude is given by : $`M_{IB}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}IB_i`$ (12) $`=`$ $`ie{\displaystyle \frac{G}{\sqrt{2}}}f_\pi V_{ud}m_l\overline{u}(p_\nu )(1+\gamma _5)\left[{\displaystyle \frac{pϵ}{pk}}{\displaystyle \frac{\overline{)}k\overline{)}ϵ+2p_lϵ}{2p_lk}}\right]v(p_l,s_{lR}),`$ where $`u`$ and $`v`$ are the Dirac spinors for the neutrino and the lepton, respectively, $`p`$ is the pion momentum, $`p_l`$ and $`s_{lR}`$ are the momentum and the polarization vector of the lepton, $`k`$ and $`ϵ`$ are the momentum and the polarization vector of the photon. We see that $`M_{IB}`$ is proportional to the lepton mass $`m_l`$, thus the decay rate is proportional to $`m_l^2`$. As we have said above, this factor is a consequence of the structure of the weak coupling. Thus, if we remove it by normalizing the radiative decay rate with respect to the non radiative one, we can emphasize the mass dependence of the radiation emission process. The differential Inner Bremsstrahlung contribution for right-handed lepton is: $`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{\mathrm{d}\mathrm{\Gamma }_{IB}^R}{\mathrm{d}y}}={\displaystyle \frac{\alpha }{4\pi }}{\displaystyle \frac{1}{(1r)^2}}{\displaystyle \frac{1}{A(1y+r)}}\{2A[1+y^22A+r(2A+r6)]`$ $`+\left[(A+2r)(1+r)^2+Ay(y4r)+2ry(1+y)y(1+y^2+5r^2)\right]\mathrm{ln}{\displaystyle \frac{y+A}{yA}}`$ $`+(1y+r)^2(y2rA)\mathrm{ln}{\displaystyle \frac{y+A2}{yA2}}\}.`$ (13) The dimensionless variable $`y`$ is defined as $$y=\frac{2E_l}{m_\pi }$$ and $$r=\frac{m_l^2}{m_\pi ^2},A=\sqrt{y^24r}.$$ The physical region for $`y`$ is: $$2\sqrt{r}y1+r.$$ (14) We see that, in eq. (13), there is a term independent of the lepton mass, the one in the first square brackets. Owing to this term, the differential decay rate doesn’t vanish in the limit $`m_l0`$. We have: $$\frac{1}{\mathrm{\Gamma }_0}\frac{\mathrm{d}\mathrm{\Gamma }_{IB}^R}{\mathrm{d}y}\frac{\alpha }{2\pi }(1y)\mathrm{as}r0.$$ (15) This indicates that there occurs an helicity flip and thus a chirality non conservation in the limit of zero lepton mass. According to the interpretation given above, this term can be interpreted as connected to the axial anomaly. It corresponds to the anomalous term present in the divergence of the axial current. Since the polarized radiative process with the right-handed $`l^+`$ is forbidden, in the limit $`m_l0`$, by the chiral invariance of the massless QED Lagrangian, the appearance of a term different from zero, in this limit, indicates the action of a cancellation mechanism, analogous to the one acting in the absorbitive part of the triangular diagram. To see how this mechanism acts, let us examine the decay rate differential with respect to the lepton energy and to the emission angle: $`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{\mathrm{d}\mathrm{\Gamma }_{IB}^R}{\mathrm{d}y\mathrm{d}\mathrm{cos}\theta }}={\displaystyle \frac{\alpha }{4\pi }}{\displaystyle \frac{1}{(1r)^2}}{\displaystyle \frac{1}{(1y+r)}}\{{\displaystyle \frac{4r}{(yA\mathrm{cos}\theta )^2}}(1+y^2yA)`$ $`+{\displaystyle \frac{4r^2}{(yA\mathrm{cos}\theta )^2}}(A+ry2)+y(1+r)A(1r)4r`$ $`+\left[(A+2r)(1+r)^2+Ay(y4r)+2ry(1+y)y(1+y^2+5r^2)\right]{\displaystyle \frac{1}{(yA\mathrm{cos}\theta )}}`$ $`+(1y+r)^2(yA2r){\displaystyle \frac{1}{(yA\mathrm{cos}\theta 2)}}\}.`$ (16) The first term is proportional to the lepton propagator squared. Thanks to this term, when we carry out the final integration over the emission angle, we obtain, besides logarithmic collinear divergences, also power collinear divergences. These power terms are essential for the convergent behaviour of the distribution. The cancellation of the chiral suppression would not take place were these terms absent. The term related to the axial anomaly (see eq. (15)) originates from this cancellation. To obtain the differential decay rate, we have to evaluate integrals of the form: $`m_l^2{\displaystyle _1^{+1}}\mathrm{d}\mathrm{cos}\theta {\displaystyle \frac{1}{(E\sqrt{E^2m_l^2}\mathrm{cos}\theta )^2}}=`$ $`{\displaystyle \frac{m_l^2}{\sqrt{E^2m_l^2}}}{\displaystyle \frac{2\sqrt{E^2m_l^2}}{m_l^2}},`$ (17) where the propagator has a power mass singularity that exactly cancels the factor $`m_l^2`$, coming from the vertex lepton-photon. The integration of the terms in equation (16) containing the lepton propagator gives rise to the collinear logarithms: $$\mathrm{ln}\frac{E+\sqrt{E^2m_l^2}}{E\sqrt{E^2m_l^2}}\mathrm{ln}\frac{m_l}{E}.$$ (18) In the differential decay rate (13) there are also terms proportional to the logarithm $$\mathrm{ln}\frac{E+\sqrt{E^2m_l^2}m_\pi }{E\sqrt{E^2m_l^2}m_\pi }.$$ (19) This one is another collinear logarithm; it diverges in the limit $`m_\pi 0`$ and $`m_l0`$ and corresponds to the possibility that the photon is emitted parallel to the pion. Finally, we point out that the chiral limit ($`m_l0`$) is not smooth. In fact we get different results depending upon whether we describe an helicity changing process using the massless theory or we take the $`m_l0`$ limit, after carrying out the integration over the final phase space. The radiative pion decay is not a good process to see this, since, owing to the angular momentum conservation in the pion vertex, this process cannot take place within a massless theory, given the structure of the $`(VA)`$ coupling of the electroweak theory. To see that the chiral limit is not a smooth one, we consider the scattering of a polarized electron with a proton (treated as a point like particle) of initial and final momenta $`q`$ and $`q^{}`$ respectively, accompanied by the emission of a real photon. We consider a left-handed incoming electron with momentum $`p`$ and spin $`s_L`$: we calculate the probability that the electron makes an helicity flip, emitting a real photon with momentum $`k`$ and polarization $`ϵ`$ and thus becoming right-handed electron with momentum $`p^{}`$ and spin $`s_R`$: $$e(p,s_L)+p(q)e(p^{},s_R)+p(q^{})+\gamma (k).$$ We study this process with the massless QED. The left-handed and right-handed spinors are given respectively by: $`u(p,s_L)`$ $`=`$ $`{\displaystyle \frac{1\gamma _5}{2}}u(p)`$ $`u(p^{},s_R)`$ $`=`$ $`{\displaystyle \frac{1+\gamma _5}{2}}u(p^{}).`$ The corresponding transition amplitude identically vanishes: $`M(e_lpe_Rp\gamma )`$ $`=`$ $`{\displaystyle \frac{e^3}{[(p^{}+k)^2+iϵ](l^2+iϵ)}}`$ (20) $`\overline{u}(p^{}){\displaystyle \frac{1\gamma _5}{2}}\overline{)}ϵ(\overline{)}p^{}+\overline{)}k)\gamma _\rho {\displaystyle \frac{1\gamma _5}{2}}u(p)\overline{u}(q^{})\gamma ^\rho u(q)=0,`$ since it contains the product of different chirality projectors. We now calculate the cross sections for processes with helicity flip using the massive QED and then we take the massless limit after having summed the transition probability over the final phase space. An example can be found in , where the cross section for the process $$e^{}(p_{},\lambda )+A(p)e^{}(p_{}^{},\lambda ^{})+\gamma (k,\lambda _\gamma )+B(q_i)$$ (21) is calculated. Here $`A`$ is the target (for example another fermion), $`\gamma `$ is a bremsstrahlung photon, assumed almost collinear with respect to the direction of the incident electron and $`B`$ is a set of particles produced in the reaction. The helicity flip cross section is given by: $$\frac{\mathrm{d}\sigma _{hf}}{\mathrm{d}x}=\sigma _0\left(s(1x)\right)\frac{\alpha }{2\pi }x,$$ (22) where $`x=k_0/E`$, $`E`$ is the energy of the incoming electron and $`\sigma _0`$ is the cross section for the Born process $$e^{}(p_{}k,\lambda )+A(p)e^{}(p_{}^{},\lambda ^{})+B(q_i).$$ (23) We see that the expression (22) does not vanish in the massless limit. The result in eq. (22) coincides with the one in eq. (15) at leading order. Indeed, in the subdominant terms are not accounted for. As it will become apparent in the following section, these terms are essential for the cancellation of mass singularities. In ref. , the authors give a different interpretation of the helicity changing processes and hence come to a different conclusion about the smoothness of the zero mass limit. ## 4 Mass singularities As it is well know, there are two types of divergences occurring in a theory when the mass of a particle goes to zero, which will be comprehensively call in the following mass singularities. The first type of divergences appears when we reach the phase space region, where the momentum of the massless particles vanishes: these are called infrared divergences. They occur for example in QED when the energy of the photon goes to zero. The Block-Nordsieck theorem assures that the infrared divergences cancel out in any inclusive cross section. The other type of mass singularities occurs in theories with massless coupled particles, like in QED when the photon couples to a fermion, in the limit of zero fermion mass. The origin is purely kinematical: when two massless particles, say with momenta $`k`$ and $`k^{}`$, move parallel to each other, they have combined invariant mass equal to zero: $$q^2=(k+k^{})^2=2EE^{}(1\mathrm{cos}\theta )0\mathrm{as}\theta 0,$$ (24) even though neither $`k`$ nor $`k^{}`$ are soft. These divergences are called collinear singularities. If we keep the fermion mass finite and integrate over the photon emission angle, the collinear divergence doesn’t occur, but the possibility of a divergence in the limit $`m0`$ results in the presence of the collinear logarithms, that is logarithms of the form $`\mathrm{ln}(E/m)`$, diverging for $`m0`$. In the case of collinear singularities, the theorem by Kinoshita, Lee and Nauenberg guarantees that these divergences cancel out if we sum the transition probability over the set of degenerate states, order by order in perturbation theory. This cancellation mechanism is analogous to that of the infrared divergences, as stated by the Block-Nordsieck theorem. Both types of mass singularities arise because the states of a theory with massless particles are highly degenerate. The infrared divergences can be interpreted as a consequences of the fact that a state with a single charged particle is degenerate with a state made of the same particle plus a number of soft photons; this correspond to the impossibility of distinguishing experimentally a charged particle from one accompanied by soft photons, owing to the finite resolution of the measurement apparatus. The situation of the collinear singularities is analogous: the state with a massless charged particle is degenerate with the states containing the same particle and a number of collinear photons. This corresponds to the fact that, as a consequence of the finite angular resolution, we cannot establish if a massless charged particle is accompanied by collinear photons. Let us now discuss the structure of mass singularities and their cancellation in the decay rates for the radiative pion and $`Z^0`$ decays and how the KLN theorem applies to this cases. ### 4.1 The pion case It is useful to separate the cases of unpolarized, right-handed and left-handed outgoing lepton. Let us consider first the mass singularity cancellation mechanism for the familiar case of unpolarized radiative $`\pi ^+`$ decay to the first order in $`\alpha `$. The differential Inner Bremsstrahlung contribution is given by: $`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{\mathrm{d}\mathrm{\Gamma }_{IB}}{\mathrm{d}y}}`$ $`=`$ $`{\displaystyle \frac{\alpha }{4\pi }}{\displaystyle \frac{1}{(1r)^2}}{\displaystyle \frac{1}{(1y+r)}}\{4A(r1)+[(1+r)^2+y(y4r)]\mathrm{ln}{\displaystyle \frac{y+A}{yA}}`$ (25) $``$ $`(1y+r)^2\mathrm{ln}{\displaystyle \frac{y+A2}{yA2}}\}.`$ One can easily see that the eq. (25) is divergent both in the collinear and in the infrared limits. The coefficients of the collinear logarithms don’t go to zero in the limit $`r0`$. There are also infrared divergences, because if we let $`y`$ reach its kinematical limit $`y^{MAX}=1+r`$, corresponding to the photon energy going to zero, the expression (25) diverges. The decay rate is made free from mass singularities in the ordinary way: the divergences cancellation occurs in the total inclusive decay rate, when we add all the first order contributions to the perturbative expansion, i.e. those relative to real and virtual photon emission. The diagrams describing the real photon emission contribution were already given in fig. 4; the diagram for the virtual correction is drawn in fig. 5. The expression for the lepton energy spectrum, including the Inner Bremsstrahlung contribution and the virtual photon one, calculated to the leading order in $`m_l/m_\pi `$, is given by : $$\frac{1}{\mathrm{\Gamma }_0}\frac{\mathrm{d}\mathrm{\Gamma }}{\mathrm{d}y}=D(y,r)[1+\frac{\alpha }{\pi }K_l(y)].$$ (26) $`D(y,r)`$ is the lepton distribution function, given, to the first order in $`\alpha `$, by: $$D(y,r)=\delta (1y)+\left[\frac{\alpha }{\pi }(L1)+O(\alpha ^2)\right]P^{(1)}(y),$$ (27) where $`L`$ is the logarithm $$L=\mathrm{ln}\frac{m_\pi }{m_l},$$ (28) diverging in the collinear limit; if the lepton is an electron, $`L5.6`$. $`P^{(1)}(y)`$ is the Gribov-Lipatov-Altarelli-Parisi kernel , which can be expressed in the form: $$P^{(1)}(y)=\frac{1+y^2}{1y}\delta (1y)_0^1dz\frac{1+z^2}{1z}.$$ (29) $`K_l(y)`$ is a finite term, free from infrared and collinear singularities, which has the expression: $$K_e(y)=1y\frac{1}{2}(1y)\mathrm{ln}(1y)+\frac{1+y^2}{1y}\mathrm{ln}y.$$ (30) The differential decay rate to order $`\alpha `$ therefore becomes: $$\frac{1}{\mathrm{\Gamma }_0}\frac{\mathrm{d}\mathrm{\Gamma }}{\mathrm{d}y}=\delta (1y)+\frac{\alpha }{\pi }(L1)P^{(1)}(y)+\frac{\alpha }{\pi }K_e(y).$$ (31) To calculate the inclusive decay rate, we have to integrate the expression in eq. (31) over $`y`$; to the leading order in the lepton mass, the physical region for $`y`$ is $`0y1`$. The kernel $`P^{(n)}`$ has the property that: $$_0^1dyP^{(n)}(y)=0;$$ (32) thus, when we calculate the inclusive decay rate, the coefficient of the collinear logarithm vanishes and the resulting expression is finite in the zero mass limit. Carrying out the integration over $`y`$, we obtain the well known inclusive decay rate to order $`\alpha `$ $$\frac{\mathrm{\Gamma }}{\mathrm{\Gamma }_0}=1+\frac{\alpha }{\pi }[\frac{15}{8}\frac{\pi ^2}{3}].$$ (33) As expected, the expression (33) is finite in the collinear limit and is also free from infrared divergences, because, as usual, the infrared divergences present in the soft photon contribution and in the virtual photon contribution have cancelled each other. Let us now discuss the mass singularities in the case of the right-handed Inner Bremsstrahlung contribution, given in eq. (13). It is easy to see that $`\mathrm{d}\mathrm{\Gamma }_{IB}^R/\mathrm{d}y`$ is finite in the limit $`r0`$, i.e. it is free from collinear singularities. In this limit the coefficients of both the collinear logarithms vanish. Indeed, as we have observed in section 3.2, only the term related to the axial anomaly survives in the zero mass limit, that is the part of the decay rate independent of the lepton mass. We observe that the right-handed Inner Bremsstrahlung contribution is free from infrared divergences as well. If we make the lepton energy $`y`$ reach its kinematical limit $`y^{MAX}`$, we obtain a finite result: $`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{\mathrm{d}\mathrm{\Gamma }_{IB}^R}{\mathrm{d}y}}0\mathrm{as}yy^{MAX}.`$ (34) The result (34) shows that the soft photon emission does not contribute to the radiative $`\pi ^+`$ decay with a right-handed lepton. This is a consequence of the fact that the soft photon contribution factorizes with respect to the Born decay rate, but this vanishes in the case of right-handed $`l^+`$ (see eq. (10)). Physically, eq. (34) is due to the fact that soft photons don’t carry spin, thus they cannot contribute to the angular momentum balance; therefore the process with the right-handed lepton emitting a soft photon is forbidden by angular momentum conservation. For the same reason of angular momentum conservation, in the right-handed case also the virtual contribution identically vanishes. The virtual photon diagram (see fig. 5) interferes with the Born one; the corresponding correction to the total decay rate for $`\pi l\nu _l`$ was calculated long ago by Kinoshita and is given by: $`\mathrm{\Gamma }_v={\displaystyle \frac{\alpha }{2\pi }}\left\{3\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }}{m_\pi }}{\displaystyle \frac{1}{2}}b(r)\left[4\mathrm{ln}{\displaystyle \frac{\lambda }{m_\pi }}\mathrm{ln}r+3\right]+{\displaystyle \frac{r}{1r}}\mathrm{ln}r+1\right\}\mathrm{\Gamma }_0,`$ (35) where $$b(r)=\frac{1+r}{1r}\mathrm{ln}r+2.$$ $`\mathrm{\Lambda }`$ is the ultraviolet cutoff and $`\lambda `$ is the infrared one. As usual, the virtual correction is factorized with respect to the Born decay rate, but, as we have already seen, if the lepton is right-handed, this is identically zero. In the right-handed case, the mass singularities cancellation occurs trough a mechanism different from the one working in the unpolarized decay rate. The infrared and the collinear limits give separately a finite result. In particular, the coefficient of the collinear logarithms is the lepton mass, instead of the usual correction factor coming from the soft and the virtual photon contributions, as in eq. (31). In this sense, since the soft and collinear radiation factorizes with respect to the Born helicity changing decay rate, the double logarithm Sudakov term can be equally factorized. It could be useful to investigate the impact of the higher order terms on the radiative correction to the Born amplitude. It is an open question if say hard collinear photons can be factorized and resummed. The particular mass cancellation mechanism occurring in the right-handed radiative decay is the consequence of the combination of two constraints: the angular momentum conservation in the pion vertex and the helicity flip in the photon-lepton vertex. The situation is completely different if we consider the radiative process with the outgoing left-handed lepton, i.e. the process without helicity flip. The differential Inner Bremsstrahlung contribution for left-handed outgoing lepton is given by: $`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{\mathrm{d}\mathrm{\Gamma }_{IB}^L}{\mathrm{d}y}}={\displaystyle \frac{\alpha }{4\pi }}{\displaystyle \frac{1}{(1r)^2}}{\displaystyle \frac{1}{A(1y+r)}}\{2A[r(2Ar+6)y^212A]`$ $`+\left[(A2r)(1+r)^2+Ay(y4r)2ry(1+y)+y(5r^2+y^2+1)\right]\mathrm{ln}{\displaystyle \frac{y+A}{yA}}`$ $`+(1y+r)^2(2rAy)\mathrm{ln}{\displaystyle \frac{y+A2}{yA2}}\}`$ (36) The expression (36) contains collinear singularities, since the coefficients of the collinear logarithms don’t vanish in the limit $`r0`$, as one can see from eq. (36). The decay rate (36) is also infrared divergent, as one can verify by taking the limit $`yy^{MAX}`$. In this case we don’t have the constraint constituted by the helicity flip in the photon-lepton vertex and in the pion vertex the angular momentum is conserved for soft and virtual photon emission. Thus, in the left-handed case, the mass singularity cancellation occurs in the ordinary way, as in the unpolarized case, i.e. in the total inclusive decay rate, obtained by adding all the order $`\alpha `$ contributions. Let us show how the cancellation takes place. As we have already seen (eq. (10)), in the Born $`\pi ^+`$ decay the outgoing lepton is left-handed, due to angular momentum conservation. Thus the unpolarized and the left-handed Born decay rate coincide: $$\mathrm{\Gamma }_0^L=\mathrm{\Gamma }_0.$$ (37) Because of the factorization with respect to the Born decay rate, also the unpolarized and the left-handed virtual contributions are equal: $$\mathrm{\Gamma }_v^L=\mathrm{\Gamma }_v.$$ (38) Expressing the left-handed Inner Bremsstrahlung contribution in terms of the unpolarized and the right-handed ones, the total contribution to order $`\alpha `$ to the left-handed process is given by: $$\mathrm{\Gamma }_{TOT}^L=(\mathrm{\Gamma }_0+\mathrm{\Gamma }_v+\mathrm{\Gamma }_{IB})\mathrm{\Gamma }_{IB}^R.$$ (39) The expression (39) is finite both in the infrared and in the collinear limit, because the mass singularities present in the terms between brackets cancel each other, as we have seen (see eq. (33) and $`\mathrm{\Gamma }_{IB}^R`$ is free from mass singularities. Let us now discuss the origin of these different cancellation mechanisms. The presence of mass singularities is a consequence of the fact that the states of a theory containing massless particles are highly degenerate. The KLN theorem states that the mass singularities disappear from the transition probability when we average it over the ensemble of degenerate states. This theorem contains the Block-Nordsieck theorem as a special case, when we consider only the cancellation of the infrared divergences. We can define two degeneration ensembles, one relative to the infrared divergences and one relative to the collinear singularities. We call them the infrared and the collinear ensembles, respectively. If we sum the transition probability over the states contained in the former, the infrared divergences cancel out and if we average it over the latter, we obtain a quantity free from both infrared and collinear singularities. The infrared ensemble is the one prescribed by the Block-Nordsieck theorem, while the second is the one prescribed by the KLN theorem and contains the first as a subset. Let us now examine how the infrared and collinear ensemble are composed for the radiative pion decay in the cases of unpolarized, left-handed and right-handed outgoing lepton. This discussion concerns the issue of the degeneration of states already addressed for the unpolarized case . This issue in the case of helicity changing processes presents peculiar features. We have seen in section 4.1 that in the unpolarized and left-handed Inner Bremsstrahlung contributions there are mass singularities, indicating that we have not summed the transition probability over the entire ensemble of degenerate states. In these cases the infrared ensemble contains, to order $`\alpha `$, the state with a pion, a charged lepton and a neutrino and all the other states differing from this for the presence of a soft virtual or real photon, i.e. a photon with an energy $`E_\gamma <\omega `$, where $`\omega `$ is an infrared cut off, tipically the measurement apparatus resolution. The collinear ensemble is constituted by all the states of the infrared ensemble plus the states with a hard photon moving parallel to the pion or the lepton. Clearly the degeneration arises in the limit $`m_l0`$. According to the KLN theorem the fact that $`\mathrm{d}\mathrm{\Gamma }_{IB}^R/\mathrm{d}y`$ is finite both in the infrared and in the collinear limits means that in the right-handed case, calculating the differential decay rate (i.e. summing over the photon polarization and integrating over the photon energy and emission angle), we have already averaged over the set of degenerate states relative to this process. Let us now consider how the degeneration ensemble for the right-handed radiative decay is composed. To obtain $`\mathrm{d}\mathrm{\Gamma }_{IB}^R/\mathrm{d}y`$ we have not averaged over the infrared subspace of the collinear ensemble, but this is enough to render the transition probability free from mass singularities. Indeed in this case the infrared ensemble is empty, owing to the constraints imposed both by the angular momentum conservation in the pion vertex and by the helicity flip in the photon-lepton vertex. Thus in this case the degeneration ensemble contains only the states with the outgoing lepton accompanied by hard collinear photons. We conclude that imposing to the outgoing lepton a polarization opposed to the one prescribed by the vertex preceding the photon emission, implies a reduction of the degeneration subspace. This fact has two consequences: the first is that both the infrared and the collinear limits are finite and the second is that these limits are disconnected, since the collinear degeneration subspace is constituted only by the states with the pion and the outgoing lepton accompanied by hard collinear photons. Thus in this case we have a particular application of the KLN theorem. ### 4.2 The $`Z^0`$ case In the radiative pion decay, due to the angular momentum conservation in the pion vertex, there is no room for a right-handed lepton. For such a channel, soft and virtual photon contributions are zero. This result is valid independently of the lepton mass. Let us now consider a more general case, by loosing the value of the angular momentum of the decaying state. As an example, we study the radiative $`Z^0`$ decay in a lepton-antilepton $`(l^{}l^+)`$ pair, in which the lepton is in a definite helicity state. The $`Z^0`$-leptons vertex is: $$i\frac{M_Z}{\sqrt{2}}\left(\frac{G}{\sqrt{2}}\right)^{1/2}\gamma _\mu (g_vg_a\gamma _5),$$ with $$g_v=14\mathrm{sin}\theta _{W}^{}{}_{}{}^{2}g_a=1$$ where $`\theta _W`$ is the Weinberg angle and $`M_Z`$ is the $`Z^0`$ mass. We have chosen this process, since, by varying the constants $`g_v`$ and $`g_a`$, we can control the structure of the $`Z^0`$-leptons coupling; thus it is possible to point out the role played by the conservation law occurring in this vertex in the chiral limit in the collinear singularity cancellation. If we set $`g_v=g_a=1`$, we require that in the limit of zero lepton mass, the $`Z^0`$ couples to a left-handed lepton. We calculate the decay rate for the process in which the lepton is right-handed. At the Born level this is given by: $`\mathrm{\Gamma }_0^R`$ $`=`$ $`{\displaystyle \frac{GM_Z^3}{48\sqrt{2}\pi }}\left\{\sqrt{14r}\left[(g_v^2+g_a^2)(1r)+3r(g_v^2g_a^2)\right]2g_vg_a(14r)\right\}.`$ (40) If in eq. (40) we set $`g_v=g_a=1`$, $`\mathrm{\Gamma }_0^R`$ vanishes in the chiral limit, since there isn’t the term related to the axial anomaly. Let us now study the decay process with the lepton emitting a real photon (see fig. 6) and evaluate the probability that the outgoing lepton is right-handed. The electromagnetic interaction doesn’t couples states with different chirality, hence the decay rate is expected to vanish for $`r0`$. The decay rate for the process described by the diagram of fig. 6, differential with respect to the lepton energy, is given by: $`{\displaystyle \frac{\mathrm{d}\mathrm{\Gamma }_\gamma ^R}{\mathrm{d}y}}={\displaystyle \frac{\alpha GM_Z^3}{96\sqrt{2}\pi ^2}}\{{\displaystyle \frac{(y2)(1y)^2}{4(1y+r)^2}}[(g_v^2+g_a^2)A+2g_vg_a(2ry)]`$ $`+`$ $`{\displaystyle \frac{(1y)}{2(1y+r)}}\left[(g_v^2+g_a^2)A(2ry)+2g_vg_a(y^22r)\right]`$ $`+`$ $`{\displaystyle \frac{2}{(y1)}}\left[2(g_v^22g_a^2)Ar+(g_v^2+g_a^2)A+2g_vg_a(4ry^2+y1)\right]`$ $`+`$ $`[g_v^2+g_a^2+2g_vg_a{\displaystyle \frac{1}{A}}(4r^2+r(yy^2+1)y)][\mathrm{ln}{\displaystyle \frac{y+A}{yA}}\mathrm{ln}{\displaystyle \frac{y+A2}{yA2}}]\}.`$ Here $$r=\frac{m_l^2}{M_Z^2},$$ $`y`$ is the usual dimensionless variable: $$y=\frac{2E_1}{M_Z}$$ where $`E_1`$ is the lepton energy and $$A=\sqrt{y^24r}.$$ The physical region for $`y`$ is $$2\sqrt{r}y1.$$ (42) The result obtained, as given by the emission of the photon by a single leg, is gauge dependent. To have a gauge independent amplitude, the contribution of the diagram b) of fig. 7 must be added. For the purpose of the polarized amplitude, however, the helicity flip contribution of the diagram b) of fig. 7 gives zero in the massless limit and is therefore negligible in our discussion. From now on we consider the case $`g_v=g_a=1`$, to have the condition of chirality conservation in the $`Z^0`$ vertex for $`m_l0`$. Taking this limit in eq. (LABEL:Rz), we see that $`\mathrm{d}\mathrm{\Gamma }_\gamma ^R/\mathrm{d}y`$ does not vanish: $$\frac{\mathrm{d}\mathrm{\Gamma }_\gamma ^R}{\mathrm{d}y}\frac{\alpha }{2\pi }(1y)\mathrm{\Gamma }_0(Z^0\nu \overline{\nu })\mathrm{as}r0,g_v=g_a=1.$$ (43) The result of this limit is the contribution related to the axial anomaly, which has the same form of the one found in the pion case. Let us now discuss the mass singularities cancellation mechanism for the $`Z^0`$ decay case. If we keep the lepton mass different from zero, the helicity is not fixed by the interaction occurring before the photon emission, even if we set $`g_v=g_a=1`$. Thus, for $`m_l0`$, the soft and virtual photons contribution are different from zero and diverge in the infrared limit. Indeed, if we let the lepton energy reach its kinematical limit, $`y^{MAX}=1`$, we see that $`\mathrm{d}\mathrm{\Gamma }_\gamma ^R/\mathrm{d}y`$ diverges. We expect the infrared divergences to cancel, if we add all the first order contributions, given by the diagrams of fig. 6 and 7 and calculate the totally inclusive decay rate. Eq. (43) shows that the collinear limit gives a finite result. Thus, we conclude that, as in the case of the radiative pion decay, evaluating $`\mathrm{d}\mathrm{\Gamma }_\gamma ^R/\mathrm{d}y`$ we have already summed over all the collinear degeneration subspace. As a consequence, the collinear and infrared limit are disconnected. If we take the limit $`y1`$ in eq. (43), we obtain: $$\frac{\mathrm{d}\mathrm{\Gamma }_\gamma ^R}{\mathrm{d}y}0\mathrm{as}m_l0,y1\mathrm{and}g_v=g_a=1.$$ (44) In general, if a quantity is finite in the collinear limit, it is finite also in the infrared limit, since the collinear subspace contains the infrared one. The eq. (44) indicates that in the massless limit the soft photon contribution is zero. Indeed, it is factorized with respect to the Born decay rate, which, for $`r0`$ and $`g_v=g_a=1`$, vanishes. As we have discussed in section 2, the presence of the anomalous term is directly connected to the emission of the photon, hence it vanishes in the infrared limit. The virtual photon contribution vanishes in the zero mass limit, as well. Indeed, it is factorized respect to the Born decay rate, which goes to zero as $`r0`$. The virtual correction factor can produce only a logarithmic collinear singularity, not a power-like one, needed for the cancellation of the chiral suppression. We observe that taking the infrared limit $`x0`$ in eq. (22), gives a finite result (indeed the cross section vanishes). This is a consequence of the fact that the cross section (22) has been calculated to the leading order in the lepton mass. From the eq. (LABEL:Rz), we see that, for $`r1`$, the infrared divergent term is given by: $$\left(\frac{\mathrm{d}\mathrm{\Gamma }_\gamma ^R}{\mathrm{d}y}\right)_{IR}\frac{4r}{(y1)}\left(\frac{2r}{y}y\frac{2}{y}\right)$$ (45) and it is proportional to the lepton mass. Performing the calculation, neglecting the mass terms, as done in ref. , means imposing the chirality conservation law in the $`Z^0`$ vertex; thus the soft photon contribution is zero and the infrared divergences disappear. To the leading order in the lepton mass, we have only the anomalous term, which vanishes in the infrared limit. As done for the pion case, we now examine the composition of the degeneration ensembles for the radiative $`Z^0`$ decay with the right-handed lepton. Eq. (44) indicates that, in the chiral limit, the soft photons don’t contribute to the process with the right-handed outgoing lepton, since, not carrying spin, they cannot contribute to the helicity flip. As we have already said, the virtual photon contribution vanishes in the zero mass limit. In the limit $`m_l0`$, also the diagram with the photon emitted by $`l^+`$ doesn’t contribute, since, clearly, it violates the chirality conservation in the $`Z^0`$ vertex. The collinear degeneration arises in the massless limit. The result (43) shows that, in this limit, the collinear ensemble is constituted only by the states with the lepton accompanied by hard collinear photons, just as in the case of the pion decay. Thus for $`m_l0`$ the infrared ensemble is empty and also the states with the antilepton accompanied by a hard collinear photon don’t contribute. ## 5 Conclusions We have shown that the Dolgov and Zakharov treatment of the axial anomaly can be extended to processes characterized by a lepton which changes helicity by emitting a photon, as it was already noticed in . The corresponding decay rates don’t vanish in the chiral limit, due to a term independent of the lepton mass; we interpret the presence of this term as related to the axial anomaly. This can be seen as a signal of the anomalous symmetry breaking in processes different from the usual ones, like the $`\pi ^0\gamma \gamma `$ decay. We have computed the rates corresponding to the $`\pi ^+`$ and $`Z^0`$ radiative decays. We have analysed their infrared and collinear limits. It results essential to keep the terms of all orders in the lepton mass, since the cancellation of the infrared and collinear divergences takes place among these terms. We have examined the connection between the polarization of the outgoing leptons and the application of the KLN theorem. We have shown that for the helicity changing processes the cancellation of the collinear singularities occurs through a mechanism different from the usual one of real and virtual compensation. The coefficients in front of the collinear terms go to zero in the chiral limit, producing the finiteness of the distribution. We have found, however, a difference between the pion case and the more general case of the $`Z^0`$ decay. The former represents a particular case due to the angular momentum conservation in the pion vertex. As a consequence, the virtual and real soft photon contribution are zero, even if the lepton mass is kept different from zero. The Inner Bremsstrahlung contribution is finite both in the infrared and in the collinear limits. In the $`Z^0`$ case, the decay rate diverges in the infrared limit, since, for $`m_l0`$, the soft photon contributions are not zero. However in the collinear limit, the result is finite, despite the fact that the collinear degenerate states arise only in the zero lepton mass limit. In this limit the virtual and real soft photon contributions do vanish. In order to make the collinear limit finite, it is therefore sufficient to sum over degenerate states made of the changing helicity lepton accompanied by a hard collinear photon. The transition probability becomes finite after summing over the photon final phase space. We noticed that in the helicity changing processes the collinear limit results disconnected from the infrared one. The contributions coming from the virtual photon emission and from the emission of photons by particles different from the one changing helicity, are zero. This situation is due to the fact that the Born part of the process fixes the fermion chirality in the zero mass limit, while, after the photon emission, it is in a state of opposite chirality; this reduces the collinear ensemble. We conclude that the collinear singularity cancellation mechanism for helicity changing processes is controlled by the anomalous breaking of the chiral symmetry. The axial anomaly implies that the collinear limit gives a finite result, independent of the fermion mass. The extension of the above remarks to other gauge theories like QCD, is possible. It could allow a more systematic and complete treatment of the infrared and collinear singularities. Acknowledgments We wish to thank V. Fadin and E. Kuraev for valuable comments and S. Forte, J. Kodaira and L. Lipatov for useful discussions.
warning/0006/hep-ph0006102.html	ar5iv	text	# LIGHT VECTOR MESON PHOTOPRODUCTION AT HIGH \|𝑡\| ## 1 Introduction My talk is based on our recent publication $`^\mathrm{?}`$. We derived pQCD predictions for helicity amplitudes of the high $`\|t\|=𝐪^2`$ diffractive light vector meson photoproduction process. Initial proton disintegates into hadron system $`X`$ which is separated from the produced vector meson $`V`$ by the rapidity gap $`\eta _0`$. The cross section for this reaction can be related to those for the photoproduction of $`V`$ off a quark and a gluon via the gluon and quark densities in a proton $`G(x,t)`$ and $`q(x,t)`$: $`{\displaystyle \frac{d^2\sigma (\gamma pVX)}{dtdx}}`$ $`=`$ $`{\displaystyle \underset{f}{}}\left(q(x,t)+\overline{q}(x,t)\right){\displaystyle \frac{d\sigma (\gamma qVq)}{dt}}+`$ (1) $`G(x,t){\displaystyle \frac{d\sigma (\gamma GVG)}{dt}};x={\displaystyle \frac{4\|t\|}{s}}\mathrm{cosh}^2{\displaystyle \frac{\eta _0}{2}}.`$ There are three independent helicity amplitudes for the parton subprocess $`\gamma qVq`$, the first and the second indices are the helicities of photon and vector meson respectively $$M_{++}(M_{}=M_{++}),M_{+\mathrm{\hspace{0.33em}0}}(M_{\mathrm{\hspace{0.33em}0}}=M_{+\mathrm{\hspace{0.33em}0}}),M_+(M_+=M_+).$$ Within usual perturbation theory a photon can split only into a chiral-even (the helicities of the quark and the antiquark are antiparallel) massless quark pair. The violation of chiral symmetry, which is well known to be a soft QCD phenomenon, generates a nonperturbative chiral-odd component of the real photon wave function. The interaction of this additional chiral contribution can however be described in pQCD since high $`t`$ quark-dipole scattering chooses a $`q\overline{q}`$ configuration with small transverse interquark distances. As a result the chiral-odd contribution can be factorized into a convolution of two nonperturbative photon and vector meson light-cone wave functions with the hard scattering amplitude. The chiral-odd wave function of a real photon have a similar form as the chiral-odd wave function of a vector meson. The photon dimentional coupling constant $`f_\gamma `$ (wich is similar to vector meson chiral-odd constant $`f_V^T`$) is a product of the quark condensate $`<\overline{q}q>`$ and its magnetic susceptibility $`^\mathrm{?}`$, $$f_\gamma =<\overline{q}q>\chi 70MeV.$$ (2) Both parameters $`<\overline{q}q>`$ and $`\chi `$ describing QCD vacuum have been tested in various QCD sum rule applications. ## 2 Results and Discussion The helicity amplitudes of the process $`\gamma qVq`$ are the sums of the chiral-even and the chiral-odd contributions $$M_{\lambda _1\lambda _2}=M_{\lambda _1\lambda _2}^{even}+M_{\lambda _1\lambda _2}^{odd}.$$ At asymptotically high $`\|t\|`$ the dominant helicity amplitude is $`M_{+\mathrm{\hspace{0.33em}0}}`$. Its chiral-even part $`M_{+\mathrm{\hspace{0.33em}0}}^{even}`$ has the minimal, $`1/𝐪^3`$, suppression. $$M_{+\mathrm{\hspace{0.17em}0}}^{even}=is\alpha _s^2\frac{32\pi }{3\sqrt{2}}eQ_Vf_V/q^3;.$$ (3) We present our results for helicity amplitudes as the ratios $`M_{\lambda _1\lambda _2}/M_{+\mathrm{\hspace{0.33em}0}}^{even}`$ $$\frac{M_{+\mathrm{\hspace{0.33em}0}}}{M_{+\mathrm{\hspace{0.33em}0}}^{even}}=1\frac{24\pi ^2f_\gamma f_V^Tm_V}{3f_V𝐪^2}\left(\mathrm{ln}\frac{1u_{min}}{u_{min}}2(12u_{min})\right)$$ (4) $$\frac{M_{++}}{M_{+\mathrm{\hspace{0.33em}0}}^{even}}=\frac{m_V}{\|𝐪\|\sqrt{2}}\left(\mathrm{ln}\frac{1u_{min}}{u_{min}}1+2u_{min}\right)+\frac{24\pi ^2f_\gamma f_V^T}{3f_V\|𝐪\|\sqrt{2}}$$ (5) $$\frac{M_+}{M_{+\mathrm{\hspace{0.33em}0}}^{even}}=\frac{3m_V}{\sqrt{2}\|𝐪\|}\frac{48\pi ^2f_\gamma f_V^Tm_V^2}{3f_V\|𝐪\|^3\sqrt{2}}\left(2\mathrm{ln}\frac{1u_{min}}{u_{min}}3(12u_{min})\right)$$ (6) The chiral-odd contributions to the amplitudes are accompanied with large numerical coefficients. In the case of $`M_{++}`$ the chiral-even and the chiral-odd parts of (5) add with the same signs. In contrast, the chiral-even and the chiral-odd parts of $`M_{+\mathrm{\hspace{0.33em}0}}`$ or $`M_+`$ enters with opposite sign which leads to an effective reduction of these amplitudes for intermediate $`\|t\|`$. Though the chiral-even and the chiral-odd contributions to $`M_{++}`$ are of the same order with respect to $`1/\|𝐪\|`$ counting the chiral-odd one can be dominant up to very large $`\|t\|`$. According to (5) for the $`t`$ range $`3÷\mathrm{\hspace{0.33em}8}\text{ GeV}^2`$ the chiral-even part constitutes only $`10÷\mathrm{\hspace{0.33em}20}\%`$ of the $`M_{++}`$ amplitude and for $`\|t\|100\text{GeV}^2`$ $`M_{++}^{even}0.72M_{++}^{odd}`$. Due to a large compensation between chiral-even and chiral-odd parts of $`M_{+\mathrm{\hspace{0.33em}0}}`$ and $`M_+`$ the non spin-flip $`M_{++}`$ amplitude dominates strongly in the region of intermediately high $`\|t\|`$. $`M_{+\mathrm{\hspace{0.33em}0}}`$ will exceed $`M_{++}`$ only at $`\|t\|>40(\text{ GeV})^2`$. For the $`\|t\|`$ interval $`3\text{ GeV}^2÷\mathrm{\hspace{0.33em}8}\text{ GeV}^2`$ $$\frac{M_{+\mathrm{\hspace{0.33em}0}}}{M_{++}}0.25÷\mathrm{\hspace{0.33em}0.35},\frac{M_+}{M_{++}}\mathrm{\hspace{0.17em}0.02}÷\mathrm{\hspace{0.33em}\hspace{0.17em}0.04}$$ (7) Both chiral-even and chiral–odd parts of the amplitudes were calculated in the leading order of $`1/\|𝐪\|`$ expansion. As usual in the QCD approach to any exclusive reaction, the question about the region of applicability of these results is open untill the power corrections have not studied. In our case the situation is more difficult because the factorization of the amplitude into hard and soft parts is violated for the chiral-even part of $`M_{++}`$ and the chiral-odd parts of $`M_{+\mathrm{\hspace{0.33em}0}}`$ and $`M_+`$. In these cases the corresponding integrals over the quark longitudinal momentum $`u`$ contain the end point logarithmic singularities. At present we simply restrict the corresponding $`u`$ integrals to the interval $`[1u_{min},u_{min}],u_{min}=m_V^2/𝐪^2`$, which corresponds to the contribution of the hard region only. The good news is, however, that the chiral-odd part $`M_{++}^{odd}`$ of the dominant in the intermediatly large $`\|t\|`$ region helicity amplitude $`M_{++}`$ is free from this end-point singularity. This chiral-odd part is numerically considerably larger than the hard contribution to its chiral-even counterpart $`M_{++}^{even}`$. Therefore we can expect that the relative uncertainty related with the uncalculated soft contributions to $`M_{++}^{even}`$ is small. ## Acknowledgments This work was supported by the Foundation Universities of Russia (grant 015.02.01.16). ## References
warning/0006/cond-mat0006079.html	ar5iv	text	# Spectral function of transverse spin fluctuations in an antiferromagnet ## I Introduction The antiferromagnetic state of the half-filled Hubbard model with nearest-neighbour (NN) hopping is characterized by an energy gap, and the spectrum of transverse spin fluctuations consists of the low-lying, collective (magnon) excitations, as well as the single-particle excitations across the gap. Generally, there is a clear distinction between these two excitations which are well separated in energy. However, recent extensions, e.g. with disorder, impurities, and next-nearest-neighbour (NNN) hopping, clearly show the existence of essentially gapless antiferromagnetism, arising from a variety of mechanisms. Thus, with increasing disorder (on-site potential disorder) the two Hubbard bands progressively broaden, until the band gap vanishes when the disorder strength $`WU`$. With low-$`U`$ impurities on the other hand, the effective charge gap becomes negligible due to nearly localized states on the impurities. And lastly, when NNN hopping is included, a band asymmetry is introduced which reduces the band gap in the AF state, and in weak coupling there exists a region of gapless AF phase in the magnetic phase diagram. This gaplessness implies that the collective magnon excitations are no longer well defined, and actually merge with the single-particle excitations, thereby necessitating a unified scheme for the evaluation of transverse spin fluctuation. While the magnon contribution to the transverse spin fluctuations was evaluated recently, and the sublattice magnetization and the Néel temperature were obtained within a renormalized spin fluctuation theory in the whole $`U/t`$ range, in this paper we describe an alternate scheme for evaluating the transverse spin fluctuations which allows both the collective excitations and the single-particle excitations to be studied on the same footing. Single-particle excitations are especially significant because it is precisely this part of the transverse spin spectral function which allows for a quantitative distinction of the antiferromagnetic state within the Hubbard model from that of an equivalent Heisenberg spin model with $`U`$-dependent, extended-range spin couplings $`J_{ij}(U)`$, but possessing only magnon excitations. A study of the relative strengths of the magnon and single-particle excitations, in terms of the integrated spectral weights in the whole $`U/t`$ range, will therefore allow for a quantitative demarkation along the $`U/t`$ axis below which single-particle excitations are the dominant contribution. Thus, while in this low-$`U`$ regime the AF state is not well described in terms of an effective Heisenberg spin model possessing only magnon excitations, in the intermediate and strong coupling limits use of an effective Heisenberg model with $`U`$-dependent spin couplings $`J_{ij}(U)`$, and generally use of a spin picture is appropriate. Yet another significance of the single-particle excitations is interestingly connected with the spin commutation relation $`[S^+,S^{}]=2S^z`$. The RPA-level ground-state expectation value $`[S^+,S^{}]_{\mathrm{RPA}}`$, involving the difference of transverse spin correlations evaluated in the AF ground state, should be identically equal to $`2S^z_{\mathrm{HF}}`$, the local magnetization at the HF level. This is deduced from the fact that both RPA and HF approximations are of the same order (O(1)) within the systematic inverse-degeneracy expansion scheme in powers of $`1/𝒩`$, where $`𝒩`$ is the number of orbitals per site. Therefore both RPA and HF become exact in the limit $`𝒩\mathrm{}`$, when all corrections of order $`1/𝒩`$ or higher vanish. The magnon contribution $`[S^+,S^{}]_{\mathrm{RPA}}^{\mathrm{magnon}}`$ was indeed found to be in excellent agreement with the HF magnetization $`2S^z_{\mathrm{HF}}`$ in the intermediate and strong coupling limits. However, a discrepancy was observed at small $`U`$ which was attributed to the neglect of the single-particle excitations, which become relatively more important in the weak coupling limit. We will show here that indeed when the single-particle excitations are included this discrepancy is exactly removed. Finally, the unified scheme for incorporating both the single particle and collective excitations in the evaluation yields the complete spectrum of transverse spin excitations in the magnetic state. Calculations with realistic models of magnetic solids can hence be used for comparison with results of scattering studies with neutrons, which are direct experimental probes into the spectrum of magnetic excitations in solids. In this regard the necessity of more realistic models which include NNN hopping etc. has been acknowledged recently from realistic band structure studies, photoemission data and neutron-scattering measurements of high-T<sub>c</sub> and related materials. ## II Hubbard model with next-nearest-neighbour hopping We consider the following Hamiltonian on a square lattice, with hopping terms $`t`$ and $`t^{}`$ between nearest-neighbour (NN) and next-nearest-neighbour (NNN) pairs of sites $`ij`$ and $`ik`$ respectively, $$H=t\underset{ij\sigma }{\overset{\mathrm{NN}}{}}a_{i\sigma }^{}a_{j\sigma }t^{}\underset{ik\sigma }{\overset{\mathrm{NNN}}{}}a_{i\sigma }^{}a_{k\sigma }+U\underset{i}{}n_in_i.$$ (1) Extension to the simple cubic lattice is considered in the Appendix. In the plane-wave basis defined by $`a_{i\sigma }=\sqrt{\frac{1}{N}}_𝐤e^{i𝐤.𝐫_i}a_{𝐤\sigma }`$, the non-interacting part of the Hamiltonian $`H_0=_{𝐤\sigma }(ϵ_𝐤+ϵ_𝐤^{})a_{𝐤\sigma }^{}a_{𝐤\sigma }`$, where $`ϵ_𝐤`$ and $`ϵ_𝐤^{}`$ are the two free-particle energies, correponding to NN and NNN hopping respectively, $`ϵ_𝐤`$ $`=`$ $`2t(\mathrm{cos}k_x+\mathrm{cos}k_y)`$ (2) $`ϵ_𝐤^{}`$ $`=`$ $`4t^{}\mathrm{cos}k_x.\mathrm{cos}k_y.`$ (3) For the NN model, the HF-level description of the broken-symmetry AF state, and the transverse spin fluctuations have been discussed earlier in the strong, intermediate, and weak coupling limits. We briefly discuss the extension for the NNN hopping. Since the NNN hopping term connects sites in the same sublattice, in the two-sublattice basis the $`ϵ_𝐤^{}`$ term appears in the diagonal matrix elements of the HF Hamiltonian $$H_{\mathrm{HF}}^\sigma (𝐤)=\left[\begin{array}{cc}\sigma \mathrm{\Delta }+ϵ_𝐤^{}& ϵ_𝐤\\ ϵ_𝐤& \sigma \mathrm{\Delta }+ϵ_𝐤^{}\end{array}\right]=ϵ_𝐤^{}\mathbf{\hspace{0.33em}1}+\left[\begin{array}{cc}\sigma \mathrm{\Delta }& ϵ_𝐤\\ ϵ_𝐤& \sigma \mathrm{\Delta }\end{array}\right]$$ (4) for spin $`\sigma `$. Here $`2\mathrm{\Delta }=mU`$, where $`m`$ is the sublattice magnetization. For the NN hopping model $`2\mathrm{\Delta }`$ is also the energy gap for single-particle excitations. Since the $`ϵ_𝐤^{}`$ term appears as a unit matrix, the eigenvectors of the HF Hamiltonian remain unchanged from the NN case, whereas the eigenvalues correponding to the quasiparticle energies are modified to, $$E_{𝐤\sigma }^{(\pm )}=ϵ_𝐤^{}\pm \sqrt{\mathrm{\Delta }^2+ϵ_𝐤^2}.$$ (5) The two signs $`\pm `$ refer to the two quasiparticle bands. The band gap is thus affected by the NNN hopping term, and it progressively decreases as $`2\mathrm{\Delta }4t^{}`$ in the weak coupling limit. As the eigenvectors of the HF Hamiltonian are unchanged, the self-consistency condition retains its form, and therefore the sublattice magnetization is independent of $`t^{}`$, provided there is a nonzero band gap, with the lower band occupied and the upper band empty. We will restrict ourselves only to this regime where there is no band overlap, though the gap may vanish when the two bands are just touching. The fermionic quasiparticle amplitudes $`a_{𝐤\sigma }`$ and $`b_{𝐤\sigma }`$ for spin $`\sigma =,`$ and the two quasiparticle bands are given by $`a_𝐤^2=b_𝐤^2=a_𝐤^2=b_𝐤^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{\Delta }}{\sqrt{\mathrm{\Delta }^2+ϵ_𝐤^2}}}\right)`$ (6) $`a_𝐤^2=b_𝐤^2=a_𝐤^2=b_𝐤^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\mathrm{\Delta }}{\sqrt{\mathrm{\Delta }^2+ϵ_𝐤^2}}}\right).`$ (7) These relationships follow from the spin-sublattice and particle-hole symmetry in the AF state of the half-filled system. However, this situation changes when with increasing NNN hopping $`t^{}`$ the two bands overlap and the band gap vanishes. This overlap indicates that the low-lying states in the upper band corresponding to double occupancy are lower in energy than the high-energy states in the lower band. This results in charge transfer and consequently partially empty and doubly occupied sites. The situation is therefore analogous to the addition of holes or electrons in the half-filled Hubbard model, with all its associated complications of spin bags, strings of upturned spins, spiral and striped phases etc. Thus the simple gapless AF state of the NNN hopping model, which is indeed a self-consistent HF solution, may actually be unstable due to same instabilities as the hole-doped AF. Numerical HF studies do show this instability, and will be further discussed elsewhere. ## III The transverse spin spectral function The spectral function for transverse spin fluctuations in the AF state is obtained from the imaginary part of the corresponding time-ordered propagator of the transverse spin operators $`S_i^{}`$ and $`S_j^+`$ at sites $`i`$ and $`j`$, $`\chi ^+(𝐪\omega )=𝑑t_ie^{i\omega (tt^{})}e^{i𝐪.(𝐫_i𝐫_j)}\mathrm{\Psi }_\mathrm{G}\|T[S_i^{}(t)S_j^+(t^{})]\|\mathrm{\Psi }_\mathrm{G}`$. In terms of the RPA expression $`[\chi ^+(𝐪\omega )]_{\mathrm{RPA}}=[\chi ^0(𝐪\omega )]/[1U\chi ^0(𝐪\omega )]`$, where $`\chi ^0(𝐪\omega )`$ is the zeroth-order particle-hole propagator, there are two sources of contribution to the spectral function, $`A(\omega )=_𝐪\mathrm{Im}\mathrm{Tr}[\chi ^+(𝐪\omega )]`$. The contribution arising from the imaginary part of $`\chi ^0(𝐪\omega )`$ is associated with the single-particle excitations across the gap (with $`\omega >2\mathrm{\Delta }`$ for the NN model), whereas that due to the vanishing of the denominator $`1U\chi ^0(𝐪\omega )=0`$ is associated with the collective magnon excitations (involving $`\omega <2\mathrm{\Delta }`$ for the NN model). As mentioned earlier, this latter contribution to the integrated spectral weight $`\pi ^1𝑑\omega A(\omega )`$, and therefore to the transverse spin correlations $`S^+S^{}_{\mathrm{RPA}}`$ and $`S^{}S^+_{\mathrm{RPA}}`$ was recently studied in the context of quantum correction to sublattice magnetization and the Néel temperature in the whole $`U/t`$ range. The evaluation of the transverse spin spectral function, including contributions from both the single-particle and collective excitations, is facilitated by expressing the $`2\times 2`$ complex matrix $`[\chi ^0(𝐪\omega )]`$ in terms of its eigenvalues $`\lambda _𝐪^n(\omega )`$ and eigenvectors $`\|\varphi _𝐪^n(\omega )`$, and we obtain $`A(\omega )`$ $`=`$ $`{\displaystyle \underset{𝐪}{}}\mathrm{ImTr}[\chi ^+(𝐪\omega )]={\displaystyle \underset{𝐪}{}}\mathrm{ImTr}{\displaystyle \frac{[\chi ^0(𝐪\omega )]}{\mathrm{𝟏}U[\chi ^0(𝐪\omega )]}}`$ (8) $`=`$ $`{\displaystyle \underset{𝐪}{}}\mathrm{ImTr}{\displaystyle \underset{n=1,2}{}}\left({\displaystyle \frac{\lambda _𝐪^n(\omega )}{1U\lambda _𝐪^n(\omega )}}\right)\|\varphi _𝐪^n(\omega )\varphi _𝐪^n(\omega )\|`$ (9) $`=`$ $`{\displaystyle \underset{𝐪}{}}\mathrm{Im}{\displaystyle \underset{n=1,2}{}}\left({\displaystyle \frac{\lambda _𝐪^n(\omega )}{1U\lambda _𝐪^n(\omega )}}\right).`$ (10) If instead of taking the trace in Eq. (6), the two diagonal matrix elements are considered separately, then we obtain the transverse spin correlations on the A(B) sublattice by integrating over frequency: $`S^+S^{}_{\mathrm{B}(\mathrm{A})}=S^{}S^+_{\mathrm{A}(\mathrm{B})}={\displaystyle \frac{d\omega }{\pi }\underset{𝐪}{}\mathrm{Im}[\chi ^+(𝐪\omega )]_{\mathrm{A}(\mathrm{B})}}`$ (11) $`={\displaystyle \frac{d\omega }{\pi }\underset{𝐪}{}\underset{n=1,2}{}\mathrm{Im}\left(\frac{\lambda _𝐪^n(\omega )}{1U\lambda _𝐪^n(\omega )}\right)\|\varphi _𝐪^n(\omega )\|_{\mathrm{A}(\mathrm{B})}^2}.`$ (12) Here the relationship $`S^+S^{}_{\mathrm{B}(\mathrm{A})}=S^{}S^+_{\mathrm{A}(\mathrm{B})}`$ between the transverse spin correlations on opposite sublattices follows from the spin-sublattice symmetry in the AF state. ## IV Computational procedure and Results To begin with the zeroth order $`2\times 2`$ matrix $`[\chi ^0(𝐪\omega )]=i\frac{d\omega }{2\pi }_𝐤^{}[G^{}(𝐤^{}\omega ^{})][G^{}(𝐤^{}𝐪,\omega ^{}\omega )]`$, is evaluated in the broken-symmetry AF state. In terms of the fermionic quasiparticle amplitudes and energies, $`[\chi ^0(𝐪\omega )]`$ is given by, $`[\chi ^0(𝐪\omega )]`$ $`=`$ $`{\displaystyle \underset{𝐤}{}}\left[\begin{array}{cc}a_𝐤^2a_{𝐤𝐪}^2\hfill & \hfill a_𝐤b_𝐤a_{𝐤𝐪}b_{𝐤𝐪}\\ a_𝐤b_𝐤a_{𝐤𝐪}b_{𝐤𝐪}\hfill & \hfill b_𝐤^2b_{𝐤𝐪}^2\end{array}\right]{\displaystyle \frac{1}{E_{𝐤𝐪}^{}E_𝐤^{}+\omega i\eta }}`$ (15) $`+`$ $`{\displaystyle \underset{𝐤}{}}\left[\begin{array}{cc}a_𝐤^2a_{𝐤𝐪}^2\hfill & \hfill a_𝐤b_𝐤a_{𝐤𝐪}b_{𝐤𝐪}\\ a_𝐤b_𝐤a_{𝐤𝐪}b_{𝐤𝐪}\hfill & \hfill b_𝐤^2b_{𝐤𝐪}^2\end{array}\right]{\displaystyle \frac{1}{E_𝐤^{}E_{𝐤𝐪}^{}\omega i\eta }}.`$ (18) Analytical evaluation of $`[\chi ^0(𝐪\omega )]`$ in the strong and intermediate coupling limits has been discussed earlier for the NN hopping model, and is described in the Appendix for the NNN hopping model. For arbitrary $`U`$, the $`𝐤`$-sum is numerically performed to evaluate the complex matrix $`[\chi ^0(𝐪\omega )]`$, which is then diagonalized to obtain the two eigenvalues $`\lambda _𝐪^n`$ and eigenvectors $`\|\varphi _𝐪^n`$. The infinitesimal $`\eta `$ is appropriately chosen according to the fineness of the $`𝐤`$ and $`𝐪`$ grids. Typically, the grid sizes $`dk`$ and $`dq`$ were taken of order 0.05 and $`\eta 0.01`$. Summing over the whole range of $`𝐪`$ values between $`0`$ and $`\pi `$ in each dimension then yields the spectral function $`A(\omega )`$ from Eq. (6), and the transverse spin correlations from Eq. (7). The integrated spectral weight is obtained by numerically integrating the spectral function $`A(\omega )`$ over frequency. The collective and single-particle contributions can be evaluated separately by integrating over the $`\omega `$-regions $`\omega <2\mathrm{\Delta }`$ and $`\omega >2\mathrm{\Delta }`$ respectively (for the NN case). As the single-particle excitations have a continuum distribution in Eq. (6), the evaluation of this contribution is less computationally intensive. However, for collective excitations, which are a set of delta functions, it is necessary to have in the numerical integration procedure fine enough $`\omega `$\- and $`q`$-grids, so that a sufficiently large number of magnon modes are picked up in the $`\omega `$ and $`q`$ sums. The method described in ref. is more efficient for evaluating the collective-excitation contribution. Before discussing the results we first consider the two limiting cases. In the limit of vanishing interaction strength $`U0`$, $`\mathrm{\Delta }0`$, we have $`\chi ^+(𝐪\omega )=\chi ^0(𝐪\omega )`$, and therefore the magnon contribution vanishes. Also, as all the quasiparticle probabilities $`a^2=b^2=1/2`$ from Eq. (5), we have from above $`\mathrm{Im}\chi ^0(𝐪\omega )={\displaystyle \underset{𝐤}{}}\left[\begin{array}{cc}\frac{1}{4}\hfill & \hfill \frac{1}{4}\\ \frac{1}{4}\hfill & \hfill \frac{1}{4}\end{array}\right]\times \pi \times `$ (21) $`[\delta (E_{𝐤𝐪}^{}E_𝐤^{}+\omega )+\delta (E_𝐤^{}+E_{𝐤𝐪}^{}\omega )].`$ (22) In the non-interacting limit, the integrated spectral weight therefore yields $$\stackrel{\mathrm{Lim}}{U0}\frac{d\omega }{\pi }\underset{𝐪}{}\mathrm{ImTr}[\chi ^+(𝐪\omega )]=1/2.$$ (23) On the other hand, in the strong coupling limit $`U\mathrm{}`$, the quasiparticle bands are narrowed to infinitesimal width, so that the imaginary part of $`\chi ^0(𝐪\omega )`$ strengthens to delta functions at the gap frequencies $`\omega =\pm U`$. Therefore, for both possible values of $`\chi ^0(𝐪\omega )`$, zero or infinity, the single-particle contribution to the imaginary part of $`\chi ^+(𝐪\omega )`$ vanishes. Thus, with increasing interaction strength, the spectrum of transverse spin fluctuations changes from predominantly single-particle excitations in weak coupling to predominantly magnon excitations in the strong coupling limit. Figures 1 through 6 depict results for the NN model. Fig. 1 shows a typical transverse spin spectral function $`A(\omega )=_q\mathrm{Im}\mathrm{Tr}\chi ^+(𝐪\omega )`$, for a moderate correlation value, showing the distinct separation between the low-lying, collective excitations $`(\omega <2\mathrm{\Delta })`$ and the single-particle excitations $`(\omega >2\mathrm{\Delta })`$. Fig. 2 shows the single-particle contribution to the spectral function for moderate gap values $`2\mathrm{\Delta }=1,2,3,4`$. The rapid decrease in the spectral weight with increasing correlation shows that the single-particle excitations are strongly suppressed in the strong correlation limit. With decreasing correlation, the separation between the low-lying collective excitations and the single-particle excitations progressively decreases. Fig. 3 shows the near merging of the collective $`(\omega <2\mathrm{\Delta })`$ and the single-particle $`(\omega >2\mathrm{\Delta })`$ excitations for $`2\mathrm{\Delta }=1`$ ($`U=2.28`$). A comparison of the integrated spectral weights as a function of $`U`$ is shown in Fig. 4, for the single-particle and collective excitations. While the collective excitations are seen to sharply fall-off in the weak coupling limit, the single-particle excitations are likewise suppressed in the strong coupling limit. For $`U>2.5`$, the collective (magnon) excitations are dominant, and thus an effective spin description of the antiferromagnetic state is appropriate down to a surprisingly low $`U`$ value. We now examine the role of the single-particle excitations in the sum rule following from the spin commutation relation $`[S^+,S^{}]=2S^z`$. As discussed in the Introduction, the RPA-level, AF ground-state expectation value $`[S_i^+,S_i^{}]_{\mathrm{RPA}}`$, including both magnon and single-particle contributions, should be identically equal to $`2S_i^z_{\mathrm{HF}}`$. The single-particle contribution to this expectation value, obtained from the transverse spin correlations $`S^+S^{}`$ and $`S^{}S^+`$ evaluated from Eq. (7), are shown in Table I for several $`U`$ values. Also shown are the magnon contributions which were obtained earlier. The total of the magnon and single-particle contributions is indeed in close agreement with the HF magnetization $`2S_i^z_{\mathrm{HF}}`$. The sharp fall-off of the gapless magnon contribution in the weak coupling limit essentially leaves only the single-particle contribution. As these excitations have a minimum-energy threshold of $`2\mathrm{\Delta }`$, the spectrum of trans- verse spin fluctuations essentially acquires a pseudo gap at low energies. This is clearly seen in Fig. 5, showing a gap-like suppression in the density of excitations at low energies. Fig. 6 shows a comparison of the spectral functions for $`\mathrm{\Delta }=0.1`$ and $`\mathrm{\Delta }=1.0`$. The drastic suppression in the contribution of the low-energy, collective excitations in going from $`U=3.29`$ ($`\mathrm{\Delta }=1.0`$) to $`U=1.19`$ $`(\mathrm{\Delta }=0.1)`$ is quite remarkable. Fig. 7 shows the effects of the NNN hopping on the spectral function. With increasing NNN hopping $`t^{}`$, the magnon spectrum shifts towards lower energy, which reflects the magnon softening. There is also a substantial increase in the low-energy spectral function due to the reduction in the energy gap with NNN hopping. For $`U=3.29`$, we have $`\mathrm{\Delta }=1`$, and the energy gap $`2\mathrm{\Delta }4t^{}`$ just vanishes when $`t^{}=0.5`$. It is therefore clear that the strong enhancement in the low-energy spectral function for $`t^{}=0.5`$ is due to the gapless single-particle excitations. This magnon softening has been analytically studied in the strong coupling limit, and the results are given in the Appendix. Due to a frustration induced by the NNN hopping term, the long-wavelength magnon-mode energy is suppressed, and for the simple cubic lattice it vanishes at $`t^{}=1/2`$. At and above this critical value the Néel temperature $`T_\mathrm{N}`$ vanishes, indicating that the Néel state is unstable. The new phase which is stabilized beyond this point is a F-AF phase, involving antiferromagnetic ordering of spins in planes (say parallel to the x-y plane), and ferromagnetic alignment of spins along the z direction. Hence $`t^{}=1/2`$ marks the phase boundary between the AF phase with ordering wavevector $`𝐐=(\pi ,\pi ,\pi )`$ and the F-AF phase with $`𝐐=(\pi ,\pi ,0)`$. For large but finite $`U`$, there is a relative stiffening of the magnon modes due to the 3rd neighbour ferromagnetic coupling induced by the NNN hopping. This reflects a reduction in the degree of frustration, so that a slightly higher $`t^{}`$ value is required to suppress the magnon velocity to zero. From a study of the $`U`$-dependence of this critical $`t^{}`$ value, the magnetic phase diagram of the three dimensional Hubbard model with NNN hopping has been obtained recently. ## V Conclusions In conclusion, we have studied the spectral function of transverse spin fluctuations in an antiferromagnet using a unified approach which includes the contributions from both single-particle excitations and collective magnon excitations. For the NN hopping model in two dimensions, the integrated spectral weight was studied in the whole $`U`$ range, and the collective magnon contribution was found to be dominant for $`U>2.5`$, so that an effective spin description of the AF state is appropriate down to a surprisingly low $`U`$ value. In the weak coupling limit the sharp fall-off of the gapless magnon contribution essentially leaves only the single-particle contribution having a minimum-energy threshold of $`2\mathrm{\Delta }`$, so that the spectrum of transverse spin fluctuations effectively acquires a gap at low energies. The evolution of the spectral function with increasing NNN hopping, which reduces the energy gap, shows a shift of the spectral function towards lower energy due to magnon softening, and also a significant rise at low energy due to single-particle excitations. ## Appendix ### Hubbard model with NNN hopping — strong coupling limit In this section we describe for completeness the AF properties of the half-filled Hubbard model with NNN hopping in the strong coupling limit. NNN hopping introduces a frustration in the system through the competing NNN AF interaction which enhances spin fluctuations and destabilizes the Néel state. Focussing on the enhancement of transverse spin fluctuations at the RPA level by NNN hopping, we examine the spectrum of the collective (magnon) excitations, and their contribution to the quantum spin-fluctuation correction $`\delta m_{\mathrm{SF}}`$ to the sublattice magnetization in two dimensions, and the reduction in the Néel temperature $`T_\mathrm{N}`$ in three dimensions. As discussed earlier, the quasiparticle amplitudes $`a_{𝐤\sigma (\pm )}`$ and $`b_{𝐤\sigma (\pm )}`$ which form the eigenvectors of the HF Hamiltonian remain unchanged by the NNN hopping $`t^{}`$, provided the charge gap is finite, which is very much so in the strong coupling limit. Therefore the only change in $`\chi ^0(𝐪\omega )`$ arises from the change in the quasiparticle energy expression given in Eq. (4), which appear in the energy denominators in Eq. (8), and we obtain, for the AA matrix element, for example $`[\chi ^0(𝐪\omega )]_{\mathrm{AA}}`$ $`=`$ $`{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{a_𝐤^2a_{𝐤𝐪}^2}{\sqrt{\mathrm{\Delta }^2+ϵ_𝐤^2}+\sqrt{\mathrm{\Delta }^2+ϵ_{𝐤𝐪}^2}+(ϵ_{𝐤𝐪}^{}ϵ_𝐤^{})+\omega }}`$ (24) $`+`$ $`{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{a_𝐤^2a_{𝐤𝐪}^2}{\sqrt{\mathrm{\Delta }^2+ϵ_𝐤^2}+\sqrt{\mathrm{\Delta }^2+ϵ_{𝐤𝐪}^2}+(ϵ_𝐤^{}ϵ_{𝐤𝐪}^{})\omega }}.`$ (25) Substituting $`a_𝐤^2=a_𝐤^21ϵ_𝐤^2/4\mathrm{\Delta }^2`$ and $`a_𝐤^2=a_𝐤^2ϵ_𝐤^2/4\mathrm{\Delta }^2`$ in the strong coupling limit, expanding the denominator in powers of $`t/\mathrm{\Delta }`$, $`t^{}/\mathrm{\Delta }`$, $`\omega /\mathrm{\Delta }`$, and systematically retaining terms only up to order $`t^2/\mathrm{\Delta }^2`$ and $`t^{}_{}{}^{}2/\mathrm{\Delta }^2`$, we obtain after performing the $`𝐤`$-sums in two dimensions for a square lattice with $`_𝐤ϵ_𝐤^2=4t^2`$, $`_𝐤ϵ_𝐤^{}_{}{}^{}2=4t^{}_{}{}^{}2`$, $`_𝐤ϵ_𝐤^{}ϵ_{𝐤𝐪}^{}=4t^{}_{}{}^{}2\mathrm{cos}q_x\mathrm{cos}q_y`$, $`[\chi ^0(𝐪\omega )]_{\mathrm{AA}}`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Delta }}}\left[1{\displaystyle \frac{4t^2}{\mathrm{\Delta }^2}}+{\displaystyle \frac{2t^{}_{}{}^{}2}{\mathrm{\Delta }^2}}(1\mathrm{cos}q_x\mathrm{cos}q_y){\displaystyle \frac{\omega }{2\mathrm{\Delta }}}\right]`$ (26) $`=`$ $`{\displaystyle \frac{1}{U}}\left[1{\displaystyle \frac{2t^2}{\mathrm{\Delta }^2}}\left(1+{\displaystyle \frac{\omega }{2J}}\right)+{\displaystyle \frac{2t^{}_{}{}^{}2}{\mathrm{\Delta }^2}}(1\mathrm{cos}q_x\mathrm{cos}q_y)\right],`$ (27) where $`2\mathrm{\Delta }=mU(12t^2/\mathrm{\Delta }^2)U`$ and $`J=4t^2/U`$. Similarly evaluating the other matrix elements,, we obtain $$[1U\chi ^0(𝐪\omega )]=\frac{2t^2}{\mathrm{\Delta }^2}\left[\begin{array}{cc}1\frac{J^{}}{J}(1\gamma _𝐪^{})+\frac{\omega }{2J}& \gamma _𝐪\\ \gamma _𝐪& 1\frac{J^{}}{J}(1\gamma _𝐪^{})\frac{\omega }{2J}\end{array}\right],$$ (28) where $`\gamma _𝐪=(\mathrm{cos}q_x+\mathrm{cos}q_y)/2`$ and $`\gamma _𝐪^{}=\mathrm{cos}q_x\mathrm{cos}q_y`$. Here $`J=4t^2/U`$ and $`J^{}=4t^{}_{}{}^{}2/U`$ are the NN and NNN spin couplings in the equivalent Heisenberg model, and the NNN term $`J^{}(1\gamma _𝐪^{})`$ directly leads to a softening of the magnon mode energies. Substituting in the RPA expression, we finally obtain for the transverse spin propagator $$[\chi ^+(𝐪\omega )]=\frac{1}{2}\left(\frac{2J}{\omega _𝐪}\right)\left[\begin{array}{cc}1\frac{J^{}}{J}(1\gamma _𝐪^{})\frac{\omega }{2J}& \gamma _𝐪\\ \gamma _𝐪& 1\frac{J^{}}{J}(1\gamma _𝐪^{})+\frac{\omega }{2J}\end{array}\right].\left(\frac{1}{\omega \omega _𝐪+i\eta }\frac{1}{\omega +\omega _𝐪i\eta }\right),$$ (29) where the magnon-mode energy $`\omega _𝐪`$ is given by $$\left(\frac{\omega _𝐪}{2J}\right)^2=\left\{1\frac{J^{}}{J}(1\gamma _𝐪^{})\right\}^2\gamma _𝐪^2.$$ (30) In the long wavelength limit $`(q0)`$, with $`\gamma _𝐪^{}=\mathrm{cos}q_x\mathrm{cos}q_y(1q_x^2/2)(1q_y^2/2)=1q^2/2`$, and $`\gamma _𝐪=(\mathrm{cos}q_x+\mathrm{cos}q_y)/21q^2/4`$, the magnon energy reduces to: $$\omega _𝐪=\sqrt{2}Jq\left(1\frac{2J^{}}{J}\right)^{1/2}$$ (31) showing the strong softening of low-energy modes by the NNN hopping. The spin-wave velocity vanishes in the limit $`J^{}/J1/2`$. The magnon density of states evaluated from Eq. (15) is shown in Fig. 8 for different values of the ratio $`J^{}/J`$. NNN hopping clearly softens the magnon energies, and transfers the magnon spectral weight from the high-energy to the low-energy region. The strong softening of the magnon-mode energies suggests an enhancement in the transverse spin fluctuations. To examine this we evaluate the transverse spin correlations as described earlier. From Eq. (14) for the transverse spin propagator, after performing the frequency integral we obtain the local transverse spin correlations $$S^+S^{}+S^{}S^+_{\mathrm{RPA}}=\underset{𝐪}{}\frac{2J}{\omega _𝐪}\left\{1\frac{J^{}}{J}(1\gamma _𝐪^{})\right\}.$$ (32) The spin-fluctuation correction to sublattice magnetization is then obtained from $$\delta m_{\mathrm{SF}}=\frac{S^+S^{}+S^{}S^+_{\mathrm{RPA}}}{S^+S^{}S^{}S^+_{\mathrm{RPA}}}1$$ (33) where the denominator $`S^+S^{}S^{}S^+_{\mathrm{RPA}}`$ is precisely 1 for the $`S=1/2`$ system due to the commutation relation $`[S^+,S^{}]=2S^z`$. The spin-fluctuation correction to sublattice magnetization, evaluated from Eqs. (17), (18) is shown in Fig. 9, showing the rapid rise in transverse spin fluctuations with the frustrating NNN spin coupling $`J^{}`$. We now consider the reduction in the Néel temperature in three dimensions due to the frustrating NNN spin coupling. For this purpose we consider the NNN hopping model on a simple cubic lattice. In this case the lattice free-particle energies are $`ϵ_𝐤`$ $`=`$ $`2t(\mathrm{cos}k_x+\mathrm{cos}k_y+\mathrm{cos}k_z),`$ (34) $`ϵ_𝐤^{}`$ $`=`$ $`4t^{}(\mathrm{cos}k_x\mathrm{cos}k_y+\mathrm{cos}k_y\mathrm{cos}k_z+\mathrm{cos}k_z\mathrm{cos}k_x),`$ (35) and simple extension of the earlier treatment for the two-dimensional case leads to the following result for the transverse spin propagator at the RPA level $$[\chi ^+(𝐪\omega )]=\frac{1}{2}\left(\frac{3J}{\omega _𝐪}\right)\left[\begin{array}{cc}1\frac{2J^{}}{J}(1\gamma _𝐪^{})\frac{\omega }{3J}& \gamma _𝐪\\ \gamma _𝐪& 1\frac{2J^{}}{J}(1\gamma _𝐪^{})+\frac{\omega }{3J}\end{array}\right].\left(\frac{1}{\omega \omega _𝐪+i\eta }\frac{1}{\omega +\omega _𝐪i\eta }\right),$$ (36) where $`\gamma _𝐪=(\mathrm{cos}q_x+\mathrm{cos}q_y+\mathrm{cos}q_z)/3`$ and $`\gamma _𝐪^{}=(\mathrm{cos}q_x\mathrm{cos}q_y+\mathrm{cos}q_y\mathrm{cos}q_z+\mathrm{cos}q_z\mathrm{cos}q_x)/3`$. The magnon-mode energy $`\omega _𝐪`$ is given by $$\omega _𝐪=3J\left[\left\{1\frac{2J^{}}{J}(1\gamma _𝐪^{})\right\}^2\gamma _𝐪^2\right]^{1/2}.$$ (37) For small $`q`$, with $`\gamma _𝐪^{}1q^2/3`$ and $`\gamma _𝐪1q^2/6`$, the magnon energy reduces to $$\omega _𝐪=\sqrt{3}Jq\left(1\frac{4J^{}}{J}\right)^{1/2}$$ (38) which vanishes in the limit $`J^{}/J1/4`$ due to the frustration effect of the NNN coupling $`J^{}`$. The softening of the low-energy magnon spectrum has a bearing on the Néel temperature, as discussed below. Within the renormalized spin-fluctuation theory, the Néel temperature $`T_\mathrm{N}`$ is obtained from the isotropy condition $`S^+S^{}+S^{}S^+_{T=T_\mathrm{N}}=\frac{2}{3}S(S+1)`$. For the NN coupling model ($`J^{}=0`$), the Néel temperature was obtained earlier as $`T_\mathrm{N}=zJ\frac{S(S+1)}{3}f_{\mathrm{SF}}^1`$ for the general case of spin $`S`$ and $`z`$ nearest neighbors on a hypercubic lattice. For the simple cubic lattice the spin-fluctuation factor $`f_{\mathrm{SF}}_𝐪1/(1\gamma _𝐪^2)=1.517`$, and for $`S=1/2`$ this reduces to $`T_\mathrm{N}/J=0.989`$. Extending this analysis to the present case, from the expression for the magnon propagator in Eq. (20) we obtain $$T_\mathrm{N}=\frac{3J}{2}\left[\underset{𝐪}{}\frac{1\frac{2J^{}}{J}(1\gamma _𝐪^{})}{\left\{1\frac{2J^{}}{J}(1\gamma _𝐪^{})\right\}^2\gamma _𝐪^2}\right]^1$$ (39) The Néel temperature, evaluated from the above equation, is shown in Fig. 10 as a function of $`J^{}`$. The rapid reduction of $`T_\mathrm{N}`$ with $`J^{}`$ and the vanishing at $`J^{}=J/4`$ is due to the enhancement of transverse spin fluctuations arising from the frustration-induced softening of the long-wavelength, low-energy magnon modes. The instability at $`J^{}=J/4`$ is towards a F-AF phase with ordering wavevector $`𝐐=(\pi ,\pi ,0)`$, involving antiferromagnetic ordering of spins in planes (say parallel to the x-y plane), and ferromagnetic alignment along the z direction.
warning/0006/hep-th0006228.html	ar5iv	text	# Bosonic D-Branes at Finite Temperature ## Abstract We derive the finite temperature description of bosonic D-branes in the thermo field approach. The results might be relevant to the study of thermical properties of D-brane systems. During the past several years, researches on the properties of D-branes, viewed either as solitons of low energy field theories or as states in perturbative spectra of strings , has continued to grow unabated. Theoretical studies predicted perturbative as well as non-perturbative dualities among string theories and M-theory. D-branes offered new approaches to gauge theories , a first microscopic description of the Beckenstein-Hawking entropy of black-holes and new insights in string cosmology . More recently, D-branes have been used to conjecture a relationship between gravity and quantum field theory and to study the stable non-BPS states of string spectra . In lights of these achivements, its is worthwhile to understand deeper the physical properties of D-branes. In this Letter we aim at constructing bosonic D-brane at finite temperature as boundary states of perturbative closed strings . Although the construction of a single D-brane works in principle for any temperature, in considering statistical ensembles one should limit the discussion to low temperatures mainly because of the following reason. The bosonic string theories contain tachyons and, at zero mass level, dilatons and gravitons. The presence of these fields make the ensemble of strings to suffer a first order phase transition at $`T_c<T_H`$, where $`T_H`$ is the Hagedorn temperature. As was argued in , the latent heat involved in the transition is large enough to break down the notion of temperature shortly above $`T_c`$ (for a review see .) The thermodynamics of a generic Dp-brane in the imaginary time formalism was studied for the first time in while the D-brane free energy in external electro-magnetic field has been computed in and in by using a boundary state formalism. In what follows, we are going to apply the thermo field theory to construct the equations that define the D-branes and the boundary states of closed strings that satisfy these equations. We note that the quasi-particle picture in this theory no longer holds at sufficient large temperatures. This might be another reason for studying the ensembles of branes at low temperatures. However, it is not clear yet what is the validity of the quasi-particle interpretation in the case of strings and to what extent it may affect the ensembles of branes at high temperature. We hope to clarify these aspects in a forthcoming paper . Consider a free bosonic closed string in the Minkowski spacetime and in the conformal gauge at zero temperature $`T=0`$. We denote by $`\alpha _n^\mu `$ and $`\beta _n^\mu `$, $`\mu =0,1,\mathrm{},25`$, $`nZ^{}`$, the right- and left-moving modes, respectively. In order to apply the thermo field method we pass to the oscillator description of string modes. This is given in terms of creation and annihilation operators for each mode $`n`$ and in each sector $`A_n^\mu `$ $`=`$ $`\alpha _n^\mu ,A_n^\mu =\alpha _n^\mu `$ (1) $`B_n^\mu `$ $`=`$ $`\beta _n^\mu ,B_n^\mu =\beta _n^\mu ,`$ (2) where $`nN^{}`$. In units such that $`\mathrm{}=1`$ the above operators describe independent harmonic oscillators of frequencies $`\omega _n=n`$. The vacuum states of right- and left-moving sectors are denoted by $`\|0_\alpha `$ and $`\|0_\beta `$, respectively, and the fundamental state of string is given by the following product $$\|0=\|0_\alpha \|0_\beta \|p=0,$$ (3) where $`\|p`$ is the eigenstate of the linear momentum of the center-of-mass corresponding to the eigenvalues $`p^\mu `$. A D$`p`$-brane is obtained by imposing the Neumann boundary conditions along $`a=0,1,\mathrm{},p`$ directions and the Dirichlet boundary conditions along $`i=p+1,\mathrm{},25`$ directions of the open string at the endpoint $`\sigma =0`$. These are equivalent with the following operatorial equations on the Hilbert space of the closed string $`(A_n^a+B_n^a)\|B_{mat}`$ $`=`$ $`(A_n^a+B_n^a)\|B_{mat}=0`$ (4) $`(A_n^iB_n^i)\|B_{mat}`$ $`=`$ $`(A_n^iB_n^i)\|B_{mat}=0,`$ (5) for oscillator modes $`n>0`$, and $$\widehat{p}^a\|B_{mat}=(\widehat{X}^iy^i)\|B_{mat}=0,$$ (6) for the momentum and position operators of the center-of-mass of string. Here, $`\{y^i\}`$ represent the coordinates of the D$`p`$-brane in the transverse space. Let us construct the counterpart of equations (5) and (6) at $`T0`$ using the thermo field approach . To this end, firstly we have to double the system. We denote by ~ the quantities that correspond to an identical copy of the original bosonic string. The corresponding operators obey the usual commutation rules, the right- and left-sector operators commute and they also commute with the operators corresponding to the original string. For example $$[A_n^\mu ,A_m^\nu ]=[\stackrel{~}{A}_n^\mu ,\stackrel{~}{A}_m^\mu ]=\delta _{n,m}\eta ^{\mu \nu }$$ (7) $$[A_n^\mu ,\stackrel{~}{A}_m^\nu ]=[A_n^\mu ,\stackrel{~}{A}_m^\nu ]=[A_n^\mu ,\stackrel{~}{B}_m^\nu ]=\mathrm{}=0.$$ (8) The extended Hilbert space of the total system is given by the direct product of the two Hilbert spaces of closed strings $$H_0=H\times \stackrel{~}{H}$$ (9) and we denoted a state from $`H_0`$ by $`\|`$. The vacuum state of the total system is given by $`\|0`$ $`=`$ $`\|0_\alpha \|0_\beta =(\|0_\alpha \|\stackrel{~}{0}_\alpha )(\|0_\beta \|\stackrel{~}{0}_\beta )`$ (10) $`=`$ $`(\|0_\alpha \|0_\beta )(\|\stackrel{~}{0}_\alpha \|\stackrel{~}{0}_\beta ),`$ (11) where the last equality is a consequence of the fact that the original string and the ~ string are independent. The first line in (11) shows explicitely the doubling of each oscillator while the second one shows the string-tilde string structure of the vacuum state. In order to obtain the fundamental state of the enlarged system we have to multiply (11) by $`\|p\|\stackrel{~}{p}`$. The thermal description of the system can be obtained from the above one by acting with Bogoliubov operators $`G_n^\alpha `$ and $`G_n^\beta `$ on each sector of the Hilbert space and on the creation and annihilation operators. The $`G`$-operators are defined as follows $`G_n^\alpha `$ $`=`$ $`i\theta (\beta _T)(A_n\stackrel{~}{A}_nA_n^{}\stackrel{~}{A}_n^{})`$ (12) $`G_n^\beta `$ $`=`$ $`i\theta (\beta _T)(B_n\stackrel{~}{B}_nB_n^{}\stackrel{~}{B}_n^{}).`$ (13) Here, $`\theta _n(\beta _T)`$ is real and depends on the statistics of the $`n`$th oscillator $$\mathrm{cosh}\theta _n(\beta _T)=(1e^{\beta _Tn})^{\frac{1}{2}}.$$ (14) Therefore, $`\theta `$ is the same in both right- and left-sectors for a given mode $`n`$. The dot in (13) means the usual scalar product in the Minkowski space $`A_n\stackrel{~}{A}_n=A_n^\mu \stackrel{~}{A}_{n\mu }`$. Since the oscillators are independent, the vacuum state at $`T0`$ has the following structure $$\|0(\beta _T)=\underset{n>0}{}e^{iG_n^\alpha }\|0_\alpha \underset{m>0}{}e^{iG_m^\beta }\|0_\beta .$$ (15) The creation and annihilation operators at $`T0`$ corresponding to (15) are obtained by acting on the zero temperature operators $`\{A^{},A,\stackrel{~}{A}^{},\stackrel{~}{A}\}`$ and $`\{B^{},B,\stackrel{~}{B}^{},\stackrel{~}{B}\}`$ with the $`G`$-operators given in (13). Taking into account the algebraic properties of $`G_n^\alpha `$ and $`G_n^\beta `$ it is easy to show that the finite temperature annihilation operators can be cast into the following form $`A_n^\mu (\beta _T)`$ $`=`$ $`e^{iG_n^\alpha }A_n^\mu e^{iG_n^\alpha }=u_n(\beta _T)A_n^\mu v_n(\beta _T)\stackrel{~}{A}_n^\mu `$ (16) $`\stackrel{~}{A}_n^\mu (\beta _T)`$ $`=`$ $`e^{iG_n^\alpha }\stackrel{~}{A}_n^\mu e^{iG_n^\alpha }=u_n(\beta _T)\stackrel{~}{A}_n^\mu v_n(\beta _T)A_n^\mu `$ (17) in the right-moving sector and $`B_n^\mu (\beta _T)`$ $`=`$ $`e^{iG_n^\beta }B_n^\mu e^{iG_n^\beta }=u_n(\beta _T)B_n^\mu v_n(\beta _T)\stackrel{~}{B}_n^\mu `$ (18) $`\stackrel{~}{B}_n^\mu (\beta _T)`$ $`=`$ $`e^{iG_n^\beta }\stackrel{~}{B}_n^\mu e^{iG_n^\beta }=u_n(\beta _T)\stackrel{~}{B}_n^\mu v_n(\beta _T)B_n^\mu `$ (19) in the left-moving sector. Here, $$u_n(\beta _T)=\mathrm{cosh}\theta _n(\beta _T),v_n(\beta _T)=\mathrm{sinh}\theta _n(\beta _T).$$ (20) Since the operators $`\widehat{p}`$, $`\widehat{X}`$, $`\widehat{\stackrel{~}{p}}`$ and $`\widehat{\stackrel{~}{X}}`$ commute with all oscillator operators, they are not affected by the transformations of the type (17) and (19). If we construct zero mode $`G`$-operators as we did for the oscillators we see that the momenta commute with them. Therefore, we take the position and momenta operators of both the string and the tilde string to be invariant under the transformations above. The corresponding eigenstates of the momenta operators are taken to be invariant, too. Then it is not difficult to see that the string coordinate operators $`X^\mu (\tau ,\sigma )(\beta _T)`$ and $`\stackrel{~}{X}^\mu (\tau ,\sigma )(\beta _T)`$ constructed from (17) and (19) are solutions of the string equations of motion ($`G_n^\alpha `$ and $`G_n^\beta `$ do not act on the two dimensional wave functions.) Therefore, we may construct the D$`p`$-brane states at $`T0`$ by imposing the Neumann and Dirichlet boundary conditions on the appropriate spacetime directions $$_\tau X^a(\tau ,\sigma )(\beta _T)\|_{\tau =0}=0,X^i(\tau ,\sigma )(\beta _T)\|_{\tau =0}=y^i,$$ (21) and similarly for $`\stackrel{~}{X}^\mu (\tau ,\sigma )(\beta _T)`$. Here, $`a=0,1,\mathrm{},p`$ and $`i=p+1,\mathrm{},25`$. When imposed on the extended Hilbert space, these relations define two sets of equations that determine D$`p`$\- and D̃$`p`$-boundary states at $`T0`$. We denote these states by $`\|B_{mat}(\beta _T)`$ and $`\|\stackrel{~}{B}_{mat}(\beta _T)`$, respectively. If we introduce the folowing matrices $$S^{\mu \nu }=(\eta ^{ab},\delta ^{ij}),\stackrel{~}{S}^{\mu \nu }=(\stackrel{~}{\eta }^{ab},\stackrel{~}{\delta }^{ij}),$$ (22) we can see that the equations defining D$`p`$\- and D̃$`p`$-branes have the following form $$[u_n(\beta _T)(A_n^\mu +S_\nu ^\mu B_n^\nu )+v_n(\beta _T)(\stackrel{~}{B}_n^\mu +\stackrel{~}{S}_\nu ^\mu \stackrel{~}{A}_n^\nu )]\|B_m(\beta _T)=0$$ (23) $$[u_n(\beta _T)(A_n^\mu +S_\nu ^\mu B_n^\nu )+v_n(\beta _T)(\stackrel{~}{B}_n^\mu +\stackrel{~}{S}_\nu ^\mu \stackrel{~}{A}_n^\nu )]\|B_m(\beta _T)=0$$ (24) and $$[u_n(\beta _T)(\stackrel{~}{A}_n^\mu +\stackrel{~}{S}_\nu ^\mu \stackrel{~}{B}_n^\nu )+v_n(\beta _T)(B_n^\mu +S_\nu ^\mu A_n^\nu )]\|\stackrel{~}{B}_m(\beta _T)=0$$ (25) $$[u_n(\beta _T)(\stackrel{~}{A}_n^\mu +\stackrel{~}{S}_\nu ^\mu \stackrel{~}{B}_n^\nu )+v_n(\beta _T)(B_n^\mu +S_\nu ^\mu A_n^\nu )]\|\stackrel{~}{B}_m(\beta _T)=0$$ (26) If for some mode $`n`$ and some critical temperature $`T_c`$ the equation $`u_n(\beta _{T_c})=v_n(\beta _{T_c})`$ holds, then the corresponding equations (23),(24) and (25),(26) are identical and give the same $`n`$th oscillator contribution to $`\|B_{mat}(\beta _T)`$ and $`\|\stackrel{~}{B}_{mat}(\beta _T)`$. This condition is actually equivalent to the limit $`T_c\mathrm{}`$ for any finite $`n`$, case in which (23),(24) and (25),(26) are identical for all $`n`$. The equations above are consistent with the zero temperature description of the D$`p`$-branes as boundary states since in the limit $`T0`$ they reduce to two copies of the equations (5). Let us look for the solutions of the equations above. We put the state $`\|B_{mat}(\beta _T)`$ under the form $$\|B_{mat}(\beta _T)=\widehat{B}_{mat}\|0(\beta _T)=\widehat{B}_{mat}(\beta _T)\|0.$$ (27) Then, by analogy with the zero temperature limit, we can write the operator $`\widehat{B}_{mat}(\beta _T)`$ in the following way $$\widehat{B}_{mat}(\beta _T)\delta ^{25p}(\widehat{X}^iy^i)\widehat{O}_{mat}\underset{k=1}{\overset{\mathrm{}}{}}e^{iG_k}.$$ (28) After some simple algebra one can see that the equations (23),(24) have solution of the form (27). Moreover, this solution is not unique since it can be constructed with the following different operators $`\widehat{O}_{mat}^1`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{A_n^{}SB_n^{}}{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}e^{\stackrel{~}{B}_r^{}\stackrel{~}{S}\stackrel{~}{A}_r^{}}`$ (29) $`\widehat{O}_{mat}^{2,3}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{A_n^{}SB_n^{}}\times 1\pm 1\times {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\stackrel{~}{B}_n^{}\stackrel{~}{S}\stackrel{~}{A}_n^{}}.`$ (30) Thus, the D$`p`$-brane states at $`T0`$ from the extended Hilbert space are degenerate in the sense of the equations above. The same considerations apply to $`\|\stackrel{~}{B}_{mat}(\beta _T)`$ states. In this case we obtain the first solutions from (30) as a consequence of the fact that the right- and left-modes commute for either the original string as well as for its copy. Each of D$`p`$\- or D̃$`p`$-branes depend on the oscillator operators of the original as well as the tilde strings. However, according to the thermofield theory, the physical quantities in a thermal D$`p`$-brane state may depend only on the vev of the operators without tilde . In this sense, the solutions derived above in (30) are equivalent. Note that one has to include in $`\|B_{mat}(\beta _T)`$ the eigenvectors of momenta operators at $`p=0`$. However, since these operators and their eigenvalues do not transform under $`G`$-transofrmations, their contribution is the same as in $`T=0`$ case. The normalization constant that enter (27) is related, at $`T=0`$, to the tension of the brane . To determine it at finite temperature, one has to perform calculations in both open and closed string sectors . The above boundary states make sense only if they belong to the physical sector of the theory. At $`T=0`$, this sector is defined using the conformal invariance of the string theory. The same definition can be applied here since the conformal invariance is not broken. Indeed, using the fact that $`G`$-operators commute, it is easy to show that the string oscillation modes statisfy the same algebra at finite temperature. Then, since $`X^\mu (\tau ,\sigma )(\beta _T)`$ and $`\stackrel{~}{X}^\mu (\tau ,\sigma )(\beta _T)`$ satisfy the string equations of motion the operators $`L_n(\beta _T)`$ and $`\stackrel{~}{L}_n(\beta _T)`$ have formally the same expansion in terms of modes $`\alpha _m(\beta _T)`$ as the Virasoro operators at zero temperature and satisfy the same algebra in both right- and left-sectors and for both the original string and the tilde string. The algebras in all of these sectors are independent of each other. Using the conformal symmetry, one can impose the physical state conditions on the Fock space of the extended system as for $`T=0`$ . Also, the problem of mantaining the conformal invariance is related to the possibility of mapping the open string sector to the closed one and to constructing the BRST invariant states. Note that the conformal invariance at finite temperature is due to the specific form of $`G`$-operators taken in (13). This is not the most general form allowed in thermo field approach . Nevertheless, the conformal symmetry is mantained by other tranformations, too. The point here is that $`G`$-operators do not mix the various oscillation mode which is a reasonable assumption at low temperature. At high temperature we might expect that $`G`$-operators involve interacting terms between various oscillators, case in which the conformal invariance is not automatically satisfied. If we stick on the definition of strings as conformal invariant field theories in two dimensions, the conformal invariance imposes some constraints on $`G`$-operators. On the other hand, breaking the conformal invariance by $`G`$-operators might be a consequence of the fact that the temperature is not well defined for string systems at $`T>T_c`$ as was showed in . In the case of super D-branes one has to repeat the same construction in both NS-NS and R-R sectors of the closed string. It is quite easy to construct the corresponding Bogoliubov transformations for NS and R spinors and the corresponding fermionic modes at finite temperature. However, it is not clear yet what is the natural extension of the supersymmetry at $`T0`$ of the theory that is supersymmetric at $`T=0`$ and whether the superconformal invariance is broken or not . The supersymmetry in right- and left-modes splits into $`\alpha ^{}p^\mu +{\displaystyle \underset{n}{}}u_n(\beta _T)_{}X_n^\mu \stackrel{susy}{}{\displaystyle \underset{t}{}}u_t(\beta _T)\psi _{(t)}^\mu `$ (31) $`{\displaystyle \underset{n}{}}v_n(\beta _T)_{}\stackrel{~}{Y}_n^\mu \stackrel{susy}{}{\displaystyle \underset{t}{}}v_t(\beta _T)\stackrel{~}{\chi }_{(t)}^\mu ,`$ (32) where by $`\stackrel{~}{Y}_n^\mu `$ and $`\stackrel{~}{\chi }_n^\mu `$ we denoted the bosonic and fermionic modes of the tilde string present in the equations of the original string. (Similar expressions should be considered for the tilde string, too.) The problem is due to the presence of $`u_n(\beta _T)`$ and $`v_n(\beta _T)`$ in the bosonic mode expansion and to similar terms $`u_t(\beta _T)`$ and $`v_t(\beta _T)`$ in the fermionic mode expansion, where $`tZ^++1/2`$ in NS sector and $`tZ^+`$ in R sector, respectively. We note that the second set of coefficients is different from the first one since it is related to Fermi-Dirac distributions of the fermionic modes. In summary, we have used the thermo field theory to find out the equations that define the bosonic D$`p`$-branes at finite temperature. We found that the corresponding boundary states from the extended Hilbert space are degenerate, but equivalent at thermical equilibrium. They do not come in pairs since the solutions are the same for the equations of the original string and for the tilde string. In this respect the states that describe D$`p`$-branes and D̃$`p`$-branes are degenerate, opposite to the quasi-particle states in field theory. ###### Acknowledgements. I have greatly benefited from discussion with M. C. B. Abdalla, N. Berkovits, A. L. Gadelha, B. M. Pimentel, H. Q. Placido, D. L. Nedel and B. C. Vallillo. I also acknowledge to M. A. De Andrade and J. A. Helayel-Neto for hospitality at GFT-UCP where part of this work was done and to A. Sen for correspondence. The work was financially supported by a FAPESP postdoc fellowship.
warning/0006/nucl-th0006077.html	ar5iv	text	# Influence of the Coulomb Interaction on the Chemical Equilibrium of Nuclear Systems at Break-Up ## Abstract The importance of a Coulomb correction to the formalism proposed by Albergo et al. for determining the temperatures of nuclear systems at break-up and the densities of free nucleon gases is discussed. While the proposed correction has no effect on the temperatures extracted based on double isotope ratios, it becomes non-negligible when such temperatures or densities of free nucleon gases are extracted based on multiplicities of heavier fragments of different atomic numbers. PACS numbers: 25.70.Mn, 25.70.Pq, 25.70.-z, 24.10.Pa The formalism by Albergo, Costa, Costanzo, and Rubbino (ACCR) offers simple and elegant prescriptions for the experimental determination of the temperature $`T`$, and of the free nucleon densities, $`\rho _{nF}`$ and $`\rho _{pF}`$ in a nuclear system at the instance of the break-up. This formalism presumes a specific scenario for the decay of excited systems, similar to the scenarios modeled by the Berlin Microcanonical Metropolis Monte Carlo and the Copenhagen Statistical Multifragmentation models. Therefore, it may be expected to be meaningful in circumstances (in an excitation energy domain), where the use of the above two more complete models can be justified. On the other hand, the ACCR formalism is incompatible with models such as the equilibrium-statistical model GEMINI and the Expanding Emitting Source Model, both of which refer to sequential decay scenarios of systems of uniform density. The purpose of the present paper is to point out that in certain circumstances, the approximation of the free nucleon gas as a collection of non-interacting nucleons, as assumed in the original ACCR approach, may not be sufficiently accurate, and that the inclusion of a proper Coulomb correction term is warranted. Regardless of the magnitude of the effects of such a correction, its inclusion is warranted already by didactical considerations. The ACCR formalism refers to a fragment production scenario in which an equilibrated freeze-out/break-up configuration emerges from initial compression and expansion stages. In this formalism, the average numbers of fragments or clusters of different mass and atomic numbers $`(A,Z)`$ are determined by the requirement of chemical equilibrium between the fragments and free nucleons - neutrons and protons (quite obviously, the fragments are then also in a state of chemical equilibrium among themselves). In Ref. 1, the chemical equilibrium between fragments $`(A,Z)`$ and gases of free neutrons and free protons is described by the equation $$\stackrel{~}{\mu }_{A,Z}B_{A,Z}=Z\stackrel{~}{\mu }_{pF}+(AZ)\stackrel{~}{\mu }_{nF},$$ (1) where $`B_{A,Z}`$ is the fragment binding energy (taken with positive sign) and $`\stackrel{~}{\mu }_{A,Z}`$, $`\stackrel{~}{\mu }_{nF}`$, and $`\stackrel{~}{\mu }_{pF}`$ are the reduced chemical potentials of a fragment $`(A,Z)`$, of a free neutron, and a free proton, respectively. Here, the qualifier “reduced” is used with respect to the term chemical potential and a tilda is used in the respective symbolic representation, to distinguish the quantities involved in Eq. 1 from the true thermodynamical chemical potentials $`\mu _{A,Z}`$, as defined via the equation $$\mu _{A,Z}=[\frac{F(V,T)}{N_{A,Z}}]_{V,T}.$$ (2) In Eq. 2, $`F(V,T)`$ is the free energy of the system at volume $`V`$ and temperature $`T`$. The quantity $`N_{A,Z}`$ is the average number of fragments $`(A,Z)`$, and the partial derivative is taken at constant temperature and volume. Note that, unlike their “reduced” counterparts $`\stackrel{~}{\mu }_{A,Z}`$, the true chemical potentials $`\mu _{A,Z}`$ include not only the binding energy term $`B_{A,Z}`$, but also the energy of interaction of the fragments $`(A,Z)`$ with the relevant mean Coulomb field. When true chemical potentials are considered, the chemical equilibrium is expressed through the following equation: $$\mu _{A,Z}=Z\mu _{pF}+(AZ)\mu _{nF}V_{coul}^{pF}(Z),$$ (3) where $`V_{coul}^{pF}(Z)`$ is the average energy of mutual Coulomb interaction of $`Z`$ free protons. The latter Coulomb interaction energy must be subtracted on the right-hand side of Eq. 3 to compensate for the fact, that in a mean-field type of “bookkeeping” of the Coulomb energy that is included in the term $`Z\mu _{pF}`$, the mutual Coulomb interaction energy is double-counted. For some purposes, one may consider the convenient notion of a reduced chemical potential $`\stackrel{~}{\mu }(A,Z)`$, which excludes the fragment binding energies and the interaction with the mean Coulomb field. In Maxwell-Boltzmann statistics, the latter quantity has a simple relationship to the density $`\rho _{A,Z}`$ of fragments $`(A,Z)`$ representing the number of fragments per unit break-up volume: $$\rho _{A,Z}=\frac{A^{\frac{3}{2}}\omega _{A,Z}(T)}{\lambda _T^3}e^{\stackrel{~}{\mu }_{A,Z}/T}.$$ (4) Eqs. 4 and Eq. 1 are the two fundamental equations used in Ref. to establish the relationship between the yields of various fragments and the characteristics of a chemically equilibrated nuclear system at break-up. In Eq. 4, $`\lambda _T=h/\sqrt{2\pi m_oT}`$ is the nucleon thermal wave-length ($`m_o`$ is the mass of a nucleon), and $`\omega _{A,Z}(T)`$ is the temperature-dependent internal partition function of fragment $`(A,Z)`$: $$\omega _{A,Z}(T)=\mathrm{\Sigma }_k(2s_k^{A,Z}+1)e^{E_k^{A,Z}/T},$$ (5) where the summation extends over all bound states of the fragment $`(A,Z)`$ with spins $`s_k^{A,Z}`$ and excitation energies $`E_k^{A,Z}`$. An inspection and comparison of Eq. 1, to the more fundamental Eq. 3 reveals a lack of symmetry of the former equation. While the fragment $`(A,Z)`$ side of the balance includes the mutual interaction energy of the constituent nucleons - the binding energy $`B_{A,Z}`$, no equivalent term is present for the free nucleons on the r.h.s. of Eq. 1. Yet, the $`Z`$ free protons do interact among themselves via long-range Coulomb interactions. Therefore, a more complete equation for the chemical equilibrium based on Eq. 3 must include the respective Coulomb interaction term. Note that, on the fragment side of the balance, the mutual Coulomb interaction energy of $`Z`$ protons is included in the binding energy term $`B_{A,Z}`$. One may note also that neither side of Eq. 1 considers explicitly (or implicitly) the Coulomb interaction energy of the $`Z`$ protons with the remaining $`(Z_{system}Z)`$ “spectator” protons of the system. Such an omission, however, may be well justified, as these two Coulomb interaction energies are to a good approximation equal to each other and, hence, cancel mutually. A more complete, Coulomb-corrected equation for the chemical equilibrium of fragments $`(A,Z)`$ and free nucleons, in terms of reduced chemical potentials $`\stackrel{~}{\mu }_{A,Z}`$ has the symmetrical form: $$\stackrel{~}{\mu }_{A,Z}B_{A,Z}=Z\stackrel{~}{\mu }_{pF}+(AZ)\stackrel{~}{\mu }_{nF}+V_{coul}^{pF}(Z).$$ (6) Here, $`V_{coul}^{pF}(Z)`$ represents the average potential energy of the mutual Coulomb interaction of $`Z`$ free protons, an equivalent of the term -$`B_{A,Z}`$. It is worth noting that, here (unlike in Eq. 3) the Coulomb interaction term enters with positive sign, as no Coulomb interaction is included in the reduced free-proton chemical potential $`\stackrel{~}{\mu }_{pF}`$. While it is clear from simple estimates that the Coulomb correction term, $`V_{coul}^{pF}(Z)`$, is of non-negligible magnitude when compared to typical temperatures of the system, it is not obvious how to actually evaluate it. A conservative estimate for the value of $`V_{coul}^{pF}(Z)`$ may be obtained by assuming that this term is equal to the Coulomb interaction energy of $`Z`$ protons uniformly distributed over a spherical volume $`V_{free}`$ = $`Z/\rho _{pF}`$, i.e., distributed with a density equal to that of the gas of free protons $`\rho _{pF}`$ at break-up. In this case, the correction term is independent of the fragment mass number $`A`$: $$V_{coul}^{pF}(Z)=\frac{3}{5}\frac{e^2Z^2}{(\frac{3Z}{4\pi \rho _{pF}})^{1/3}}1.39Z^{5/3}\rho _{pF}^{1/3}(MeV),$$ (7) where $`\rho _{pF}`$ is expressed in units of $`fm^3`$. The presence of the Coulomb correction term $`V_{coul}^{pF}(Z)`$ in Eq. 6 modifies the basic equation 5 of Ref. for the average number of fragments $`(A,Z)`$ per unit break-up volume, $`\rho _{A,Z}`$. It can now be written more accurately as $$\rho _{A,Z}=\frac{A^{3/2}\lambda _T^{3(A1)}\omega _{A,Z}(T)}{2^A}\rho _{pF}^Z\rho _{nF}^{AZ}e^{[B_{A,Z}+V_{coul}^{pF}(Z)]/T},$$ (8) replacing Eq. 4. In Eq. 8, $`\rho _{nF}`$ and $`\rho _{pF}`$ are the densities (i.e., numbers per unit break-up volume) of free neutrons and free protons, respectively. A similar result was obtained earlier based on a more rigorous macrocanonical description of a decaying nuclear system in a freeze-out configuration. It is worth noting that, in more complete theoretical descriptions of equilibrated freeze-out configurations, offered by the Berlin and Copenhagen models, the effects of the Coulomb interaction of the free protons are accounted for in a rigorous fashion, but remain largely transparent to the model users. In practical applications of Eq. 8, ratios of properly selected densities, $`\rho _{A_1,Z_1}/\rho _{A_2,Z_2}`$, are taken and identified with the ratios of the respective experimental yields of fragments $`Y_{A_1,Z_1}/Y_{A_2,Z_2}`$. Such ratios are free of some model parameters (e.g., of the densities of free neutron and proton gases, in the case of double isotope ratios), providing often a simple link between observable yields and selected characteristics of the break-up state. It is clear from Eq. 8 that the introduction of the Coulomb correction term $`V_{coul}^{pF}(Z)`$ is of no consequence when ratios of yields are taken for fragments with identical atomic numbers $`Z`$, i.e., ratios of experimental fragment yields of the type $`Y_{A+1,Z}/Y_{A,Z}`$. In such cases, the corresponding Coulomb correction terms for the two isotopes involved cancel each other. As a result, this correction has no effect on the outcome of an experimental evaluation of break-up temperatures based on double-isotope ratios \- the most common use of the ACCR approach. A similar cancellation does not, however, occur in cases when, e.g., an experimental “thermometer” is constructed from isotone ratios, $`Y_{A+1,Z+1}/Y_{A,Z}`$, or in cases when relative densities of free neutron and proton gases, $`\rho _{nF}/\rho _{pF}`$, are determined based on an isobaric ratio $`Y_{A,Z}/Y_{A,Z+1}`$. To assess the significance of the proposed Coulomb correction, several examples are considered below. In these examples, it is assumed that $`T`$=3.3 MeV (as found for the system S+Ag at E/A=22 MeV) and $`\rho _{pF}=5/(4/3\pi 8.0141)=0.0011`$ fm<sup>-3</sup> (which corresponds to 5 protons in a break-up volume of radius $`R_{breakup}=2.0141^{1/3}`$ fm, as for the system <sup>32</sup>S+<sup>109</sup>Ag). First, consider an evaluation of the ratio of densities of free neutron and free proton gases from the observed isobaric ratios $`Y_{A,Z}/Y_{A,Z+1}`$. Based on Eq. 8, one has $$\frac{\rho _{nF}}{\rho _{pF}}=R_{raw}F_{Coul}^{n/p},$$ (9) where $`R_{raw}`$ is the value of this ratio deduced in the absence of the Coulomb correction (i.e., given by the original ACCR formalism) and $`F_{Coul}^{n/p}`$ is a correction factor resulting from the Coulomb term proposed in the present paper: $$F_{Coul}^{n/p}=e^{[V_{coul}^{pF}(Z+1)V_{coul}^{pF}(Z)]/T},$$ (10) Using Eqs. 10 and 7 and the values of $`T`$=3.3 MeV and $`\rho _{pF}=0.0011`$ fm<sup>-3</sup>, one obtains $`F_{Coul}^{n/p}`$ = 1.1, for the case of the isobaric ratio $`Y_{3,1}/Y_{3,2}`$ (tritium - helium-3), $`F_{Coul}^{n/p}`$ = 1.28, for the case of the isobaric ratio $`Y_{13,6}/Y_{13,7}`$, and $`F_{Coul}^{n/p}`$ = 1.59 when the isobaric ratio $`Y_{34,16}/Y_{34,17}`$ is utilized. This example demonstrates that even in the favorable case of light isobars $`{}_{}{}^{3}H`$ and <sup>3</sup>He, the correction factor is large enough to mandate an inclusion of the proposed Coulomb correction term in the equation for the chemical equilibrium. Certainly, this correction factor is quite sizeable when yields of heavier isobars are utilized for the evaluation of the relative densities of free neutron and free proton gases. It is worth noting that, according to Eq. 10, the Coulomb correction factor, $`F_{Coul}^{n/p}`$, for the relative densities of gases of free neutrons and protons is always greater than unity, since the exponent $`[V_{coul}^{pF}(Z+1)V_{coul}^{pF}(Z)]`$ is positive. This fact reflects the role of the Coulomb energy in a neutron-enrichment of the free nucleon gas, a rather trivial effect that should not be confused with an isospin fractionation driven by an isospin-dependent equation of state of nuclear matter. As a second example, consider the evaluation of the break-up temperature $`T`$ based on a double isotone ratio: $$R_{isotone}=\frac{Y_{A_1,Z_1}/Y_{A_1+1,Z_1+1}}{Y_{A_2,Z_2}/Y_{A_2+1,Z_2+1}}.$$ (11) In such a case, the breakup temperature $`T`$ is ultimately evaluated from the experimentally determined value of the ratio $$\frac{\mathrm{\Delta }B+\mathrm{\Delta }V_{coul}^{pF}}{T}=\frac{\mathrm{\Delta }B}{T_{raw}}=ln(aR_{isotone}),$$ (12) where $`\mathrm{\Delta }B=B_{A_1+1,Z_1+1}B_{A_1,Z_1}+B_{A_2,Z_2}B_{A_2+1,Z_2+1}`$ and $`\mathrm{\Delta }V_{coul}^{pF}=V_{coul}^{pF}(Z_1+1)V_{coul}^{pF}(Z_1)+V_{coul}^{pF}(Z_2)V_{coul}^{pF}(Z_2+1)`$. The parameter $`a`$ in Eq. 12 accounts for intrinsic partition functions $`\omega `$ (see Eq. 5) of the isotones involved, and $`T_{raw}`$ is the break-up temperature obtained using the original ACCR approach. Eq. 12 allows one to express the relevant Coulomb correction factor $`F_{Coul}^T`$ as $$F_{Coul}^T=\frac{T}{T_{raw}}=1+\frac{\mathrm{\Delta }V_{coul}^{pF}}{\mathrm{\Delta }B}.$$ (13) Using Eq. 13, the value of $`\rho _{pF}=.0011`$ fm<sup>-3</sup>, and values of binding energies from the mass tables, one obtains $`F_{Coul}^T=0.96`$ in the case when the experimental double isotone ratio $`(Y_{13,7}/Y_{12,6})/(Y_{4,2}/Y_{3,1})`$ is employed, $`F_{Coul}^T=1.19`$ in the case of the isotone ratio $`(Y_{13,7}/Y_{12,6})/(Y_{10,5}/Y_{9,4})`$, and $`F_{Coul}^T=0.31`$ in the case of the isotone ratio $`(Y_{14,8}/Y_{13,7})/(Y_{3,2}/Y_{2,1})`$. Again, the estimated magnitude of the Coulomb correction factor well warrants an inclusion of the Coulomb correction term in the ACCR formalism. In summary, the importance of a Coulomb correction term to the equation for the chemical equilibrium between fragments and the gas of free nucleons has been demonstrated. The correction term restores the symmetry of the equation defining the equilibrium, when the mutual interaction energies of nucleons in both, bound and free states, are consistently accounted for as done with the more rigorous Eq. 3. While the proposed correction term has no effect on the determination of break-up temperatures based on double isotope ratios and may be small in some cases, its effects on the determination of break-up temperature from isotone ratios and on the determination of the relative densities of free proton and neutron gases may be quite substantial in some other cases. This correction is certainly important in systematic studies of various experimental “thermometers” that, by design, include a large variety of isotonic ratios. Useful discussions with S. Albergo are gratefully acknowledged. This work was supported by the U.S. Department of Energy grant No. DE-FG02-88ER40414.
warning/0006/astro-ph0006302.html	ar5iv	text	# Is there ERE in diffuse galactic light at high latitude? ## 1 Introduction To study the optical properties of interstellar grains two types of approaches are usually followed. The first is to compare the strength of a nebula surface brightness to the value of the source radiation field at the cloud position. More frequently the color of the cloud is compared to that of the illuminating source. Both methods have led some authors (Lynds lynds62 (1962), Guhathakurta & Cutri gu2 (1994), Gordon, Witt and Friedmann gordon (1998) and references therein) to conclude that not all interstellar cloud emission can be accounted for solely by the scattering of starlight. According to these authors, comparison of nebula and source emissions cannot explain the strength of the emission in the $`R`$ and $`I`$ bands, and the red color emission of some, if not most, interstellar clouds. Hence, the clouds must be emitting themselves, which implies the presence of a particular class of dust grain, assumed to absorb UV ambient radiation and to re-emit in the red. The phenomenom is called ERE (Extended Red Emission). Unfortunately, all the attempts to identify ERE carriers have failed. The directions where ERE is found, which for many years consisted of only one direction (the Red Rectangle), now cover nearly the entire sky. The phenomenon is thought to be so important that up to $`50\%`$ of interstellar red emission is attributed to ERE (Gordon, Witt and Friedmann gordon (1998), hereafter GWF). ERE appears in a wide range of different environments and astrophysical objects, including cirrus, nebulae, planetary nebulae, HII regions, novae, and other galaxies. Complete references can be found in the introduction of GWF. Diffuse galactic light (DGL) was identified at the beginning of the century (Struve & Elvey struve36 (1936)) as diffuse interstellar radiation which remains after the subtraction of direct starlight and the light of solar and terrestial origins. It was attributed to starlight scattered by interstellar grains. Toller (toller81 (1981)) has found that the visible surface brightness of the DGL in high latitude directions and the HI column density are roughly proportional. The GWF paper is devoted to detecting ERE in DGL and is divided into two parts. The first part of the GWF paper describes the method the authors have followed to estimate the $`B`$ and $`R`$ surface brightnesses of direct starlight in different directions of the galaxy. In two high galactic latitude regions, region (a) and region (b), the result is substracted from the $`B`$ and $`R`$ Pioneer satellite measurements to yield an estimate of the DGL. The second part of the GWF paper (§ 4 on), in which the authors attempt to demonstrate the existence of ERE in the diffuse interstellar medium, is constructed on the following idea: if we have a reliable model of the starlight scattered from the diffuse interstellar medium, and if the observed DGL is in excess of this estimate, the difference must arise from a non-scattering emission process. To estimate the scattered part of DGL, GWF use the Witt-Peterson model (‘WP’ model hereafter) of the galaxy. From HI data, the model sets up a vast three dimensional representation of the galaxy which GWF use to calculate the amount of scattered-DGL received from any direction. The large grains responsible for the scattered optical DGL also have a thermal $`100\mu `$m emission. Following Toller’s search for an HI counterpart to the DGL, IRAS images are a powerful tool to visualize the medium in which the DGL originates. For the regions of the sky in which GWF have separated DGL from starlight, IRAS images can be used to check the pertinence and accuracy of the WP model. In both regions, the medium has a complicated structure. It is impossible for the WP model, or any other model, to approach the three dimensional structure of these mediums with the necessary precision solely from an estimate of the HI column density. It is surprising that the large differences between the results of the model and the observations are attributed either to variations of the properties of interstellar grains with longitude when the $`B`$ band is concerned, or to ERE for the $`R`$ band. Is it even possible, or furthermore necessary, to model the part of the DGL which comes from the scattering of starlight? Since the interstellar grains are forward scattering light (Henyey & Greenstein henyey (1941)), DGL in one direction must arise from the scattering of the light of stars in close proximity to this direction. Evaluation of the scattered part of the DGL can be restricted to a local estimate of the radiative transfer of the background direct starlight through the interstellar medium. The radiative transfer does not need to be calculated for the entire galaxy. Existence of ERE could have been probed through a direct and local comparison of DGL and direct starlight. Such a comparison (section 5) is more reliable since it will use GWF data only. IRAS images help to understand the variation of the DGL and the starlight surface brightnesses in fields (a) and (b), in relation to the interstellar medium structure. There are clear correlations between IRAS $`100\mu `$m images and the variations of DGL and direct starlight in the two regions where it was isolated by GWF. These correlations will be discussed in section 6. The data presented in the GWF paper are reviewed in section 2. The difficulties which arise when considering the GWF method and the WP model are detailed in section 3 and section 4. An alternative explanation of the GWF data, supported by the IRAS images of fields (a) and (b), is proposed from section 5 to section 6. To conclude, all of GWF data will be explained by scattering of background starlight by the nearby interstellar cirrus. The relative colors of Pioneer observations, the direct starlight and the DGL can be explained with the absence of ERE. ## 2 Data To compare the color of DGL and of direct starlight, GWF have determined the absolute emissions in the $`B`$ and $`R`$ bands in two regions of the sky and calculated the $`R/B`$ ratio of the emissions. This method is difficult to apply to groundbased observations since it requires the substraction of all foreground (zodiacal light and scattering in the earth’s atmosphere, see Leinert et al. leinert98 (1998) for a review) sky emissions. The difficulty is increased by the low level of emission of the diffuse interstellar medium at high galactic latitude. Any such attempt should use space observations and substract the direct emission of the stars, which can be evaluated from star catalogs. Toller (toller81 (1981), see also Leinart et al.) did this work in 192 regions of the sky using Pioneer measurements. In boxes $`5^{}\times 5^{}`$ large, and along two cuts at galactic longitude, $`l=0`$ and $`l=100`$ (region (a) and region (b), figure 4), GWF have separated the respective $`B`$ and $`R`$ emissions of bright stars, stars with $`\mathrm{mag}>6`$, and the DGL. Bright stars with $`m_V<6.5`$ were removed during the Pioneer data reduction. The remaining emission, comprising stars of $`m_V>6.5`$, galaxies, and DGL, corresponds to the Pioneer curves of the figures 1, 2 and 3. The ISGL curves correspond to the expected emission of stars and galaxies with magnitude greater than 6.5. They are calculated from the GWF ‘Master Catalog’ which is a compilation of different star catalogs. Not all stars have a $`B`$ and an $`R`$ measured magnitude. To determine the $`R`$ and $`B`$ magnitudes of a star with only a $`V`$ measurement, GWF use an estimate of the $`A_V`$ value in the star direction. Determination of $`A_V`$ assumes a reddening proportional to the star distance. Figure 1 gives the Pioneer surface brightnesses $`P_B`$ and $`P_R`$, and the ISGL integrated intensities per unit solid angle, $`isgl_B`$ and $`isgl_R`$, as measured from earth, for the $`B`$ and $`R`$ bands, along cuts (a) and (b). ## 3 The GWF’s argumentation ### 3.1 The GWF analysis of figures 1, 2 and 3 GWF’s analysis of the curves, figures 1, 2 and 3, is short enough to be reproduced here: “From Figure 14, it is obvious that the diffuse ISM is redder (larger $`PR/PB`$ ratio) than either the Pioneer measurements or the ISGL. As the scattered component of the diffuse ISM (DGL) is bluer (see § 4.1) than Pioneer measurements, this requires that the nonscattered component is present in the red diffuse ISM intensity. (GWF, p.531)” It is difficult to accept GWF’s premise at face value. Why should the scattered component of the diffuse ISM (DGL) be bluer than Pioneer measurements or than the ISGL? Where in GWF has it been proven that DGL must be bluer than Pioneer measurements? In GWF’s introduction it is ascerted that: “The DGL will have a bluer color than the integrated starlight because scattering is more efficient at shorter wavelengths \[in the blue than in the red\] (GWF p.523)”. The following sentence states: “So, if the diffuse ISM color (red/blue ratio) is as red as or redder than the integrated starlight and other sources of excess red light can be positively excluded, ERE is present (GWF p.523)”. Note that the first of these GWF assumptions is evidently restricted to low column density mediums since an increase of column density will see all blue starlight absorbed before red light and the scattered light color turn to the red. If GWF had related the two sentences to suggest that the source of the scattered light in one direction is the background starlight in that direction, I would fully agree with this assumption since interstellar grains are known to scatter light preferentially in the forward direction. The elimination of the contribution of the bright stars ($`mag<6.5`$) may not modify the ISGL color enough to invalidate the calculations. This is justified, for instance, if the brightest stars are close to the sun or by the small number of bright stars at high latitude. GWF don’t seem to make this connection between DGL and background starlight. Their approach of the problem is to calculate what should be the scattered starlight in each direction, taking into account the transfer of starlight from all directions and throughout the whole galaxy. ### 3.2 The ‘WP’ model #### 3.2.1 Description of the WP model § 4.1 of GWF -where the reader is supposed to find the demonstration that the scattered-DGL should have a bluer color than Pioneer’s and ISGL’s- is dedicated to the comparison of the Witt-Petersson (‘WP’) model with the data. According to GWF, the WP model represents the galaxy as a ‘gigantic reflection nebula’. Calculation of the radiative transfer through the nebula will give a model of the scattered light. This scattered-DGL model will be compared to the observations. In order to compute the radiative transfer of starlight through the ‘gigantic nebula’ the structure of the medium in which light is scattered and the radiation field at each position of the galaxy need to be defined. The result given by the simulation of the radiative transfer through the WP model gives an expected radiation field for all directions of the galaxy. The WP model of the galaxy is constructed from HI surveys. For each direction of space, the HI column density is converted into dust optical depth for the $`B`$ and $`R`$ bands by multiplication with an appropriate factor. The total column density is also divided into interstellar clouds with a spectrum of sizes and optical depths given by Witt et al. (witt97 (1997)). The radiation field GWF use is deduced from Pioneer data, with the bright stars re-integrated. It comprises the DGL. #### 3.2.2 Results Concerning field (a), the model can hardly fit the data in the $`B`$ band, where scattering is supposed to be the only process involved. For field (b), whatever grain parameters are chosen, the model is a factor of 2 to 4 above the data, which GWF interpret as changes in grain properties with longitude. Despite the fact that the WP model does not fit the data when it should, the conclusion of GWF is that the $`PR/PB`$ ratio is a factor of 2 over the WP model expectation. Since $`PB`$ contains starlight and scattered starlight, it is concluded that additional red emission must arise in the diffuse interstellar medium. ## 4 Some problems with the GWF argument ### 4.1 Remarks on the ‘WP’ model Not surprisingly, GWF’s model does not concur with the observations and as such the following questions arise: What if the ‘WP model’ is an inappropriate one on which to model the galaxy? Can the WP model give a representation of the galaxy with enough accuracy to justify the radiative transfer calculations? When the GWF model doesn’t fit the data, why is the discrepancy attributed to the variations of grain properties and not to the inadequacy of the model? Why is it necessary to have a representation of the whole galaxy to model the DGL in two limited regions? Why isn’t the model restricted to fields (a) and (b)? Grains are known to strongly forward scatter starlight, which implies that most of the DGL we receive from one direction originates from the scattering of the light of the stars close to that direction. The radiation field which GWF use in their Monte Carlo simulation of the radiative transfer through the Witt representation of the galaxy includes DGL which should be the result of the simulation. Why isn’t direct starlight, which GWF have estimated in their ‘Master Catalog’ the input radiation field of the simulation? GWF never proved that the observed DGL should be bluer than the ISGL. Only Monte Carlo simulations with the WP model as a representation of the galaxy and an arbitrary radiation field give a theoretical color of the scattered starlight bluer than that of the ISGL. Is it necessary to introduce any model at all? The search for ERE can be restricted to a comparison of the observed color of the ISGL (corrected for reddening) and of the DGL in each of the $`5^{}\times 5^{}`$ areas. In regard to the method and the uncertainties which accompany this model one can at best find GWF conclusions doubtful. ### 4.2 Is the DGL redder than the ISGL? The high values of $`S_R/S_B`$ in figure 3, figure 14 in GWF, are due to the considerable amplification of the error by the successive substractions and divisions which led to the calculation of the DGL color. The problem introduced by the amplification of the error through these operations can be overcome if it is remarked that $`S_R/S_B>isgl_R/isgl_B`$ is equivalent to $`P_R/P_B>isgl_R/isgl_B`$ since $`P_R=S_R+isgl_R`$ and $`P_B=S_B+isgl_B`$. Hence, comparison of the DGL color and the ISGL color can be reduced to a comparison of the Pioneer and the ISGL colors, which can be done with much better accuracy. In both fields the $`2`$ curves $`P_R/P_B`$ and $`isgl_R/isgl_B`$ are very close to each other. Concerning field (b), we have within the error margin: $`P_R/P_B=isgl_R/isgl_B`$, which would yield: $$P_R/P_B=isgl_R/isgl_B=S_R/S_B$$ (1) Except for the two lower latitude points, the same remark applies to field (a). The redder color of the Pioneer data at the two low latitude points will be interpreted in section 6. ## 5 Relations between the colors of DGL, ISGL and Pioneer ### 5.1 IRAS data #### 5.1.1 The diffuse interstellar medium and IRAS images Because grains scatter starlight in the forward direction, it will be considered that DGL in one direction is the light of background stars in the same direction scattered by foreground dust. This is also justified by the large areas which GWF have considered. It is a wide-spread idea that the diffuse galactic scattered light is due to a diffuse interstellar medium, which the WP model attempts to represent. But the large grains which scatter starlight in the visible have a thermal emission, and, considering the importance of the scattering, must be detected on the IRAS images. Therefore, the diffuse medium which, according to GWF, is responsible for the DGL in fields (a) and (b) can be identified to the infrared cirrus shown in figure 4. The GWF areas are in approximate correspondence with the $`5^{}\times 5^{}`$ rectangles in each of the images. These areas sample a medium with evident structure. It is this medium that the WP model pretends to reconstruct from its HI emission. #### 5.1.2 IRAS images of fields (a) and (b) In all probability the mediums which compose the cirrus of field (b) (figure 4, right) has similar properties in each of the GWF areas. The decrease of the $`100\mu `$m surface brightness with absolute latitude in field (b) is associated with the decrease of the radiation field. At high absolute latitude, the IRAS surface brightness varies from $`5`$ MJy/sr to $`1.5`$ MJy/sr. With an $`100\mu `$m to visible extinction ratio $`I_{100}/A_V`$ of $`18`$ MJy/sr/mag (Boulanger & Pérault bou88 (1988)), $`A_V`$ is less than $`0.2`$ on the average. The areas of field (a) with $`b>45^{}`$ have similar properties as field (b). They have a low $`100\mu `$m surface brightness, a small $`A_V`$ (on average), and a clear small scale structure. These areas are the outermost parts of the HI loop at the edge of the Scorpio Centaurus region. The (a) areas at lower latitude are in the densest parts of the HI loop and have high surface brightnesses. Zeta Oph ($`m_V=2.6`$), a few degrees apart, may contribute to the heating of the region and to the enhancement, up to $`30`$ MJy/sr, of its infrared emission. In the low latitude areas there is an increase of the average HI column density (de Geus degeus88 (1988)), hence of the visible extinction. CO is detected at the brightest IRAS positions (Laureijs et al. laureijs95 (1995)). GWF data can be separated in two. The field (b) and the high latitude regions of field (a) have a low column density and a low visible extinction on average. The low latitude points of field (a) are clearly a different kind of medium with much higher column density and visible extinction. ### 5.2 The color of the ISGL In regions where there is interstellar matter, the direct starlight is reddened. In these regions the color of the direct starlight, $`isgl_R/isgl_B`$, is redder than $`\sigma _R/\sigma _B`$, where $`\sigma _R`$ and $`\sigma _B`$ are the surface brightnesses of direct starlight corrected for interstellar reddening. Let $`f`$ be the filling factor of interstellar matter in one of the areas GWF have considered. Assume that the ‘diffuse interstellar matter’ has an average low $`A_V\tau _V`$ value. This is true in region (b) and in the high latitude regions of (a). We have for the regions with low $`100\mu `$m emission: $`isgl_R`$ $`=`$ $`(1f)\sigma _R+(1\tau _R)f\sigma _R`$ (2) $`=`$ $`\sigma _R(1f\tau _R)`$ $`isgl_B`$ $`=`$ $`\sigma _B(1f\tau _B),`$ (3) and $`{\displaystyle \frac{isgl_R}{isgl_B}}`$ $`=`$ $`{\displaystyle \frac{\sigma _R}{\sigma _B}}{\displaystyle \frac{1f\tau _R}{1f\tau _B}}`$ (4) $``$ $`{\displaystyle \frac{\sigma _R}{\sigma _B}}(1+f(\tau _B\tau _R))`$ Expressed as a function of $`A_V`$ (Cardelli et al cardelli89 (1989)), equation 4 takes the simple form: $$\frac{isgl_R}{isgl_B}=\frac{\sigma _R}{\sigma _B}(1+0.6fA_V)$$ (5) The measured color of the stars is of course redder than $`\sigma _R/\sigma _B`$, but the change of color will not be important. The reddening of direct starlight, $`1+0.6fA_V`$, for a medium of $`I_{100}2`$ MJy/sr and $`A_V0.1`$ will be of order $`1+0.05f`$. Within the margin of error estimated by GWF for the ISGL, $`0.1`$, we can adopt $`isgl_R/isgl_B=\sigma _R/\sigma _B`$. Local increases of column density, due to the small scale structure of the medium, will not modify this approximation. The high resolution images of MCLD123.5+24.9 presented in Zagury et al. (zagury99 (1999)) shows the presence of high density clumps with a small surface coverage ($`5^{}\times 5^{}`$). In field (b) and for the (a)-field regions where $`b>40^{}`$, high density clumps must occupy only a small fraction of the surface. In such clumps the number of stars diminishes with $`A_V`$ and the color of ISGL in the lower column density medium, which probably occupies most of the volume in regions (a) and (b), will determine the ISGL color of the $`5^{}\times 5^{}`$ areas. The result does not hold for the low latitude points of field (a). There is a net increase of the average column density which can be evaluated by: $`{\displaystyle \frac{isgl_R}{isgl_B}}`$ $`=`$ $`{\displaystyle \frac{\sigma _R}{\sigma _B}}e^{\tau _B\tau _R}`$ (6) $`=`$ $`{\displaystyle \frac{\sigma _R}{\sigma _B}}e^{0.5A_V}`$ While the reddening is, as seen before, small for small $`A_V`$ values, for $`A_V=0.5`$ the ISGL is $`1.3`$ times redder than the stars’ color. The latter remarks also show that GWF’s estimate of the $`B`$ and $`R`$ magnitudes of the stars may have been biased when the $`A_V`$ value of the stars were used to construct the ‘Master Catalog’. To estimate $`A_V`$, GWF assume an average extinction of $`0.6`$ mag/kpc. This approximation applies only if the interstellar matter in the star direction has a low column density. In the case of larger column densities, the GWF estimate of the color of starlight will be bluer than it is in reality. The distance to the cirrus of field (a) should be $`200`$ pc at most, since it belongs to the Scorpio Centaurus region. For an $`A_V`$-value of 0.5 for instance, the extinction of all stars at closer distance than $`300`$ pc will be underestimated. The effect increases with column density and will be more pronounced in the low latitude regions of field (a). It will affect the $`B`$ band more than the $`R`$ band since the $`E(BV)/A_V`$ coefficient used by GWF is twice $`E(VR)/A_V`$. The color of the ISGL estimated by GWF from the Master catalog will be bluer than it should. It may explain the drop of the ISGL color of the lowest latitude point of field (a), figure 3. ### 5.3 The Pioneer color The Pioneer surface brightness comprises the ISGL ($`isgl\sigma `$) and the DGL. If the scattering volume is a medium of low optical detph $`\tau _B`$ in the $`B`$ band, the DGL surface brightness will be at most: $`\omega \tau _B`$, where $`\omega `$ is the albedo, assumed to be constant ($`0.6`$) at optical wavelengths. Use of equations 2 and 3 gives: $`P_R=\sigma _R(1f(1\omega )\tau _R)`$ (7) $`P_B=\sigma _B(1f(1\omega )\tau _B)`$ (8) Then: $`{\displaystyle \frac{P_R}{P_B}}`$ $`=`$ $`{\displaystyle \frac{\sigma _R}{\sigma _B}}{\displaystyle \frac{(1f(1\omega )\tau _R)}{(1f(1\omega )\tau _B)}}`$ (9) $``$ $`{\displaystyle \frac{\sigma _R}{\sigma _B}}(1+f(1\omega )(\tau _B\tau _R))`$ $`{\displaystyle \frac{P_R}{P_B}}`$ $`=`$ $`{\displaystyle \frac{isgl_R}{isgl_B}}(1\omega (\tau _B\tau _R))`$ (10) The Pioneer color is in between the ISGL and the DGL colors. Equation 9 can be simplified with $`\omega 0.6`$ and Cardelli, Clayton and Mathis (cardelli89 (1989)) relations: $`{\displaystyle \frac{P_R}{P_B}}`$ $`=`$ $`{\displaystyle \frac{\sigma _R}{\sigma _B}}(1+0.3fA_V)`$ (11) $`{\displaystyle \frac{P_R}{P_B}}`$ $`=`$ $`{\displaystyle \frac{isgl_R}{isgl_B}}(10.15fA_V)`$ (12) Equations 11 and 12 show that within GWF error margin equality between $`P_R/P_B`$, $`isgl_R/isgl_B`$, $`\sigma _R/\sigma _B`$ and $`dgl_R/dgl_B`$, will be satisfied for mediums of low column density. ### 5.4 Effect of the cirrus small scale structure on the observed colors Small scale structure affects the surface brightness of starlight. Stars with no interstellar matter on their line of sight will have little or no reddening, while stars behind a cirrus are reddened. The clumpiness of the regions GWF have chosen, revealed by the IRAS images, will certainly affect the evaluation of the starlight surface brightness since GWF have considered an average reddening by unit distance (section 2) for all stars, regardless of the increase of interstellar matter in some directions. As pointed out in section 5.2, the color of direct starlight may be redder than estimated by GWF. Small scale structure also modifies the color of the DGL since small clumps of different $`A_V`$ can be mixed in the beam. Even if the optical depth averaged over the beam of the observations is small, the existence of clumps of higher $`A_V`$ than average cannot be discarded. These clumps will redden the color of the DGL from the color of a low column density medium. Regions of higher than average $`A_V`$, such as the low latitude areas of field (a), are more likely to show this effect. These effects of the small scale structure of the interstellar medium on GWF data cannot be totally leaved over. Each of the GWF areas is large compared to the size at which the interstellar medium is stuctured. IRAS images show that this size is less than a few arcminutes. Higher resolution observations (Falgarone et al. falgarone (1998), Zagury et al. zagury99 (1999)) give sizes of a few arcseconds at most for the entities which compose the interstellar medium. ### 5.5 The comparison of Pioneer and ISGL colors The preceding sections show that a precise comparison of Pioneer and direct starlight data will sharply depend on the reliability of the GWF master catalog and on the cirrus structure. For the GWF regions of low average optical depth, plot 3 is in accordance, within the error margin given by GWF, with equality 1. Concerning the two lower latitude point of the field (a) different reasons may explain the drop in ISGL color. These points are different from the others since the region has a much higher $`A_V`$ and cannot be considered as a ‘diffuse medium’. This was not taken into account by GWF for the estimate of the direct starlight $`B`$ and $`R`$ magnitudes. The medium also contains regions with more extinction which may modify the DGL color and redden the Pioneer color. ## 6 A qualitative comparison of IRAS image and the GWF data A remarkable difference between fields (a) and (b) can be seen in figure 1. In field (b), $`isgl_R`$ and $`isgl_B`$ decrease with increasing (absolute) latitude, as is expected. The $`B`$ surface brightness follows a $`1/\mathrm{sin}\|b\|`$ law from $`b=28`$ to $`b=55`$, with a higher value than expected at $`b=22`$. It is slightly under but close to the prediction of the Besançon Galactic Model (figure 5). The $`R/B`$ color closely follows the model. In field (a) the ISGL is nearly constant and $`isgl_R`$ and $`isgl_B`$ have parallel variations. Low latitude points in (a) have lower ISGL than in (b) and than what is predicted ($`75\mathrm{S}_{10}(\mathrm{V})_{\mathrm{G2V}}`$) by the Besançon Galactic Model, figure 5, while higher latitude points ($`\|b\|>45^{}`$) have the same values at both longitudes. The ratio of the DGL to the total Pioneer emission in the $`R`$ band, $`S_R/P_R`$ (calculated from figure 1 and figure 2) is between $`0.45`$ and $`0.5`$ for the $`3`$ lowest latitude points in (a), while for the higher latitudes points and for all points in (b), it is between $`0.2`$ and $`0.3`$. The lower than expected starlight emission for field (a) low latitude points, along with the relative increase of DGL $`R`$ emission, are easily interpreted as extinction effects due to the average increase of interstellar matter along the line of sight (section 5.1). The increase of the column density increases the extinction of starlight. Corresponding to this extinction of starlight, there is a sharp rise in the DGL emission between the 3 (a) points at latitude $`42.5`$, $`37.5`$ and $`32.5`$, figure 2. The $`R`$ emission between latitudes $`37.5`$ and $`32.5`$ reaches a ceiling which corresponds to a decrease of $`B`$ emission, $`S_B`$. This can also be understood by the average increase of column density: absorption starts to dominate scattering in the $`B`$ band and the $`R`$ surface brightness is increased. For the two low latitude points of field (a) in figure 3 departure from relation 1 seems certain. In this region we are clearly outside the low column density approximation, supposed by GWF, hence outside the framework of their study. The understanding of the ISGL color in this region deserves further investigations but two reasons may contribute to the blue color of the ISGL. One reason is the GWF method to estimate the ISGL magnitude which applies to low column density mediums only (section 5.2). The second reason is if Zeta Oph, a few degrees apart, participates in the illumination of the region. Note that illumination by Zeta Oph ($`m_V=2.6`$) on such a large scale may be difficult to justify. The correlation between the IRAS images of both fields (figure 4) and the DGL $`R`$ emission (figure 2) makes it likely that most of the emission measured by Pioneer comes from the infrared cirrus. In the visible, the cirrus scatter the light of background stars. The relation between the color of the radiation field due to the stars and the color of the DGL is given by equation 1. The DGL optical emission at low latitude in field (a) may in part be due to illumination by the star Zeta Oph, a few degrees apart. ## 7 Conclusion GWF have compared the color of the diffuse galactic light at high galactic latitude, deduced from observations of the Pioneer satellite, to the color of the light scattered by a diffuse medium with a certain cloud size distribution. GWF conclude that ERE is required to explain the amount of diffuse galactic emission at high latitude and can represent up to $`50\%`$ of this emission. In GWF it is said that, “An accurate calculation of the DGL should include the effects of multiple scattering, the cloudiness of the interstellar medium, and the observed anisotropy of the illuminating radiation field.” It should also include a reliable description of the interstellar medium. It would be very surprising if the WP model or any other, as sophisticated as it might be, can deduce from the measure of the HI emission in one direction the organisation of the interstellar matter in that direction. It is also surprising that when the WP model does not match the data it is supposed to match, changes in the properties of the interstellar grains are invoked rather than the validity of the model. It seems more plausible that observations do not agree with the WP model because this model does not properly describe the reality of the interstellar medium. The exact structure of the interstellar medium the WP model should reproduce is shown at a scale of $`2^{}`$ by the IRAS $`100\mu `$m images. These images reveal the thermal emission of the same large grains responsible for the scattering of background starlight in the visible. They would have been a better basis to model the interstellar medium than the HI emission used in GWF. This medium may be of low column density on the average but is extremely structured and may have local column density enhancements. Some regions, the lower latitude points of field (a), have a clear increase of their average column density. It is very unlikely that the WP model reproduces this medium with the accuracy necessary to justify their conclusions. The uncertainty of the DGL color as calculated in GWF is extremely large compared to the errors on Pioneer integrated light color or on the direct starlight color. Comparison of the DGL and the starlight colors should be replaced by a comparison of the starlight to the Pioneer colors, both of which are determined with less relative error. Within the error margin given in GWF, there is equality between the color of the integrated starlight and the color given by the Pioneer measurements. It implies equality with the color of the DGL. There may be a tendency for the starlight surface brightness calculated by GWF to be slightly bluer than Pioneer integrated light. This tendency correlates with mediums of higher $`A_V`$ and can be attributed to two effects. Starlight is redder than estimated by GWF who assume a low and equal reddening in all directions. It modifies the GWF estimation of the $`R`$ and/or $`B`$ magnitude of the stars. Small scale structure will redden the color of DGL and of Pioneer from the color of a low column density medium. A better understanding of the optical emission in the high latitude directions comes from the consideration that the DGL in these directions is the light of background stars scattered by the nearby infrared cirrus. For the GWF fields (a) and (b), there are relations between the cirrus observed on the $`100\mu `$m IRAS image, the visible emission of background stars, and the DGL. The IRAS $`100\mu `$m emission in the GWF field (b) and the high latitude regions of field (a) attests to low column densities in the average. The corresponding DGL surface brightness is low. Increase of the column density of dust in the low latitude areas of field (a) attenuates starlight and increases the $`R`$ surface brightness of the DGL, while the increase in the $`B`$ surface brightness is limited because of absorption. ERE is not needed to explain the DGL at high galactic latitude (this paper) or in bright nebulae (Zagury, ‘Is there ERE in bright nebulae?’, submitted). Nor is it needed to explain the emission in the Red Rectangle (Zagury, in preparation), the milestone of ERE. There might be no ERE at all in interstellar space, which can be inferred from present day observations.
warning/0006/astro-ph0006404.html	ar5iv	text	# Stellar Evolution with Rotation V: Changes in all the Outputs of Massive Star Models ## 1 Introduction Due to the many well known successes of the standard theory of stellar evolution, rotation has generally been considered as a secondary effect. This was justified in many cases. However, since some years a number of very significant discrepancies between model predictions and observations have been found for massive stars, and also for red giants of lower masses. These difficulties have been listed and examined (cf. Maeder 1995a ). Recently, large excesses of \[N/H\] have been found in A–type supergiants, particularly in the SMC (Venn et al. Ve98 (1998); Venn Ve99 (1999)), where the excesses may reach up to an order of magnitude with respect to the average \[N/H\] local value in the SMC. These excesses, not predicted by current models, are the signatures of mixing effects over the entire star. This extensive mixing changes the size of the reservoir of nuclear fuel available for evolution, and thus the lifetimes, the tracks and all the model outputs (cf. also Langer 1992; Langer 1997; Heger et al. 2000). Also the chemical abundances, the yields and the final stages may be modified. Rotation appears as a natural driver for this mixing or, at least, as one of the first mechanisms whose consequences for mixing has to be explored. This is especially true for massive stars which are known to be fast rotators (see e.g. Penny 1996; Howarth et al. 1997). The physical effects of stellar rotation are numerous. The basic equations of stellar structure need to be modified. The meridional circulation and its interaction with the horizontal turbulence, the diffusion effects produced by shear turbulence, the inhibition of the transport mechanisms by the $`\mu `$–gradients, the transport of the angular momentum and of the chemical elements, the loss of mass and angular momentum at the surface, the enhancement of mass loss by rotation, etc… are some of the effects to be included in realistic models. In addition to the model physics, rotation also brings a number of new numerical problems in the stellar code, such as the 4<sup>th</sup> order equation for the transport of the angular momentum, which has to be coupled to the equations of stellar structure. In this work, we apply the investments in the treatment of the physical and numerical effects of rotation made in the previous works and we explore their consequences for the outputs of stellar models. In Sect. 2, we describe the various effects considered. The non–rotating stellar models are briefly discussed in Sect. 3. The evolution of the internal rotation during the evolution is examined in Sect. 4. Sect. 5 discusses the effects of rotation on the evolutionary tracks in the HR diagram, on the lifetimes and isochrones. The evolution of the rotational velocities at the stellar surface is discussed in Sect. 6 and the effects on the surface abundances are analysed in Sect. 7. ## 2 Physical ingredients of the models Let us briefly summarize here the basic physical ingredients of the numerical models of rotating stars we are constructing here. ### 2.1 Shellular rotation The differential rotation which results from the evolution and transport of the angular momentum as described by Eq. (3) below, makes the stellar interior highly turbulent. The turbulence is very anisotropic, with a much stronger geostrophic–like transport in the horizontal direction than in the vertical one (Zahn Za92 (1992)), where stabilisation is favoured by the stable temperature gradient. This strong horizontal transport is characterized by a diffusion coefficient $`D_\mathrm{h}`$, which is quite large as will be shown below. The horizontal turbulent coupling favours an essentially constant angular velocity $`\mathrm{\Omega }`$ on the isobars. This rotation law, constant on shells, applies to fast as well as to slow rotators. As an approximation, it is often represented by a law of the form $`\mathrm{\Omega }=\mathrm{\Omega }(r)`$ (Zahn Za92 (1992); see also Endal and Sofia ES76 (1976)). ### 2.2 Hydrostatic effects In a rotating star, the equations of stellar structure need to be modified (Kippenhahn and Thomas KippTh70 (1970)). The usual spherical coordinates must be replaced by new coordinates characterizing the equipotentials. The classical method applies when the effective gravity can be derived from a potential $`\mathrm{\Psi }=\mathrm{\Phi }\frac{1}{2}\mathrm{\Omega }^2r^2\mathrm{sin}^2\theta `$, i.e. when the problem is conservative. There, $`\mathrm{\Phi }`$ is the gravitational potential which in the Roche approximation is $`\mathrm{\Phi }=\frac{GM_r}{r}`$. If the rotation law is shellular, the problem is non–conservative. Most existing models of rotating stars apply, rather inconsistently, the classical scheme by Kippenhahn and Thomas. However, as shown by Meynet and Maeder (MM97 (1997)), the equations of stellar structure can still be written consistently, in term of a coordinate referring to the mass inside the isobaric surfaces. Thus, the problem of the stellar structure of a differentially rotating star with an angular velocity $`\mathrm{\Omega }=\mathrm{\Omega }(r)`$ can be kept one–dimensional. ### 2.3 Surface conditions The distribution of temperature at the surface of a rotating star is described by the von Zeipel theorem (vZ24 (1924)). Usually, this theorem applies to the conservative case and states that the local radiative flux $`F`$ is proportional to the local effective gravity $`g_{\mathrm{eff}}`$, which is the sum of the gravity and centrifugal force, $$F=\frac{L(P)}{4\pi GM_{}(P)}g_{\mathrm{eff}},$$ (1) with $`M_{}(P)=M(1\frac{\mathrm{\Omega }^2}{2\pi G\overline{\rho }})`$; $`L(P)`$ is the luminosity on an isobar and $`\overline{\rho }`$ the mean internal density. The local $`T_{\mathrm{eff}}`$ on the surface of a rotating star varies like $`T_{\mathrm{eff}}(\vartheta )g_{\mathrm{eff}}(\vartheta )^{\frac{1}{4}}`$. We define the average stellar $`T_{\mathrm{eff}}`$ by $`T_{\mathrm{eff}}^4=L/(\sigma S(\mathrm{\Omega }))`$, where $`\sigma `$ is Stefan’s constant and $`S(\mathrm{\Omega })`$ the total actual stellar surface. Of course, for different orientation angles $`i`$, the emergent luminosity, colours and spectrum will be different (Maeder and Peytremann MP70 (1970)). In the case of non–conservative rotation law, the corrections to the von Zeipel theorem depend on the opacity law and on the degree of differential rotation, but they are small, i.e. $`1\%`$ in current cases of shellular rotation (Kippenhahn Kipp77 (1977); Maeder Mae99 (1999)). Wether or not the star is close to the Eddington limit, the von Zeipel theorem keeps the same form as the one given by Maeder (1999, see also Maeder & Meynet 2000a ). In this last work, it is in particular shown that the expression for the Eddington factor in a rotating star needs to be consistently written and that it contains a term depending on rotation. ### 2.4 Changes of the mass loss rates $`\dot{M}`$ with rotation Observationally, a growth of the mass flux of OB stars with rotation, i.e. by 2–3 powers of 10, was found by Vardya (Var85 (1985)), while Nieuwenhuijzen and de Jager (Nieu88 (1988)) concluded that the $`\dot{M}`$–rates seem to increase only slightly with rotation for O– and B–type stars. On the theoretical side, Friend and Abbott (1986) find an increase of the $`\dot{M}`$–rates which can be fitted by the relation (Langer La98 (1998)) $$\dot{M}(v)=\dot{M}(v=0)\left(\frac{1}{1\frac{v}{v_{\mathrm{crit}}}}\right)^\xi $$ (2) with $`\xi 0.5`$ and $`v_{\mathrm{crit}}`$ the equatorial velocity at break–up. This expression is used in most evolutionary models and this is also what is done in the present work. The critical rotation velocity of a star is often written as $`v_{\mathrm{crit}}^2=\frac{GM}{R_{\mathrm{eb}}}(1\mathrm{\Gamma })`$, where $`R_{\mathrm{eb}}`$ is the equatorial radius at break–up velocity and $`\mathrm{\Gamma }=L/L_{\mathrm{Edd}}`$ is the ratio of the stellar luminosity to the Eddington luminosity (cf. Langer 1997, 1998). Glatzel (1998) has shown that when the effect of gravity darkening is taken into account, the above expression for $`v_{\mathrm{crit}}`$ does not apply. Glatzel (1998) gives $`v_{\mathrm{crit}}^2=\frac{GM}{R_{\mathrm{eb}}}`$, which we adopt here. The problem is now being further examined by Maeder & Meynet (2000a), who critically discuss the Eddington factors, their dependence on rotation, the expression of the critical velocity, the dependence of the mass loss rates on rotation. These various new results will be applied in subsequent works, particularly for the study of the effects of rotation on the formation of W–R stars. Further improvements, based either on the observations or on the theory, to account for the anisotropic winds which selectively remove the angular momentum need also to be performed. The above expression gives the change of the mass loss rates due to rotation. As reference mass loss rates in the case of no rotation, we use the recent data by Lamers and Cassinelli (LaCa96 (1996)); for the domain not covered by these authors we use the results by de Jager et al. (Ja88 (1988)). During the Wolf–Rayet phase we use the mass loss rates proposed by Nugis et al (1998) for the WNL stars (mean, clumping–corrected rates from radio data $`\dot{M}(\mathrm{WNL})=\mathrm{3\; 10}^5`$ M y<sup>-1</sup>). For the WNE and WC stars we use the prescription devised by Langer (1989), modified according to Schmutz (1997) for taking into account the clumping effects in Wolf-Rayet stellar winds ($`\dot{M}=2.410^8(M/\mathrm{M}_{})^{2.5}`$ M y<sup>-1</sup>). These mass loss rates are smaller by a factor 2–3 than the mass loss rates used in our previous stellar grids (Schaller et al. 1992; Meynet et al. 1994). ### 2.5 Transport of the angular momentum For shellular rotation, the equation of transport of angular momentum in the vertical direction is in lagrangian coordinates (cf. Zahn Za92 (1992); Maeder and Zahn MZ98 (1998)) $`\rho {\displaystyle \frac{d}{dt}}\left(r^2\mathrm{\Omega }\right)_{M_r}=`$ (3) $`{\displaystyle \frac{1}{5r^2}}{\displaystyle \frac{}{r}}\left(\rho r^4\mathrm{\Omega }U(r)\right)+{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{}{r}}\left(\rho Dr^4{\displaystyle \frac{\mathrm{\Omega }}{r}}\right).`$ $`\mathrm{\Omega }(r)`$ is the mean angular velocity at level $`r`$. The vertical component $`u(r,\theta )`$ of the velocity of the meridional circulation at a distance $`r`$ to the center and at a colatitude $`\theta `$ can be written $`u(r,\theta )=U(r)P_2(\mathrm{cos}\theta ),`$ (4) where $`P_2(\mathrm{cos}\theta )`$ is the second Legendre polynomial. Only the radial term $`U(r)`$ appears in Eq. (3). The quantity $`D`$ is the total diffusion coefficient representing the various instabilities considered and which transport the angular momentum, namely convection, semiconvection and shear turbulence. As a matter of fact, a very large diffusion coefficient as in convective regions implies a rotation law which is not far from solid body rotation. In this work, we take $`D=D_{\mathrm{shear}}`$ in radiative zones, since as extra–convective mixing we consider shear mixing and meridional circulation. In case the outward transport of the angular momentum by the shear is compensated by an inward transport due to the meridional circulation, we obtain the local conservation of the angular momentum. We call this solution the stationary solution. In this case, $`U(r)`$ is given by (cf. Zahn Za92 (1992)) $$U(r)=\frac{5D}{\mathrm{\Omega }}\frac{\mathrm{\Omega }}{r}.$$ (5) The full solution of Eq. (3) taking into account $`U(r)`$ and $`D`$ gives the non–stationary solution of the problem. In this case, $`\mathrm{\Omega }(r)`$ evolves as a result of the various transport processes, according to their appropriate timescales, and in turn differential rotation influences the various above processes. This produces a feedback and, thus, a self–consistent solution for the evolution of $`\mathrm{\Omega }(r)`$ has to be found. The transport of angular momentum by circulation has often been treated as a diffusion process (Endal and Sofia ES76 (1976); Pinsonneault et al. Pin89 (1989); Heger et al. He00 (2000)). From Eq. (3), we see that the term with $`U`$ (advection) is functionally not the same as the term with $`D`$ (diffusion). Physically advection and diffusion are quite different: diffusion brings a quantity from where there is a lot to other places where there is little. This is not necessarily the case for advection. A circulation with a positive value of $`U(r)`$, i.e. rising along the polar axis and descending at the equator, is as a matter of fact making an inward transport of angular momentum. Thus, we see that when this process is treated as a diffusion, like a function of $`\frac{\mathrm{\Omega }}{r}`$, even the sign of the effect may be wrong. The expression of $`U(r)`$ given below (Eq. 12) involves derivatives up to the third order, thus Eq. (3) is of the fourth order, which makes the system very difficult to solve numerically. In practice, we have applied a Henyey scheme to make the calculations. Eq. (3) also implies four boundary conditions. At the stellar surface, we take (cf. Talon et al. Ta97 (1997); Denissenkov et al. Den99 (1999)) $`{\displaystyle \frac{\mathrm{\Omega }}{r}}=0\mathrm{and}U(r)=0`$ (6) and at the edge of the core we have $`{\displaystyle \frac{\mathrm{\Omega }}{r}}=0\mathrm{and}\mathrm{\Omega }(r)=\mathrm{\Omega }_{\mathrm{core}.}`$ (7) We assume that the mass lost by stellar winds is just embarking its own angular momentum. This means that we ignore any possible magnetic coupling, as it occurs in low mass stars. This is not unreasonable in view of the negative results about the detection of magnetic fields in massive stars (Mathys Math99 (1999)). It is interesting to mention here, that in case of no viscous, nor magnetic coupling at the stellar surface, i.e. with the boundary conditions (6), the integration of Eq. (3) gives for an external shell of mass $`\mathrm{\Delta }M`$ (Maeder 1999, paper IV) $`\mathrm{\Delta }M{\displaystyle \frac{d}{dt}}(\mathrm{\Omega }r^2)={\displaystyle \frac{4\pi }{5}}\rho r^4\mathrm{\Omega }U(r).`$ (8) This equation is valid provided the stellar winds are spherically symmetric (see paper IV), an assumption we do in this work. When the surface velocity approches the critical velocity, it is likely that there are anisotropies of the mass loss rates (polar ejection or formation of an equatorial ring) and thus the surface condition should be modified according to the prescriptions of Maeder (Mae99 (1999)). For now, these effects are not included in these models. Their neglect should not affect too much the results presented here since the critical velocity is reached only in some rare circumstances. ### 2.6 Mixing and transport of the chemical elements A diffusion–advection equation like Eq. (3) should normally be used to express the transport of chemical elements. However, if the horizontal component of the turbulent diffusion $`D_\mathrm{h}`$ is large, the vertical advection of the elements can be treated as a simple diffusion (Chaboyer and Zahn Cha92 (1992)) with a diffusion coefficient $`D_{\mathrm{eff}}`$. As emphasized by Chaboyer and Zahn, this does not apply to the transport of the angular momentum. $`D_{\mathrm{eff}}`$ is given by $$D_{\mathrm{eff}}=\frac{rU(r)^2}{30D_h},$$ (9) where $`D_\mathrm{h}`$ is the coefficient of horizontal turbulence, for which the estimate is $$D_\mathrm{h}\|rU(r)\|$$ (10) according to Zahn (1992). Eq. (8) expresses that the vertical advection of chemical elements is severely inhibited by the strong horizontal turbulence characterized by $`D_\mathrm{h}`$. Thus, the change of the mass fraction $`X_i`$ of the chemical species $`i`$ is simply $`\left({\displaystyle \frac{dX_i}{dt}}\right)_{M_r}=`$ (11) $`\left({\displaystyle \frac{}{M_r}}\right)_t\left[(4\pi r^2\rho )^2D_{\mathrm{mix}}\left({\displaystyle \frac{X_i}{M_r}}\right)_t\right]+\left({\displaystyle \frac{dX_i}{dt}}\right)_{\mathrm{nucl}}.`$ The second term on the right accounts for composition changes due to nuclear reactions. The coefficient $`D_{\mathrm{mix}}`$ is the sum $`D_{\mathrm{mix}}=D_{\mathrm{shear}}+D_{\mathrm{eff}}`$ and $`D_{\mathrm{eff}}`$ is given by Eq. (9). The characteristic time for the mixing of chemical elements is therefore $`t_{\mathrm{mix}}\frac{R^2}{D_{\mathrm{mix}}}`$ and is not given by $`t_{\mathrm{circ}}\frac{R}{U}`$, as has been generally considered (Schwarzschild 1958). This makes the mixing of the chemical elements much slower, since $`D_{\mathrm{eff}}`$ is very much reduced. In this context, we recall that several authors have reduced by large factors, up to 30 or 100, the coefficient for the transport of the chemical elements, with respect to the transport of the angular momentum, in order to better fit the observed surface compositions (cf. Heger et al. He00 (2000)). This reduction of the diffusion of the chemical elements is no longer necessary with the more appropriate expression of $`D_{\mathrm{eff}}`$ given here. When the effects of the shear and of the meridional circulation compensate each other for the transport of the angular momentum (stationary solution, see Sect. 2.5), the value of $`U`$ entering the expression for $`D_{\mathrm{eff}}`$ is given by Eq. (5). ### 2.7 Meridional circulation Meridional circulation is an essential mixing mechanism in rotating stars and there is a considerable litterature on the subject (see ref. in Tassoul Tass90 (1990)). The velocity of the meridional circulation in the case of shellular rotation was derived by Zahn (Za92 (1992)). The effects of the vertical $`\mu `$–gradient $`_\mu `$ and of the horizontal turbulence on meridional circulation are very important and they were taken into account by Maeder and Zahn (MZ98 (1998)). Contrarily to the conclusions of previous works (e.g. Mestel Mes65 (1965); Kippenhahn and Weigert KippW90 (1990); Vauclair Vau99 (1999)), the $`\mu `$–gradients were shown not to introduce a velocity threshold for the occurence of the meridional circulation, but to progressively reduce the circulation when $`_\mu `$ increases. The expression by Maeder and Zahn (MZ98 (1998)) is $`U(r)={\displaystyle \frac{P}{\rho gC_PT[_{\mathrm{ad}}+(\phi /\delta )_\mu ]}}`$ (12) $`\left\{{\displaystyle \frac{L}{M_{}}}(E_\mathrm{\Omega }+E_\mu )\right\}.`$ $`P`$ is the pressure, $`C_P`$ the specific heat, $`E_\mathrm{\Omega }`$ and $`E_\mu `$ are terms depending on the $`\mathrm{\Omega }`$– and $`\mu `$–distributions respectively, up to the third order derivatives and on various thermodynamic quantities (see details in Maeder and Zahn, MZ98 (1998)). The term $`_\mu `$ in Eq. (12) results from the vertical chemical gradient and from the coupling between the horizontal and vertical $`\mu `$–gradients due to the horizontal turbulence. This term $`_\mu `$ may be one or two orders of magnitude larger than $`_{\mathrm{ad}}`$ in some layers, so that $`U(r)`$ may be reduced by the same ratio. This is one of the important differences introduced by the work by Maeder and Zahn (MZ98 (1998)). Another difference is that the classical solution usually predicts an infinite velocity at the interface between a radiative and a semiconvective zone with an inverse circulation in the semiconvective zone. Expression (12) gives a continuity of the solution with no change of sign from semiconvective to radiative regions. Finally, we recall that in a stationary situation, $`U(r)`$ is given by Eq. (5), as seen above. ### 2.8 Shear turbulence and mixing In a radiative zone, shear due to differential rotation is likely to be a most efficient mixing process. Indeed shear instability grows on a dynamical timescale that is of the order of the rotation period (Zahn Za92 (1992)). The usual criterion for shear instability is the Richardson criterion, which compares the balance between the restoring force of the density gradient and the excess energy present in the differentially rotating layers, $$Ri=\frac{N_{\mathrm{ad}}^2}{(0.8836\mathrm{\Omega }\frac{d\mathrm{ln}\mathrm{\Omega }}{d\mathrm{ln}r})^2}<\frac{1}{4},$$ (13) where we have taken the average over an isobar, $`r`$ is the radius and $`N_{\mathrm{ad}}`$ the Brunt-Väisälä frequency given by $`N_{\mathrm{ad}}^2={\displaystyle \frac{g\delta }{H_P}}\left[{\displaystyle \frac{\phi }{\delta }}_\mu +_{ad}_{\mathrm{rad}}\right].`$ (14) When thermal dissipation is significant, the restoring force of buoyancy is reduced and the instability occurs more easily, its timescale is however longer, being the thermal timescale. This case is referred to as “secular shear instability”. The criterion for low Peclet numbers $`Pe`$ (i.e. of large thermal dissipation, see below) has been considered by Zahn (Za74 (1974)), while the cases of general Peclet numbers $`Pe`$ have been considered by Maeder (1995b ), Maeder and Meynet (MM96 (1996)), who give $`Ri={\displaystyle \frac{g\delta }{(0.8836\mathrm{\Omega }\frac{d\mathrm{ln}\mathrm{\Omega }}{d\mathrm{ln}r})^2H_P}}`$ (15) $`\left[{\displaystyle \frac{\mathrm{\Gamma }}{\mathrm{\Gamma }+1}}(_{ad})+{\displaystyle \frac{\phi }{\delta }}_\mu \right]<{\displaystyle \frac{1}{4}}`$ The quantity $`\mathrm{\Gamma }=Pe/6`$, where the Peclet number $`Pe`$ is the ratio of the thermal cooling time to the dynamical time, i.e. $`Pe=\frac{v\mathrm{}}{K}`$ where $`v`$ and $`\mathrm{}`$ are the characteristic velocity and length scales, and $`K=(4acT^3)/(3C_P\kappa \rho ^2)`$ is the thermal diffusivity. A discussion of shear–driven turbulence by Canuto (Ca98 (1998)) suggests that the limiting $`Ri`$ number may be larger than $`\frac{1}{4}`$. To account for shear transport and diffusion in Eqs. (3) and (11), we need a diffusion coefficient. Amazingly, a great variety of coefficients $`D_{\mathrm{shear}}=\frac{1}{3}v\mathrm{}`$ have been derived and applied, for example: – 1. Endal and Sofia (ES78 (1978)) apply the Reynolds and the Richardson criterion by Zahn (Za74 (1974)). They estimate $`D_{\mathrm{shear}}`$ from the product of the velocity scale height of the shear flow and of the turbulent velocities of cells at the edge of Reynolds critical number. – 2. Pinsonneault et al. (Pin89 (1989)) notice that the amount of differential rotation permitted by the secular shears is proportional to a critical number, which they treat as an adjustable parameter. To account for the effects of the $`\mu `$–gradient in mixing, of the loss of angular momentum, etc… they introduce several adjustable parameters in the equations for the transport of the chemical elements and of the angular momentum. – 3. Chaboyer et al. (1995a ) use a coefficient derived from the velocity and path length from Zahn (Za74 (1974)). Following Pinsonneault et al. (Pin89 (1989)), they also introduce two adjustable parameters for adjusting the transports of chemical elements and angular momentum respectively. We notice that thanks to the reduction of the diffusion produced by the horizontal turbulence (cf. Eq. 9 above), it is no longer necessary to arbitrarily reduce the vertical transport of the chemical elements. – 4. Zahn (Za92 (1992)) defines the diffusion coefficient corresponding to the eddies which have the largest $`Pe`$ number so that the Richardson criterion is just marginally satisfied. However, the effects of the vertical $`\mu `$–gradient are not accounted for and the expression only applies to low Peclet numbers. – 5. The same has been done by Maeder and Meynet (MM96 (1996)), who considers also the effect of the vertical $`\mu `$–gradient, the case of general Peclet numbers and, in addition they account for the coupling due to the fact that the shear also modifies the local thermal gradient. This coefficient has been used by Meynet and Maeder (MM97 (1997)) and by Denissenkov et al. (Den99 (1999)). – 6. The comparisons of model results and observations of surface abundances have led many authors to conclude that the $`\mu `$–gradients appear to inhibit the shear mixing too much with respect to what is required by the observations (Chaboyer et al. 1995a b; Meynet and Maeder MM97 (1997); Heger et al. He00 (2000)). Namely, the observations of O–type stars (Herrero et al. He92 (1992), He99 (1999)) show much more He– and N–enrichments than predicted by the models which apply Richardon’s criterion with the $`_\mu `$ term. Thus, instead of using a gradient $`_\mu `$ in the criterion for shear mixing, Chaboyer et al. (1995a ) and Heger et al. (He00 (2000)) write $`f_\mu _\mu `$ with a factor $`f_\mu =0.05`$ or even smaller. This procedure is not satisfactory since it only accounts for a small fraction of the existing $`\mu `$–gradients in stars. The problem is that the models depend at least as much (if not more) on $`f_\mu `$ than on rotation, i.e. a change of $`f_\mu `$ in the allowed range (between 0 and 1) produces as important effects as a change of the initial rotational velocity. This situation has led to two other more physical approaches discussed below. Also Heger et al. (2000) introduce another factor $`f_c`$ to adjust the ratio of the transport of the angular momentum and of the chemical elements like Pinsonneault et al. (1989). – 7. Talon and Zahn (TZ97 (1997)) account for the horizontal turbulence, which has a coefficient $`D_\mathrm{h}`$ and which weakens the restoring force of the gradient of $`\mu `$ in the usual Richardson criterion. This allows some mixing of the chemical elements to occur as required by the observations (Talon et al. Ta97 (1997)). – 8. Indeed, around the convective core in the region where the $`\mu `$–gradient inhibits mixing, there is anyway some turbulence due to both the horizontal turbulence and to the semiconvective instability, which is generally present in massive stars. This situation has led to the hypothesis (Maeder Mae97 (1997)) that the excess energy in the shear, or a fraction $`\alpha `$ of it of the order of unity, is degraded by turbulence on the local thermal timescale. This progressively changes the entropy gradient and consequently the $`\mu `$–gradient. This hypothesis leads to a diffusion coefficient $`D_{\mathrm{shear}}`$ given by $`D_{\mathrm{shear}}=4{\displaystyle \frac{K}{N_{\mathrm{ad}}^2}}\left[{\displaystyle \frac{1}{4}}\alpha \left(0.8836\mathrm{\Omega }{\displaystyle \frac{d\mathrm{ln}\mathrm{\Omega }}{d\mathrm{ln}r}}\right)^2(^{})\right].`$ (16) The term $`^{}`$ in Eq. (16) expresses either the stabilizing effect of the thermal gradients in radiative zones or its destabilizing effect in semiconvective zones (if any). When the shear is negligible, $`D_{\mathrm{shear}}`$ tends towards the diffusion coefficient for semiconvection by Langer et al. (La83 (1983)) in semiconvective zones. When the thermal losses are large ($`^{}=`$), it tends towards the value $$D_{\mathrm{shear}}=\alpha (K/N_{\mathrm{ad}}^2)\left(0.8836\mathrm{\Omega }\frac{d\mathrm{ln}\mathrm{\Omega }}{d\mathrm{ln}r}\right)^2,$$ (17) given by Zahn (Za92 (1992)). Eq. (16) is completed by the three following equations expressing the thermal effects (Maeder Mae97 (1997)) $`D_{\mathrm{shear}}=2K\mathrm{\Gamma }={\displaystyle \frac{_{rad}+(\frac{6\mathrm{\Gamma }^2}{1+\mathrm{\Gamma }})_{ad}}{1+(\frac{6\mathrm{\Gamma }^2}{1+\mathrm{\Gamma }})}},`$ (18) $`^{}={\displaystyle \frac{\mathrm{\Gamma }}{\mathrm{\Gamma }+1}}(_{\mathrm{ad}}).`$ (19) The system of 4 equations given by Eqs. (16), (18) and (19) form a coupled system with 4 unknown quantities $`D`$, $`\mathrm{\Gamma }`$, $``$ and $`^{}`$. The system is of the third degree in $`\mathrm{\Gamma }`$. When it is solved numerically, we find that as a matter of fact the thermal losses in the shears are rather large in massive stars and thus that the Peclet number $`Pe`$ is very small (of the order of 10<sup>-3</sup> to 10<sup>-4</sup>, see Sect. 4.2). For very low Peclet number $`Pe=6\mathrm{\Gamma }`$, the differences $`(^{})`$ are also very small as shown by Eq. (19). Thus, we conclude that Eq. (16) is essentially equivalent, at least in massive stars, to the original Eq. (17) above, as given by Zahn (Za92 (1992)). We may suspect that this is not necessarily true in low and intermediate mass stars since there the $`Pe`$ number may be larger. Of course, the Reynolds condition $`D_{\mathrm{shear}}\frac{1}{3}\nu Re_c`$ must be satisfied in order that the medium is turbulent. The quantity $`\nu `$ is the total viscosity (radiative + molecular) and $`Re_c`$ the critical Reynolds number estimated to be around 10 (cf. Denissenkov et al. Den99 (1999); Zahn Za92 (1992)). The numerical results in Sect. 4 will show the values of the various parameters and also indicate that the conditions for the occurence of turbulence are satisfied. The physical treatment around the core also depends on the choice of the criterion for convection. In current literature, there are at least three basic sets of assumptions: a) Ledoux criterion, which leads to small cores, b) Schwarzschild’s criterion, which gives “medium size” cores, c) Schwarzschild’s criterion and overshooting, which gives large cores. Here, we choose to apply the intermediate solution b), which means that the above equation 16 is only applied in fully radiative regions. Some consequences of this choice of convective criterion are discussed in Sect. 3 below. ### 2.9 Initial compositions, opacities, nuclear reactions and other model ingredients For purpose of comparison, we adopted here the same physical ingredients as for the solar metallicity models computed by Meynet et al. (me94 (1994)). The only exceptions, apart from the inclusion of the effects induced by rotation described above, are the following: we use the Schwarzschild criterion for convection without overshooting and the mass loss rates are as indicated in Sect. 2.4. ## 3 The non–rotating stellar models For comparison purpose, we have computed non–rotating stellar models with exactly the same physical ingredients as the rotating ones. The corresponding evolutionary tracks and lifetimes are presented in Fig. 8 and Table 1. These stellar models have similar properties as older ones computed with the Schwarzschild criterion for convection. For instance, our non–rotating 9 M stellar model has similar H– and He–burning lifetimes as the 9 M model computed with the Schwarzschild criterion by Bertelli et al. (1985). Our models also well agree with more recent computations. Indeed our MS lifetimes are similar within about two percents to the ones obtained by Heger et al. (2000) for their non–rotating stellar models. These authors used the Ledoux criterion with semiconvection. However, since during the MS phase the convective core mass decreases, one expects, for this phase, only small differences between the models computed with the Ledoux and the Schwarzschild criterion. These comparisons show that the present non–rotating stellar models well agree with results obtained with different stellar evolutionary codes. Many papers (e.g. Bertelli et al. 1985; Maeder & Meynet 1989; Chin & Stothers 1990; Langer & Maeder 1995; Canuto 2000) have discussed the effects of different criteria for convection on massive star evolution. Since our main aim in this paper is to emphasize the effects of rotation, we shall restrict ourselves to briefly mention some differences with previous grids of models computed by the Geneva group. In the present non–rotating stellar models, computed without overshooting, the convective core masses are smaller than in the models by Schaller et al. (1992) and Meynet et al. (1994). As a consequence, the stars evolve at smaller luminosities. The turn–off point at the end of the MS occurs for higher T<sub>eff</sub> and lower luminosities, reducing the extent of the MS width. The MS lifetimes are reduced. For initial masses below about 25 M, the He–burning lifetime is increased implying an increase of the ratio $`t_{\mathrm{He}}/t_\mathrm{H}`$ of the lifetimes in the He– and H–burning phases. For the higher initial masses, the effects of the stellar winds become very important and dominate the effects due to a change of the criterion for convection. As a numerical example, in our non–rotating 20 M stellar model, the convective core mass is reduced by 1–1.2 M during the whole H–burning phase compared to its value given by Meynet et al. (1994). At the end of the MS phase, the mass of the helium core has decreased by $``$ 20% with respect to its value in models with a moderate overshoot. The position of the turn–off point has a $`\mathrm{log}L`$/L decreased by 0.05 dex and a $`\mathrm{log}T_{\mathrm{eff}}`$ increased by 0.04 dex. The MS lifetime is decreased by about 10%, the ratio $`t_{\mathrm{He}}/t_\mathrm{H}`$ passes from 10% in the models of Meynet et al. (1994) to 14% in the present case. ## 4 The evolution of the internal rotation law $`\mathrm{\Omega }(r)`$ ### 4.1 The initial convergence of the rotation law Let us first mention that even in some recent works the assumption of solid body rotation is often considered. However, it is much more physical to examine the evolution of $`\mathrm{\Omega }(r)`$ resulting from the transport of angular momentum by shears, meridional circulation and convection, from the effects of central contraction and envelope expansion and from the losses of angular momentum by stellar winds. In particular, the large losses of angular momentum at the surface lead to a redistribution of the angular momentum in the interior by meridional circulation, shear turbulent diffusion and convection. The whole problem is self–consistent, since the meridional circulation and shear transport depend in turn on the degree of differential rotation.Thus, the full system of equations has to be solved in order to provide the internal evolution of $`\mathrm{\Omega }`$. Fig. 1 shows the results for the initial convergence on the zero age main sequence of $`\mathrm{\Omega }(r)`$ in a 20 M model with solar composition. The initial model on the zero–age main sequence is supposed to rotate uniformly and with an initial surface velocity $`v_{\mathrm{ini}}`$ of 300 km s<sup>-1</sup>. Very short time steps of the order of 1920 yr are taken in the initial calculations for obtaining the internal equilibrium rotation. The results are shown in Fig. 1 for every 10<sup>th</sup> model, i.e. with time intervals of 19 200 yr. We see the very fast initial changes of $`\mathrm{\Omega }(r)`$, mainly due to the meridional circulation. The circulation velocity, which is very large, is also positive everywhere at the beginning (cf. Fig. 2). This means that the circulation rises along the polar axis and goes down at the equator, thus bringing angular momentum inwards. As a consequence the angular velocity of the core increases, while that in the envelope goes down. It is important to realize that here these fast changes of the rotation law are not due, as later in the evolution, to core contraction and envelope expansion. We see that the changes are big at the beginning and then smaller. The profiles of $`\mathrm{\Omega }(r)`$ rapidly converges towards an equilibrium–profile in a time of about 1 to 2% of the MS lifetime $`t_{\mathrm{MS}}`$. This is in agreement with the results for 10 and 30 M stars by Denissenkov et al. (Den99 (1999)). As noted by these authors, the timescale for the adjustment of the rotation law, i.e. $`t_{\mathrm{circ}}R/U`$, is short with respect to the MS lifetime for any significant rotation. Later on during the bulk of MS evolution, the profiles $`\mathrm{\Omega }(r)`$ will evolve more slowly. As emphasized in Sect. 2.6, this timescale $`t_{\mathrm{circ}}`$ is not characteristic for the mixing of chemical elements. We also notice that the degree of differential rotation in the equilibrium profile is rather modest (cf. Zahn, Za92 (1992)). Fig. 2 shows the corresponding initial evolution of the vertical component $`U(r)`$ of the velocity of the meridional circulation in the 20 M model. As seen above, initially this velocity is large and positive everywhere, which explains the fast evolution of $`\mathrm{\Omega }(r)`$ in Fig. 1. Then, $`U(r)`$ decreases and becomes negative in the very external layers. The physical reason for that is the so called Gratton–Öpik term of the form $`\frac{\mathrm{\Omega }^2}{2\pi G\rho }`$ which is contained in the expression of $`E_\mathrm{\Omega }`$ in Eq. (12). When the density is very low, as in the outer regions, this negative term becomes important and produces an inverse circulation. This means that the circulation has two big cells, an internal cell rising along the polar axis and an external one descending at the pole. The velocity $`U(r)`$ also converges towards an equilibrium distribution, characterized by small velocities $`U(r)`$, a result in agreement with Denissenkov et al. (Den99 (1999)). It is interesting to compare the stationary solution of $`U(r)`$ given by Eq. (5) and the non–stationary solution given by Eq. (12). This is done in Fig. 3, where the curve in continuous line corresponds to the asymptotic distribution reached in the non–stationary regime illustrated in Fig. 2. This curve shows positive values of $`U(r)`$ in the inner regions and negative values in the outer ones. On the contrary, the solution (broken line) obtained in the stationary approximation given by Eq. (5) is always positive. This is a logical consequence of the approximation made: as $`\mathrm{\Omega }(r)`$ decreases outwards, only positive values of $`U(r)`$ are possible. Said in other words, the outward transport by shears can only be compensated by the inward transport by meridional circulation, i.e. by positive values of $`U(r)`$. However, in non–stationary models, this compensation is not achieved since the two curves in Fig. 3 are very different. One concludes that the stationary approximation which has been used in some works (cf. Urpin et al. Urp96 (1996)) is not satisfactory. The stationary solution can also be viewed as giving a velocity which is the inverse of a velocity which could be associated to the diffusive transport by shears. From the comparison of the two curves in Fig. 3, one thus immediately concludes that the meridional circulation is much more efficient for the transport of the angular momentum than the shear instabilities. The justification for this claim is the following one. In the stationary situation, shear transport and meridional circulation compensate each other and are thus of the same magnitude. The fact that $`U(r)`$ in the full non–stationary case treatment is much larger than in the stationary case as shown in Fig. 3 implies that the effects of meridional circulation on the transport of angular momentum are much larger than the effects of the shear. This conclusion can also be deduced from the comparison of the timescales for the two processes. In general one has that $`t_{\mathrm{circ}}=\frac{R}{U}<t_{\mathrm{shear}}=\frac{R^2}{D_{\mathrm{shear}}}`$. For the chemical elements, the transport by the shear instabilities is more efficient than the transport by the meridional circulation as discussed in Sect. 2.6 where we have seen that $`D_{\mathrm{eff}}`$ is reduced by horizontal turbulence. ### 4.2 Internal evolution of $`\mathrm{\Omega }(r)`$ The evolution of $`\mathrm{\Omega }(r)`$ during the MS evolution of a 20 M star is shown in Fig. 4. We notice that initially (i.e. for high values of the central H–content $`X_\mathrm{c}`$), the degree of differential rotation is small, then, it grows during the evolution. The ratio between the central and surface value of $`\mathrm{\Omega }`$ remains however small during the MS evolution. The rotation rate does not vary by more than about a factor two throughout the star. Although this degree of differential rotation is not large, it plays an essential role in the shear effects and the related transports of angular momentum and of chemical elements. The same remark also applies to the effects of differential rotation on the transport by meridional circulation, since $`U(r)`$ critically depends on the derivatives of $`\mathrm{\Omega }`$ in the star (cf. Zahn 1992; Maeder & Zahn 1998). We notice a general decrease of $`\mathrm{\Omega }(r)`$, even in the convective core which is contracting. The main reason is mass loss at the stellar surface, which removes a substantial fraction of the total angular momentum. This makes $`\mathrm{\Omega }(r)`$ to decrease with time everywhere, because of the internal transport mechanisms, which ensure some coupling of rotation. The same behaviour was obtained by Langer (1998) in the case of rigid rotation (i.e. in the case of maximum coupling). In the present models, shear transports the angular momentum outwards, which reduces the core rotation, while circulation makes an inward transport in the deep interior and an outwards transport in the external region. This is responsible for the progressive flattening of the curves above $`r=4R_{}`$ in Fig. 4. From the end of the MS evolution onwards, i.e. when $`X_c0.05`$, central contraction becomes faster and starts dominating the evolution of the angular momentum in the center. The central $`\mathrm{\Omega }`$ grows quickly (cf. dotted curve in Fig. 4) and this will in principle continue in the further evolutionary phases until core collapse. The evolution of the angular momentum after the H–burning phase ($`X_c=0`$) is mainly dominated by the local conservation of the angular momentum, since the secular transport mechanisms have little time to operate. This is the assumption we are making in the present models. During these phases, the chemical elements continue to be rotationally mixed in the radiative layers mainly through the effect of the shear. In the further evolutionary phases, the strong central contraction will lead to very large central rotation, unless some other processes as fast dynamical instabilities or magnetic braking occur in these stages. The evolution of $`U(r)`$ during MS evolution is shown in Fig. 5. The outer zone with inverse circulation progressively deepens in radius during MS evolution due to the Gratton–Öpik term, because a growing part of the outer layers has lower densities. Also, we notice that the velocities $`U(r)`$ are small in general (cf. Urpin et al. Urp96 (1996)) and tests have shown us that, contrary to the classical result of the Eddington–Sweet circulation (cf. Mestel 1965), $`U(r)`$ depends rather little on the initial rotation. The deepening of the inverse circulation has the consequence that the stationary and non–stationary solutions differ more and more as the evolution proceeds, since as said above no inverse circulation is predicted by the stationary solution. Thus we conclude that the stationary solutions are much too simplified. Also, in Fig. 5, we see that the values of $`U(r)`$ become more negative in the outer layers for the model at the end of the MS phase, when central contraction occurs. Again, this effect is due to the Gratton–Öpik term. The results of Fig. 5 also shows that the application of a diffusive treatment to meridional circulation transport is unappropriate. Paradoxically, the problem is less serious in regions which exhibit an inverse circulation, since there the signs of the effects are at least the same in the two treatments (however even there the equations for advection and diffusion are different). The various diffusion coefficients inside a 20 M star are shown in Fig. 6, when the hydrogen mass fraction $`X_c`$ at the center is equal to 0.20. We notice that in general $`K>D_h>D_{\mathrm{shear}}>D_{\mathrm{eff}}>\nu `$. Following the thermal diffusivity $`K`$, the coefficient of horizontal diffusion $`D_h`$ is the second largest one, this is consistent and necessary for the validity of the assumption of shellular rotation (Zahn, Za92 (1992)). Since we have $`D_{\mathrm{shear}}=2K\mathrm{\Gamma }`$ according to Eq. (18), Fig. 6 shows that $`\mathrm{\Gamma }`$ is typically of the order of 10<sup>-3</sup> to 10<sup>-4</sup>. The coefficient $`D_{\mathrm{eff}}`$, expressing the transport of the chemical elements by the meridional circulation with account of the reduction by the horizontal turbulence, is only slightly smaller (by a factor of about 3) than $`D_{\mathrm{shear}}`$ in the interior. One must be careful that $`D_{\mathrm{eff}}`$ concerns only the tranport of the elements and not that of angular momentum, for which meridional circulation is the largest effect. Fig. 6 also shows the values of the total viscosity $`\nu `$, and we notice that consistently the Prandtl number $`\nu /K`$ is of the order of 10<sup>-5</sup> to 10<sup>-6</sup>. ## 5 HR diagram, lifetimes and age estimates ### 5.1 The Main–Sequence evolution Evolutionary tracks of 20 M models at solar metallicity for different initial velocities are plotted on Fig. 7. The $`T_{\mathrm{eff}}`$ for a rotating star has been defined in Sect. 2.3. On and near the ZAMS, rotation produces a small shift of the tracks towards lower luminosities and $`T_{\mathrm{eff}}`$. This effect is due to both atmospheric distorsions and to the lowering of the effective gravity (see e.g. Kippenhahn and Thomas KippTh70 (1970); Maeder and Peytremann MP70 (1970); Collins and Sonneborn co77 (1977)). At this stage the star is still nearly homogeneous. When evolution proceeds, the tracks with rotation become more luminous than for no rotation. This results from essentially two effects. On one side, rotational mixing brings fresh H–fuel into the convective core, slowing down its decrease in mass during the MS. For a given value of the central H–content, the mass of the convective core in the rotating model is therefore larger than in the non–rotating one and thus the stellar luminosity is higher (Maeder ma87 (1987); Talon et al. Ta97 (1997); Heger et al. He00 (2000)). As a numerical example, in the $`v_{\mathrm{ini}}=`$ 300 km s<sup>-1</sup> models the He–cores at the end of the MS are about 20% more massive than in their non–rotating counterparts. This is equal to the increase obtained by a moderate overshooting (see Sect. 3). On the other side, rotational mixing transports helium and other H–burning products (essentially nitrogen) into the radiative envelope. The He–enrichment lowers the opacity. This contributes to the enhancement of the stellar luminosity and favours a blueward track. Indeed, in Fig. 7, one sees that when the mean velocity on the MS becomes larger than about 130 km s<sup>-1</sup>, the $`T_{\mathrm{eff}}`$ at the end of the MS increases when rotation increases. For sufficiently low velocities, rotation acts as a small overshoot, extending the MS tracks towards lower $`T_{\mathrm{eff}}`$. This results from the fact that, at sufficiently low rotation, the effect of rotation on the convective core mass overcomes the effect of helium diffusion in the envelope. Indeed, for small rotation, the time required for helium mixing in the whole radiative envelope is very long, while hydrogen just needs to diffuse through a small amount of mass to reach the convective core. Since not all the stars have the same initial rotational velocity, one expects a dispersion of the luminosities at the end of the MS. For the 20 M models shown on Fig. 7 one sees a difference of $`0.3`$ in $`M_{\mathrm{bol}}`$ between the luminosities of the low and fast rotators at the end of the MS. Rotation induces also a scatter of the effective temperatures at the end of the MS. In reality, the dispersion results from both different initial velocities and also, for a given initial velocity, from different angles between the axis of rotation and the line of sight. Indeed due to the von Zeipel theorem (vZ24 (1924)) the star appears bluer seen pole–on than equator–on. When integrated over the visible part of the star, the effects due to orientation can reach a few tenths of a magnitude in luminosity and a few hundredths in $`\mathrm{log}T_{\mathrm{eff}}`$ (cf. Maeder and Peytreman 1970). Figure 8 shows the evolutionary tracks of non–rotating and rotating stellar models for initial masses between 9 and 120 M. For the rotating stellar models, the initial velocity is $`v_{\mathrm{ini}}`$ = 300 km s<sup>-1</sup>. There is little difference between tracks with $`v_{\mathrm{ini}}`$ = 200 or 300 km s<sup>-1</sup> (see also Talon et al. Ta97 (1997)). If the effects behaved like $`v_{\mathrm{ini}}^2`$, there would be larger differences. The present saturation effect occurs because outward transport of angular momentum by shears are larger when rotation is larger, also a larger rotation produces more mass loss, which makes a larger reduction of rotation during the evolution. Let us note however that for some surface abundance ratios as N/C or N/O (see Table 1), the increase from $`v_{\mathrm{ini}}`$ = 200 to 300 km s<sup>-1</sup> produces significant changes. Thus, the similarity of the evolutionary tracks does not necessarily imply the similarity of the surface abundances for these elements. Rotation reduces the MS width in the high mass range (M $`<40`$ M). Let us recall that when the mass increases, the ratio of the diffusion timescale for the chemical elements to the MS lifetime decreases (Maeder ma98 (1998)). As a consequence, starting with the same $`v_{\mathrm{ini}}`$ on the ZAMS, massive stars will be more mixed than low mass stars at an identical stage of their evolution. This reduces the MS width since greater chemical homogeneity makes the star bluer. Moreover, due to both rotational mixing and mass loss, their surface will be rapidly enriched in H–burning products. These stars will therefore enter the Wolf–Rayet phase while they are still burning their hydrogen in their core. This again reduces the MS width. For initial masses between 9 and 25 M, the MS shape is not much changed by rotation at least for $`v_{\mathrm{ini}}300`$ km s<sup>-1</sup>. ### 5.2 The post–Main-Sequence evolution The post–MS evolution of the most massive stars (M $`40`$ M) which become W–R stars will be discussed in a forthcoming paper. We shall just mention one point of general interest here: for low or moderate rotation, the convective core shrinks as usual during MS evolution, while for high masses ($`M>40`$ M) and large initial rotations ($`\frac{\mathrm{\Omega }}{\mathrm{\Omega }_{\mathrm{crit}}}0.5`$), the convective core grows in mass during evolution. This latter situation occurs in the fast rotating 60 M model shown on Fig. 8. These behaviours, i.e. reduction or growth of the core, determine whether the star will follow respectively the usual redwards MS tracks in the HR diagram, or whether it will bifurcate to the blue (cf. Maeder ma87 (1987); Langer la92 (1992)) towards the classical tracks of homogeneous evolution (Schwarzschild sc58 (1958)) and likely produce W–R stars. The stars with initial masses between 15 and 25 M become red supergiants (RSG). Rotation does not change qualitatively this behaviour but accelerates the redwards evolution, especially for the 15 and 20 M models. As a numerical example, for an initial $`v_{\mathrm{ini}}`$ = 300 km s<sup>-1</sup>, the model stars burn all their helium as red supergiants at $`T_{\mathrm{eff}}`$ below 4000 K, while the non–rotating models spend a significant part of the He–burning phase in the blue part of the HR diagram: for the non–rotating 15 and 20 M models, respectively 25 and 20% of the total He–burning lifetime is spent at $`\mathrm{log}T_{\mathrm{eff}}4.0`$. The behaviour of the rotating models results mainly from the enhancement of the mass loss rates. This effect prevents the formation of a big intermediate convective zone and therefore favours a rapid evolution toward the RSG phase (Stothers and Chin st79 (1979); Maeder ma81 (1981)). Let us note that the dispersion of the initial rotational velocities produces a mixing of the above behaviours. Very interestingly, for the 12 M model a blue loop appears when rotation is included. This results from the higher luminosity of the rotating model. The higher luminosity implies that the outer envelope is more extended, and is thus characterized by lower temperatures and higher opacities at a given mass coordinate. As a consequence, in the rotating model during the first dredge–up, the outer convective zone proceeds much more deeply in mass than in the non–rotating star. Typically in the non–rotating model the minimum mass coordinate reached by the outer convective zone is 6.6 M while in the rotating model it is 2.6 M. This prevents temporarily the extension in mass of the He–core and enables the apparition of a blue loop. Indeed the lower the mass of the He–core is, the lower its gravitational potential. According to Lauterborn et al. (la71 (1971), see also the discussion in Maeder and Meynet ma89 (1989)), a blue loop appears when the gravitational potential of the core $`\mathrm{\Phi }_c`$ is inferior to a critical potential $`\mathrm{\Phi }_{\mathrm{crit}}`$ depending only on the actual mass of the star which is about the same for the rotating and non–rotating model. This explains the appearance of a blue loop in the 12 M rotating model. For the 9 M model, the minimum mass coordinate reached by the outer convective zone is not much affected by rotation and the models with and without rotation present very similar blue loops. ### 5.3 Masses and mass–luminosity relations When rotation increases, the actual masses at the end of both the MS and the He–burning phases become smaller (cf. Tables 1 and 2). Typically the quantity of mass lost by stellar winds during the MS is enhanced by 60–100% in rotating models with $`v_{\mathrm{ini}}`$ = 200 and 300 km s<sup>-1</sup> respectively. For stars which do not go through a Wolf–Rayet phase, the increase is due mainly to the direct effect of rotation on the mass loss rates (in the present models through the formula proposed by Friend and Abbott Fr86 (1986)) and to the higher luminosities reached by the tracks computed with rotation. The fact that rotation increases the lifetimes also contributes to produce smaller final masses. For the most massive stars (M $`60`$ M), the present rotating models enter the Wolf–Rayet phase already during the H–burning phase (see also Maeder 1987; Fliegner & Langer 1995; Meynet 1999, 2000b). This reduces significantly the mass at the end of the H–burning phase. As indicated in Sect. 5.1, the initial distribution of the rotational velocities implies a dispersion of the luminosities at the end of the MS. This effect introduces a significant scatter in the mass–luminosity relation (Langer 1992; Meynet me98 (1998)), in the sense that fast rotators are overluminous with respect to their actual masses. This is especially true in the high mass star range in which the luminosity versus mass relation flattens. This may explain some of the discrepancies between the evolutionary masses and the direct mass estimates in some binaries (Penny et al. pe99 (1999)). Let us end this section by saying a few words about the mass discrepancy problem (see e.g. Herrero et al. He00 (2000)). For some stars, the evolutionary masses (i.e. determined from the theoretical evolutionary tracks) are greater that the spectroscopically determined masses. Interestingly, according to Herrero et al. (He20 (2000)), only the low gravity objects present (if any) a mass discrepancy. Even if most of the problem has collapsed and was shown to be a result of the proximity of O–stars to the Eddington limit (Lamers and Leitherer la93 (1993); Herrero et al. He99 (1999)) and of the large effect of metal line blanketing not usually accounted for in the atmosphere models of massive stars (Lanz et al. la96 (1996)), some discrepancy seems to be still present. The remaining mass discrepancy may arise, in part, from the use of non–rotating models for determining the evolutionary masses. High rotation produces larger He–cores and He–rich envelopes, all that implies overluminous stars with lower gravity as evolution proceeds. This can occur even for the slow rotating objects because they could have had a sufficiently high initial rotational velocity. The observation that the mass discrepancies are found only for low gravity objects may reflect the fact that rotation implies more and more important changes in the $`\mathrm{log}g_{\mathrm{eff}}`$ versus $`\mathrm{log}T_{\mathrm{eff}}`$ plane when evolution proceeds. Indeed, looking at Fig. 16 below one sees that at high gravity (i.e. at an early evolutionary stage), one does not expect any mass difference when using rotating or non–rotating tracks. In contrast, the mass difference between the rotating and the non–rotating tracks become more and more important in the high mass star range and for the low values of $`\mathrm{log}g_{\mathrm{eff}}`$. The rotating 40 M track crosses the non–rotating 60 M model on Fig. 16 indicating that differences of $``$30% in mass are quite possible. For the fast rotating objects, a part of the mass discrepancy may also be due to a possible underestimate of the gravity (see Sect. 7.1). Indeed, an underestimate of the gravity would also imply an underestimate of the spectroscopically determined masses (Herrero et al. 2000). ### 5.4 Lifetimes and isochrones Table 1 presents some properties of the models. Column 1 and 2 give the initial mass and the initial velocity $`v_{\mathrm{ini}}`$ respectively. The mean equatorial rotational velocity $`\overline{v}`$ during the MS phase is indicated in column 3. This quantity is defined by $$\overline{v}=1/t_H_0^{t_H}v(t)𝑑t,$$ where $`t_H`$ is the duration of the H–burning phase given in column 4, except for those stars which enter the Wolf–Rayet phase while still burning their hydrogen in their core, i.e. for the rotating 60 and 120 M models, for which $`t_H`$ has been replaced by the duration of the O–type star phase. The H–burning lifetimes $`t_H`$, the masses M, the equatorial velocities $`v`$, the helium surface abundance $`Y_s`$ and the surface ratios (in mass) N/C and N/O at the end of the H–burning phase are given in columns 4 to 9. The columns 10 to 13 present some properties of the models when the star is a blue supergiant (BSG) or an LBV star. For stellar models with M$`40`$ M, the “BSG stage” in Table 1 corresponds either to the stage when $`\mathrm{log}T_{\mathrm{eff}}=4.0`$ during the first crossing of the Hertzsprung–Russel diagram or to the bluest point on the blue loop if any for M$`12`$ M. For non–rotating stellar models with M $`60`$ M, the “LBV stage” corresponds to the point when the star has lost half of the matter ejected by the stellar winds between the end of the MS and the entrance into the W–R phase. In the case of rotating models, the “LBV stage” corresponds to the period during which the surface velocity becomes critical and huge mass loss rates ensue. Again, here we choose a model in the middle of this phase. The columns 14 to 19 present some characteristics of the stellar models at the end of the He–burning phase and $`t_{He}`$ is the He–burning lifetime. In Table 2 some properties of 20 M models at the end of the MS are indicated, $`v_{\mathrm{ini}}`$ and $`\overline{v}`$ have the same meaning as above. From Table 1 one sees that for $`Z=0.020`$ the lifetimes are increased by about 20–30% when the mean rotational velocity on the MS increases from 0 to $``$200 km s<sup>-1</sup>. This modest increase is explained by the fact that even if there is more fuel available in the core, the luminosity is also increased. From the data presented in Table 2, one can deduce a nearly linear relation between the relative enhancement of the MS lifetime, $`\mathrm{\Delta }t_\mathrm{H}`$ and $`\overline{v}`$, where $$\mathrm{\Delta }t_\mathrm{H}(\overline{v})=[t_\mathrm{H}(\overline{v})t_\mathrm{H}(0)]/t_\mathrm{H}(0).$$ One obtains, $$\frac{\mathrm{\Delta }t_\mathrm{H}(\overline{v})}{t_\mathrm{H}(0)}=0.0013\overline{v},$$ where $`\overline{v}`$ is in km s<sup>-1</sup>. This relation reproduces the values of $`\mathrm{\Delta }t_H(\overline{v})/t_H(0)`$ from table 2 with an accuracy better than 5%. It also applies with the same accuracy to the values listed in Table 1 for the masses between 15 and 40 M. The He–burning lifetimes are less affected by rotation than the MS lifetimes. The changes are less than 10%. The ratios $`t_{He}/t_H`$ of the He to H–burning lifetimes are only slightly decreased by rotation and remain around 10–15%. On Fig. 9 isochrones for ages between about 8 and 20 10<sup>6</sup> yr, computed from non–rotating and rotating stellar models are presented. The “rotating” isochrones are computed from the models with an initial rotational velocity $`v_{\mathrm{ini}}`$ of 200 km s<sup>-1</sup> on the ZAMS. At a given age, the upper part of the “rotating” isochrones are bluer and more luminous. Typically, for an age equal to about $`2010^6`$ years ($`\mathrm{log}`$ age =7.3) the reddest point on the MS is shifted by 0.015 dex in $`\mathrm{log}T_{\mathrm{eff}}`$ and by - 0.6 in $`M_{\mathrm{bol}}`$ when rotation is taken into account. An isochrone with rotation is almost identical to an isochrone without rotation with log age smaller by 0.1 dex. This has for consequence that rotation slightly increases the age associated to a given cluster. From Fig. 9, one sees that the “rotating” isochrone for an age equal to $`2010^6`$ years has the same luminosity at the turn–off than the “non-rotating” isochrone for an age equal to $`1610^6`$ years. Thus rotation increases the age estimate by about 25%. We may wonder whether this effect explains the age difference between the estimates based on the upper MS and the estimates based on the lithium content of the very low mass stars (see e.g. Martin et al. 1998; Barrado y Navascués et al. 1999). One must also account for the dispersion of rotational velocities and possibly of the orientation angles. These two effects introduce some dispersion in the way stars are distributed in the HR diagram and thus affect the interpretation of the clusters’ observed sequences (cf. Maeder ma71 (1971)). If a bluewards track occurs, as for very massive stars with fast rotation, the larger core and mixing lead to much longer lifetimes in the H–burning phase. In this case, the fitting of time–lines becomes hazardous. ## 6 Evolution of the rotational velocities ### 6.1 Model results for stars with a large mass loss Figs. 10 and 11 show the evolution of $`v`$ and of $`\frac{\mathrm{\Omega }}{\mathrm{\Omega }_{\mathrm{crit}}}`$ as a function of the age for the present models. Firstly, we notice that for models without mass loss, as shown for the 20 M with $`\dot{M}`$ = 0, $`v`$ and $`\frac{\mathrm{\Omega }}{\mathrm{\Omega }_{\mathrm{crit}}}`$ go up fastly so that the critical velocity would be reached near the end of the MS phase. The current model of 20 M with mass loss show a significant decrease of $`v`$, while the critical ratio remains almost constant during most of the MS phase. Figs. 10 and 11 show how fastly rotation decreases at the surface of the most massive stars, which lose a lot of mass. Consistently we see that the reduction of the surface rotation is much larger for the more massive stars. This is of course a consequence of the removal of large amounts of angular momentum by the stellar winds. The effect is amplified by the increase of the mass loss in fast rotators (Eq. 2). We see that the decrease of $`\frac{\mathrm{\Omega }}{\mathrm{\Omega }_{\mathrm{crit}}}`$ is so strong that it will prevent a massive star to reach the critical velocity near the end of the MS phase. If the star makes extended excursions in the HR diagram at the end of the MS like is the case for the 60 M model, then it may reach the critical velocity. The specific case of stars close to the $`\mathrm{\Omega }`$–limit will be examined in a future study, since this requires some further theoretical developments. Let us mention here that the present results differ from those obtained by Sackmann and Anand (Sa70 (1970)) and Langer (La97 (1997), La98 (1998)). Indeed, these authors find that the star reaches the break–up limit during the MS phase. As an example, in a 60 M star, even a model with an initial $`v_{\mathrm{ini}}`$ of 100 km s<sup>-1</sup> reaches the break–up limit near the end of the MS–phase (Langer 1998). This result leads Langer to conclude that most massive stars may reach the break–up limit or the so–called $`\mathrm{\Omega }`$–limit during their MS evolution. Let us however emphasize here that such a conclusion is based on a particular definition of $`v_{\mathrm{crit}}`$ still subject to discussion (see Sect. 2.4), on the assumption of solid body rotation, and on models not accounting for the effects of rotationnally induced mixing. We see here that modifying these hypothesis (and also using other prescriptions for the mass loss rates) lead to very different results. One of the first step to clarify the situation is to determine which expression for $`v_{\mathrm{crit}}`$ (cf. Glatzel 1998; Langer 1998) is the correct one and how rotation affects the mass loss rates. These developments, now in progress, will be particularly needed for the study of the evolution of the most massive stars, like a 120 M model, which have a high value of the Eddington factor. Also this is important for the formation and evolution of W–R stars, which will be studied in a further work. Fig. 12 shows the evolution of $`v`$ with age for models of a 20 M star with different initial velocities $`v_{\mathrm{ini}}`$ from 50 to 580 km s<sup>-1</sup>. We see that the decrease in the surface $`v`$ is larger for larger initial rotation. This is a consequence of the larger mass loss rates in fast rotators. We notice some convergence (cf. Langer 1998) of the curves for the large initial $`v_{\mathrm{ini}}`$. This convergence would be more pronounced and affect the models with lower $`v_{\mathrm{ini}}`$ if the dependence of $`\dot{M}`$ vs $`v`$ would be stronger than given by Eq. (2). This certainly reduces the scatter of $`v`$ near the end of the MS phase, but does not produce a full convergence. ### 6.2 Results for stars with lower mass loss The models of 12 and 15 M show curves in Figs. 10 and 11 with little reduction of $`v`$, while there are slight increases of the critical ratio $`\frac{\mathrm{\Omega }}{\mathrm{\Omega }_{\mathrm{crit}}}`$ during MS evolution. A peak is reached during the overall contraction phase at the end of the MS phase. Then, $`v`$ and the critical ratio go down as the star moves to the red supergiant phase. During such a fast evolution, we may say that the rotation evolves almost like the case of simplified models, where the angular momentum is conserved locally. The large growth of the radius just implies a decrease of $`v`$ and of $`\frac{\mathrm{\Omega }}{\mathrm{\Omega }_{\mathrm{crit}}}`$. Later during the blue loops, where the Cepheid instability strip is crossed, the rotation velocity becomes very large again and could easily become close to critical. This behaviour, also found by Heger & Langer (1998), results from the stellar contraction which concentrates a large fraction of the angular momentum of the star (previously contained in the extended convective envelope of the RSG) in the outer few hundredths of a solar mass. This result suggests that rotation may also somehow influence the Cepheid properties, in addition to the increase of stellar luminosity discussed above. Amazingly, we notice that the stars with initially low mass loss during the MS phase, like for stars with M $`12`$ M, have more chance to reach the break–up velocities and thus huge mass loss than the more massive stars which lose a lot of mass on the MS. It is somehow surprising that little mass loss during the MS may favour large mass loss rates at the end of the MS phase. This may explain why stars close to break–up, like the Be stars, do not form among the O–type stars, but mainly among the B–type stars, where we see that the ratio $`\frac{\mathrm{\Omega }}{\mathrm{\Omega }_{\mathrm{crit}}}`$ may increase during the MS phase. Another related observation is the fact that the relative number of Be–stars with respect to B–type stars is much higher in the LMC and SMC than in the Galaxy (cf. Maeder et al. MGM99 (1999)). In the LMC and SMC, due to the lower metallicity, the average mass loss rates are lower and thus these stars may keep higher rotation in general and thus form more Be stars. This explanation does not exclude differences in the distribution of $`v_{\mathrm{ini}}`$ as well. ### 6.3 The rotational velocities in the HR diagram The evolution of $`v`$ in the HR diagram is shown in Fig. 13. Starting with models having a velocity of 300 km s<sup>-1</sup> on the zero–age main sequence, we give some lines of constant $`v`$ over the HR diagram. On the MS, we notice in particular that the decrease is much faster for the most massive stars than for stars with M $``$ 15 M. This difference remains also present in the domain of B–supergiants. During the crossing of the HR diagram, the rotational velocities decrease fastly, to become very small, i.e. of the order of a few km s<sup>-1</sup>, in the red supergiant phase. It is beyond the scope of this paper to make detailed comparisons of the evolution of the distribution of the velocities over the HR diagram, however, we may notice a few points. The fact that the average $`\overline{v}`$ is lower for O–type stars than for the early B–type stars (Slettebak, Sl70 (1970)) may be the consequence of the higher losses of mass and angular momentum in the most massive stars. Also, we remark that the increase of $`\overline{v}`$ from O–stars to B–stars is larger for the stars of luminosity class IV than for class V (Fukuda, Fu82 (1982)). This is consistent with our models, which show (cf. Fig. 13) that the differences of $`\overline{v}`$ beween O– and B–type stars are much larger at the end of the MS phase. Another fact in the observed data is the strong decrease of $`\overline{v}`$ for the massive supergiants of OB–types. This is predicted by all stellar models (cf. also Langer La98 (1998)) due to the growth of the stellar radii. Further detailed comparisons may perhaps provide some new tests and constraints. ## 7 Evolution of the surface abundances The chemical abundances offer a very powerful test of internal evolution and they give strong evidences in favour of some additional mixing processes in massive stars. A review of the observations may be found in Maeder and Meynet (2000). Here we shall concentrate on the discussion of the theoretical results and we shall compare them with some recent observations. The most striking feature appearing in Tables 1 and 2 as well as on Fig. 14 is the change of surface abundances in rotating stellar models (cf. Langer 1992). The He–, N–enrichments and the related C– and O–depletions at the surface already occur on the MS. The more massive the star is, the more pronounced are the enrichments, supporting the expectation that mixing becomes more and more efficient when the mass increases (Maeder ma98 (1998)). The same is true when the initial rotational velocity increases. During the crossing of the Hertzsprung–Russel diagram, the evolution is sufficiently short for not allowing any important change of the surface abundances (see Fig. 14). Further changes occur when the star becomes a red supergiant and undergoes the first dredge–up. In contrast, non–rotating stellar models show no change of the surface abundances until the first dredge–up in the red supergiant stage (see Table 1). This implies that these models predict some enrichment neither during the MS nor in the blue supergiant phase unless a blue loop is formed. Moreover, as can be seen in Fig. 14, the ratios obtained at the end of the He–burning phase are significantly lower than the ones obtained in rotating models. ### 7.1 He–enrichments and $`v\mathrm{sin}i`$ Fig. 15 shows the evolution in the $`ϵ`$ versus surface velocity plane where $`ϵ=\frac{n(\mathrm{H}_\mathrm{e})}{\mathrm{n}(\mathrm{H})+\mathrm{n}(\mathrm{He})}`$, $`n(\mathrm{H}_\mathrm{e})`$ and $`n(\mathrm{H})`$ being the abundances in number of helium and hydrogen at the surface. The theoretical tracks correspond to the equatorial velocities $`v`$, while the observed points from Herrero et al. (He92 (1992), He99 (1999), He20 (2000)) are $`v\mathrm{sin}i`$. This implies that the observed values are smaller on the average by a factor $`\pi /4`$ with respect to the theoretical values. In Fig. 15, the tracks go generally from the bottom on the right to the top on the left. Along a given track, the He–enrichment increases when the velocity decreases. The shaded zone in Fig. 15 corresponds to the MS for the models with $`v_{\mathrm{ini}}`$ = 300 km s<sup>-1</sup>. The higher initial mass, the greater the He–enrichments which can be reached during the MS. We see also that during the MS phase, for initial masses inferior to about 60 M, all the tracks with $`v_{\mathrm{ini}}`$ = 300 km s<sup>-1</sup> follow more or less the same $`ϵ`$ versus $`v`$ relation. However, the relation is changed when the initial velocity is different (see the fast rotating 20 and 60 M models). The surface He–enrichments on the MS generally depend on the following factors: the initial mass, the initial metallicity (Maeder & Meynet 2000b; Meynet 2000a), the initial velocity and the age of the star. A low surface velocity at a given evolutionary stage does not exclude that in the past the star was a fast rotator. The slow rotation may result from the loss of angular momentum by stellar winds and/or from the increase of radius of the star. On the other hand, a high velocity in the past implies He– and N–enrichments of the surface. Let us compare in Fig. 15 the theoretical tracks in the $`ϵ`$ versus $`v`$ plane with the observations of OB stars performed by Herrero et al. (He92 (1992), He99 (1999), He20 (2000)). Recent works (McErlean et al. mc98 (1998); Smith and Howarth sm98 (1998)) indicate that accounting for a microturbulent velocity line broadening in the model atmosphere reduces the derived He–abundances for supergiants later than O9. However, according to Villamariz and Herrero (vi99 (1999)) this effect cannot explain all the observed overabundances, especially for the earlier types. On Fig. 16 the observed points are plotted in the $`\mathrm{log}g_{\mathrm{eff}}`$ versus $`\mathrm{log}T_{\mathrm{eff}}`$ plane where $`g_{\mathrm{eff}}`$ is the effective surface gravity. We estimate $`g_{\mathrm{eff}}=\frac{GM}{R^2}\mathrm{\Omega }^2R\mathrm{sin}\theta `$ for the average orientation angle. Let us note that for the models plotted in Fig. 16, there is little difference between the effective gravities at the pole and at the equator. Indeed the ratio between these two gravities never exceeds 1.3 which means a vertical dispersion of about 0.1 dex in Fig. 16. Non–rotating and rotating evolutionary tracks are superposed to the observed points in Fig. 16. Since most of the enriched stars are in the vicinity of the 120 and 60 M tracks (see Fig. 16), we can wonder wether the changes of the surface abundances can be explained as an effect of mass loss only. It does not seem to be the case, because the part of the track shown on Fig. 16 for the non–rotating 60 M model presents no surface He–enrichment. In the case of the non–rotating 120 M model only the part of the track with $`\mathrm{log}g_{\mathrm{eff}}`$ inferior to 3.3 has $`ϵ>0.14`$. Thus the He–enrichments cannot be accounted for by current evolutionary models as was already pointed out by Herrero et al. (He92 (1992), cf. also Maeder 1987). In the following we shall suppose that these enhancements are due to rotation. Let us consider four groups of stars. In the first group we place all the stars presenting no He–enrichment at their surface ($`ϵ0.12`$) and having $`v\mathrm{sin}i<200`$ km s<sup>-1</sup> (empty squares on Figs. 15 and 16), in the second one are the enriched stars ($`ϵ>`$ 0.12, solid triangles) having $`v\mathrm{sin}i<200`$ km s<sup>-1</sup>. The third group (which is empty at present) consists of the fast rotators with no He–enrichment. These stars would occupy the bottom right corner in Fig. 15. Finally, the fourth group contains the He–enriched stars with $`v\mathrm{sin}i200`$ km s<sup>-1</sup> (solid circles). Group 1 The non–enriched stars with low rotation can be interpreted either as stars with small initial velocities or as young fast rotators whose surface has not yet been enriched in helium by rotational mixing (see also Herrero et al. He99 (1999)). In this respect let us mention that all the stars having $`\mathrm{log}g_{\mathrm{eff}}>3.7`$, and which therefore are probably not too evolved, present no He–enrichment. Group 2 The most striking feature of the second group of stars ($`ϵ>0.12`$ and $`v\mathrm{sin}i<200`$ km s<sup>-1</sup>) is the fact that they are distributed in a relatively narrow range of $`v\mathrm{sin}i`$ between 80 and 160 km s<sup>-1</sup>. This may result from the following facts: firstly there exists a minimum value of the initial velocity for rotational mixing to be able to drive changes of the helium surface abundances during the H–burning phase. Secondly the observed distribution also reflects the way the surface velocity declines during the H–burning phase and the narrow range of observed velocities may result from some convergence effect as the one mentionned in Sect. 6.1 (see also Fig. 12). We see that the high $`ϵ`$ values reached by some of these stars (superior to $``$0.16) would be compatible with their high initial mass implied from their position in Fig. 16. Group 3 Very interestingly no stars are observed with a $`v\mathrm{sin}i200`$ km s<sup>-1</sup> and no He–enrichment. Keeping in mind that the observed stars do not represent a statistical complete sample, one can nevertheless wonder why no stars are observed in this zone. A possibility would be that these very fast rotators do not exist, but this is not realistic. Indeed stars with velocities as high as 300 km s<sup>-1</sup> and more have been observed on the MS (e.g. Penny pe96 (1996); Howarth et al. ho97 (1997); see also Fig. 15). Moreover if such stars were not formed how to explain the stars in Group 4, having a very high $`v\mathrm{sin}i`$ and an important He–enrichment (the solid circles in Fig. 15) ? These stars are likely the descendants of very fast MS rotating stars (see the discussion in the subsection Group 4 below). The fact that no stars are observed in the Group 3 may indicate that the change of the surface abundances occur within a small fraction of the visible MS life. For the fast rotating 60 M model plotted on Fig. 15, the surface retains its initial composition ($`ϵ<0.11`$) only during a third of its H–burning lifetime. More rapidly rotating models would still reduce the fraction of the MS time spent with no change of the surface abundances. The lack of fast rotators with normal surface composition may be due to the fact that young massive fast rotators are still embedded in the cloud from which they formed. The models of massive star formation with accretion (Bernasconi and Maeder be96 (1996)) indicate that when the stars become visible, they have already burnt some fraction of their central hydrogen, thus some transport of He and N to the stellar surface may have occured. Group 4 Could the stars belonging to the fourth group ($`ϵ>0.12`$ and $`v\mathrm{sin}i200`$ km/ s) be formed by stars previously in the low velocity range and which have been accelerated for instance by contraction on a blue loop ? The answer is likely no. Indeed blue loops cannot accelerate the surface beyond the critical velocity which, for a 12 M blue supergiant at the tip of a blue loop, is of the order of 250 km s<sup>-1</sup>. In addition some stars in the group 4 are classified as MS stars (see the stars labeled with a V on Figs. 15 and 16). Another possibility would be to consider the stars in Group 4 as secondary stars in close binary systems which would have been accelerated through the process of mass accretion ? But these stars are not observed to belong to binary systems. Therefore the most reasonable hypothesis is to consider the stars in Group 4 as the natural descendants of very fast MS rotating stars. Their chemical enrichment at the surface is very fast. Their chemical structure as a result of the strong rotational mixing is probably near homogeneity. This view is supported by the fact that high values of $`v\mathrm{sin}i`$ are observed for very high values of $`ϵ`$ implying that the surface velocity does not decrease too much in the course of the evolution. This can be accounted for if the star remains compact, i.e. in the blue part of the HR diagram as is the case for a strongly mixed star (Maeder ma87 (1987); see also the fast rotating 60 M track in Fig. 15). Another effect could also be important in that respect, i.e. the anisotropy of the stellar winds when stars are rotating near break–up. For the hot stars, the von Zeipel (vZ24 (1924)) theorem implies that most of the mass is ejected from the pole (Maeder Mae99 (1999)). This prevents the loss of important angular momentum and maintains a high surface velocity. If the stars in the fourth group are nearly homogeneous objects, one would expect higher effective gravities than observed. Typically, the fast rotating well mixed 60 M model remains at a high value of $`\mathrm{log}g_{\mathrm{eff}}`$, at least for the portion of the evolution computed here (see Fig. 16). Moreover, as noted above, the maximum dispersion in $`\mathrm{log}g_{\mathrm{eff}}`$ due to orientation effects is around 0.1 dex. Could the $`\mathrm{log}g_{\mathrm{eff}}`$ be underestimated ? Even if it is difficult on the base of the present data to ascertain such a point of view, one can note that the $`\mathrm{sin}i`$ for these fast rotating stars cannot be too far from 1, otherwise one would obtain surface velocities above the break–up value. This means that these stars are seen essentially equator–on. The $`\mathrm{log}g`$ in the equator band is inferior to the surface averaged $`\mathrm{log}g`$ and thus the observed $`\mathrm{log}g`$ might be underestimated. Of course a quantitative analysis requires a detailed study of how the variation of $`\mathrm{log}g`$ with the latitude affects the spectroscopically determined gravities. ### 7.2 Comparisons with the surface abundances of supergiants Most blue supergiants present surface enrichments. For instance, Walborn (wa76 (1976), wa88 (1988)) showed that ordinary OB supergiants have He– and N– enrichments as a result of CNO processing. Only the small group of peculiar OBC–stars has the normal cosmic abundance ratios (cf. also Howarth and Prinja ho89 (1989); Herrero et al. He92 (1992); Gies and Lambert gi92 (1992); Lennon le94 (1994)). A possibility to explain the He– and N–enrichments in supergiants would be that blue supergiants are on a blue loop after a first red supergiant stage where dredge–up has occurred producing the observed surface enrichments. However as we shall see below, this does not appear as the good explanantion at least for some of the observed enrichments. Fig. 17 illustrates the changes of the nitrogen to carbon ratios N/C from the ZAMS to the red supergiant stage for stars in the mass range from 9 to 20 M. The N/C ratio appears as the most sensitive observable parameter. For non–rotating stars, the surface enrichment in nitrogen only occurs when the star reaches the red supergiant phase; there, CNO elements are dredged–up by deep convection. The behaviour is the same as for the He–enrichments discussed above. For rotating stars, N–excesses already occur during the MS phase and they are larger for higher rotation and initial stellar masses. At the end of the MS phase of the 12 M model, the N/C ratio is enhanced by factors 2.4 and 4.3 for $`v_{\mathrm{ini}}`$ = 200 and 300 km s<sup>-1</sup> respectively. These factors are increased up to 4 and 7.1 for the 20 M models with the same $`v_{\mathrm{ini}}`$. On Fig. 17 we plot also some observations of supergiants performed by Gies and Lambert (gi92 (1992)), Lennon (le94 (1994)) and Venn (1995a b). The plotted values are $`\mathrm{\Delta }\mathrm{log}\mathrm{N}/\mathrm{C}=[\mathrm{N}/\mathrm{C}]_{}[\mathrm{N}/\mathrm{C}]_B`$, where the abundance ratios in number are measured at the surface of the star ($``$), and at the surface of main sequence B–stars ($`B`$) supposed to have retained their pristine cosmic abundances. Gies and Lambert (gi92 (1992)) provide “LTE” and “NLTE” N/C ratios. We plot here the smaller “LTE” ratios since, for the stars in common with the Lennon sample, they are similar to Lennon’s results (see the discussion in Venn 1995b ). From the positions of the observed stars in the $`\mathrm{log}g_{\mathrm{eff}}`$ versus $`\mathrm{log}T_{\mathrm{eff}}`$ diagram, one obtains that the range of initial masses for the A–type supergiants shown on Fig. 17 are between 5 and 20 M. From Fig. 17, one notes first that at the $`T_{\mathrm{eff}}`$ of the observed supergiants, non–rotating stellar models predict no enrichment at all unless there is a blue loop. However it is not likely that all the observed points are at the tip of a blue loop. Indeed some of the stars have initial masses above 15 M and, at solar metallicity, current grids of models (Arnett ar91 (1991); Schaller et al. sch92 (1992); Alongi et al. al93 (1993); Brocato and Castellani br93 (1993)) only predict blue loops for masses equal or lower than 15 M. Moreover many of the observed stars have N/C ratios too low to result from a first dredge–up episode (see for instance the position of the blue loop of the non–rotating 9 M model in Fig. 17). Therefore Venn (1995b ) suggested that at least those stars presenting the lowest N/C ratios are on their way from the MS to the red giant branch and have undergone some mixing in the early stage of their evolution. If such stars are not at all accounted for by standard evolutionary tracks, rotating models can naturally reproduce their observed surface abundances as can be seen on Fig.17. Moreover as already noted above, theory predicts larger excesses for higher masses, a result in agreement with the suggestion of Takeda and Takeda–Hidai (ta95 (1995)) recently confirmed by McErlean et al. (mc99 (1999)). ## 8 Conclusion Mass loss by stellar winds and rotational mixing in the stellar interior are certainly the two hydrodynamical phenomena which most deeply affect the evolution of massive stars. Far from being a small refinement in the physics of stellar interior, rotation appears as an essential ingredient of future grids of stellar models. In particular among the important points wich are not discussed here but which will be studied in more details in forthcoming papers are the effects of rotation on the population of red and blue supergiants at various metallicities, on the evolutionary scenarios leading to the formation of Wolf–Rayet stars (Maeder 1987; Fliegner and Langer 1995; Meynet 2000b) and on the stellar yields (Heger et al. 2000). These questions are important for a better understanding of starbursts regions and of the chemical evolution of galaxies.
warning/0006/cond-mat0006093.html	ar5iv	text	# Theory of Magnetic Anisotropy in III_{1-𝑥}⁢Mn_𝑥⁢V Ferromagnets ## I Introduction The discovery of carrier-mediated ferromagnetism in $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ and doped $`\mathrm{II}_{1x}\mathrm{Mn}_x\mathrm{VI}`$ diluted magnetic semiconductors (DMS’s) has opened up a broad and relatively unexplored frontier for both basic and applied research. Experiments in $`\mathrm{Ga}_{1x}\mathrm{Mn}_x\mathrm{As}`$ and $`\mathrm{In}_{1x}\mathrm{Mn}_x\mathrm{As}`$ have demonstrated that these ferromagnets have remarkably square hysterisis loops with coercivities typically $`40\mathrm{O}\mathrm{e}`$, and that the magnetic easy axis is dependent on epitaxial growth lattice-matching strains. In this paper we discuss the magnetic anisotropy properties of $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ DMS ferromagnets, predicted by a mean-field-theory of the exchange interaction coupling between localized magnetic ions and valence band free carriers. We use phenomenological four or six band envelope function models, depending on the carrier density $`p`$, in which the valence band holes are characterized by Luttinger , spin-orbit splitting, and strain-energy parameters. The physical origin of the anisotropy energy in our model is spin-orbit coupling in the valence band. Our work is based in part on theoretical descriptions developed by Gaj and Bastard to explain the optical properties of undoped, paramagnetic DMS’s. As the critical temperature is approached, the mean-field-theory we employ reduces to an earlier theory that invokes generalized RKKY carrier-mediated interactions between localized spins. The two approaches differ, however, in their description of the magnetically ordered state. As this work was nearing completion we learned of a closely related study which uses the same mean-field theory to address critical temperature trends in this material class and which also comments on the mean-field theory’s ability to address magnetic anisotropy physics. We are aware of three elements of the physics of these materials which make the predictions of our mean-field theory uncertain: i) we do not account for the substantial disorder which is usually present in these ferromagnets; ii) we do not include effects due to interactions among the itinerant holes and iii) we do not account for correlations between localized spin-configurations and itinerant hole states. The importance of each of these deficiencies is difficult to judge in general, and probably depends on adjustable material parameters. In our view, it is likely that there is a substantial range in the parameter space of interest where the predictions of the present theory are useful. We expect that important progress can be made by comparing this simplest possible theory of carrier-induced DMS ferromagnetism with experiment. This work has two objectives. Most importantly, we have attempted to shed light on how various adjustable material parameters can influence magnetic anisotropy. Secondarily, we have made an effort to estimate the magnetic anisotropy energy in those cases where experimental information is presently available. Our hope here is to initiate a process of careful and quantitative comparison between mean-field theory and experiment, partially to help judge the efficiency of this approximation in predicting other physical properties. Even in the mean-field theory, we find that the magnetic anisotropy physics of these materials is rich. We predict easy axis reorientations as a function of hole density, exchange interaction strength, temperature, and strain and identify situations under which $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ ferromagnets are remarkably hard. In Section II we detail our mean-field-theory of the ordered state. The theory simplifies in the limit of low-temperature and low hole densities. Our results for this limit, presented in Section III, predict a $`111`$ easy axis in the absence of strain, and a magnetic anisotropy energy which is approximately $`10\%`$ of the free carrier band-energy density. This value is extremely large for a cubic metallic ferromagnet; typical ratios in transition metal ferromagnets are smaller than $`10^6`$ for example. The anisotropy energy in this limit varies as the free-carrier density to the $`5/3`$ power and is independent of the exchange coupling strength. Explicit results for the strain-dependence of the magnetic anisotropy in the same limit are presented in Sec. IV. We find that unrelaxed lattice-matching strains due to epitaxial growth contribute a uniaxial anisotropy which favors magnetization orientation along the growth direction when the substrate lattice constant is larger than the ferromagnetic semiconductor lattice constant and an in-plane orientation in the opposite case. Unfortunately, perhaps, the simple low density limit does not normally apply in situations where high critical temperatures are expected. The more complicated, and more widely relevant, general case is discussed in Section V. We find that magnetic anisotropy has a non-trivial dependence on both temperature and exchange coupling strength and that easy axis reversals occur, in general, as a function of either parameter. According to our theory, anisotropy energy densities comparable to those in typical metallic ferromagnets are possible when the exchange coupling is strong enough to depopulate all but one of the spin-split valence bands, even with saturation magnetization values smaller by more than an order of magnitude. In the limit of large hole densities, we find that the anisotropy energy of strain-free samples is proportional to hole density $`p^1`$ and exchange coupling $`J_{pd}^4`$. We find that in typical situations a strain $`e_0`$ of only $`1\%`$ is sufficient to overwhelm the cubic anisotropy of strain-free samples. We conclude in Section VI with a discussion of the implications of these calculations for the interpretation of present experiments, and with some suggestions for future experiments which could further test the appropriateness of the model used here. ## II Formal Theory Our theory is based on an envelope function description of the valence band electrons, and a spin representation for their kinetic-exchange inteaction with localized $`d`$ electrons on the Mn<sup>++</sup> ions: $$=_m+_b+J_{pd}\underset{i,I}{}\stackrel{}{S}_I\stackrel{}{s}_i\delta (\stackrel{}{r}_i\stackrel{}{R}_I),$$ (1) where $`i`$ labels a valence band hole and $`I`$ labels a magnetic ion. In Eq. (1), $`_m`$ describes the coupling of magnetic ions with total spin quantum number $`J=5/2`$ to an external field (if one is present), $`\stackrel{}{S}_I`$ is a localized spin, $`\stackrel{}{s}_i`$ is a hole spin, and $`_b`$ is either a four or a six-band envelope-function Hamiltonian for the valence bands. In this paper we do not consider external magnetic fields so $`_m0`$. The four-band Kohn-Luttinger model describes only the total angular momentum $`j=3/2`$ bands, and is adequate when spin-orbit coupling is large and the hole density $`p`$ is not too large. As discussed later, in the case of GaAs, a four-band model suffices for $`p10^{18}\mathrm{cm}^3`$. In III<sub>1-x</sub>Mn<sub>x</sub>V semiconductors, the four $`j=3/2`$ bands are separated by a spin-orbit splitting $`\mathrm{\Delta }_{so}`$ from the two $`j=1/2`$ bands. In the relevant range of hole and $`\mathrm{Mn}^{++}`$ densities, no more than four bands are ever occupied. Nevertheless, mixing between $`j=3/2`$ and $`j=1/2`$ bands does occur, and it can alter the balance of delicate cancellations which often controls the net anisotropy energy. The exchange interaction between valence band holes and localized moments is believed to be antiferromagnetic, i.e. $`J_{pd}>0`$. For GaAs, experimental estimates of $`J_{pd}`$ fall between $`0.04\mathrm{eVnm}^3`$ and $`0.15\mathrm{eVnm}^3`$, with more recent work suggesting a value toward the lower end of this range. The form of the valence band for Bloch wavevectors near the zone center in a cubic semiconductor follows from $`𝐤𝐩`$ perturbation theory and symmetry considerations. The four band ($`j=3/2`$) and six band ($`j=3/2`$ and 1/2) models are known as Kohn-Luttinger Hamiltonians and their explicit form is given in the Appendix. The eigenenergies are measured down from the top of the valence band, i.e., they are hole energies. The Kohn-Luttinger Hamiltonian contains the spin-orbit splitting parameter $`\mathrm{\Delta }_{so}`$ and three other phenomenological parameters, $`\gamma _1`$, $`\gamma _2`$ and $`\gamma _3`$. These are accurately known for common semiconductors. For GaAs and InAs, the two materials in which $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ ferromagnetism has been observed, $`\mathrm{\Delta }_{so}=0.34`$ eV and 0.43 eV, and $`(\gamma _1,\gamma _2,\gamma _3)=(6.85,2.1,2.9)`$ and $`(19.67,8.37,9.29)`$ respectively. Most of the specific illustrative calculations discussed below are performed with GaAs parameters. Our calculations are based on the Kohn-Luttinger Hamiltonian and on a mean-field theory in which correlations between the local-moment configuration and the itinerant carrier system are neglected. We comment later on limits of validity of this approximation. There are a number of equivalent ways of developing this mean-field theory formally. In the following paragraphs we present a view which is convenient for discussing magnetic anisotropy. In the absence of an external magnetic field, the partition function of our model may be expressed exactly as a weighted sum over magnetic impurity configurations specified by a localized spin quantization axis, $`\widehat{M}`$, and azimuthal spin quantum numbers $`m_I`$: $$Z=\underset{m_I}{}\mathrm{exp}(F_b[m_I]/k_BT),$$ (2) where $`F_b[m_I]`$ is the valence band free energy for holes which experience an effective Zeeman magnetic field $$\stackrel{}{h}(\stackrel{}{r})[m_I]=J_{pd}\widehat{M}\underset{I}{}m_I\delta (\stackrel{}{r}\stackrel{}{R}_I).$$ (3) The mean-field theory consists of replacing $`\stackrel{}{h}(\stackrel{}{r})[m_I]`$ by its spatial average for each magnetic impurity configuration, thereby neglecting correlations between spin-distributions in local-moment and hole subsystems. The effective Zeeman magnetic field experienced by the holes then depends only on $`\widehat{M}`$, the direction of the local-moment orientation, and the mean azimuthal quantum number averaged over all local moments, $`M`$: $$\stackrel{}{h}_{MF}(M)=J_{pd}N_{Mn}M\widehat{M}h\widehat{M},$$ (4) where $`N_{Mn}=N_I/V`$ is the number of magnetic impurities per unit volume. The mean-field partition function is $$Z_{MF}(M)=\mathrm{exp}((N_ITs(M)F_b(\stackrel{}{h}))/k_BT),$$ (5) where the entropy per impurity is defined by $$s(M)=k_B\underset{N_I\mathrm{}}{lim}\frac{\mathrm{ln}[_{m_I}\delta (_Im_IN_IM)]}{N_I},$$ (6) and $`F_b(\stackrel{}{h})`$ is the free-energy of a system of non-interacting fermions with single-particle Hamiltonian $`H_bh\widehat{M}\stackrel{}{s}`$. Following standard ‘large number’ arguments $`s(M)`$ is readily evaluated by considering an auxiliary system consisting of magnetic impurities coupled only to an external magnetic field $`H`$. For this model problem, a familiar exercise gives the result $$M(H)=JB_J(x),$$ (7) where $`x=g_L\mu _BHJ/k_BT`$, $`g_L`$ is the Landé g-factor of the ion, $`\mu _B`$ is the electron Bohr magneton, and $$B_J(x)=\frac{2J+1}{2J}\mathrm{coth}\left[(2J+1)x/2J\right]\frac{1}{2J}\mathrm{coth}(x/2J),$$ (8) is the Brillouin function. The Brillouin function is a one-to-one mapping between reduced fields $`x`$ in the interval $`[0,\mathrm{}]`$ and reduced magnetizations $`M/J`$ in the interval $`[0,1]`$; the inverse function $`B_J^1`$ maps $`M/J`$ to $`x`$. Since the magnetization maximizes $`s(M)+g_L\mu _BHM/k_BT`$, $$\frac{ds(M)}{dM}=g_L\mu _BH/k_BT.$$ (9) Eq. (9) can be used to eliminate $`H`$ and arrive at the following explicit expression for $`s(M)`$: $$s(M)=k_B_{B_J^1(M/J)}^{\mathrm{}}𝑑xx\frac{dB_J(x)}{dx}.$$ (10) The $`J=5/2`$ result for $`s(M)`$ is illustrated in Fig. 1. The entropy per impurity vanishes for $`M=J=5/2`$ because there is a single configuration with $`_Im_I=N_IJ`$, and approaches $`\mathrm{ln}(2J+1)1.79`$ for $`M0`$. The mean-polarization of the localized spins at a given temperature and for a given orientation of the local moments is determined by minimizing the mean-field free energy $`F_{MF}(M)`$ $`=`$ $`k_BT\mathrm{ln}Z_{MF}(M)`$ (11) $`=`$ $`F_b(\stackrel{}{h}=N_{Mn}J_{pd}M\widehat{M})k_BTN_Is(M),`$ (12) with respect to $`M`$. Setting the derivative to zero gives $$\frac{ds(M)}{dM}=\frac{J_{pd}}{k_BTV}\frac{dF_b(h\widehat{M})}{dh}.$$ (13) Comparing with Eq. (9), it follows that $`F_{MF}(M)`$ is minimized by $`M=JB_J(g_L\mu _BH_{eff}J/k_BT)=JB_J(x_{eff})`$ where $$x_{eff}\frac{g_L\mu _BH_{eff}}{k_BT}=\frac{J_{pd}}{Vk_BT}\frac{dF_b(h\widehat{M})}{dh}.$$ (14) It follows that $`h=N_{Mn}J_{pd}M`$ is determined by solving the self-consistent equation $$h=N_{Mn}J_{pd}JB_J\left[x_{eff}(h)\right].$$ (15) Note that $$\frac{dF_b(h\widehat{M})}{dh}=\stackrel{}{S}_{tot}\widehat{M},$$ (16) where $`S_{tot}`$ is the total hole spin and the angle brackets indicate a thermal average for the non-interacting valence band system. Since the valence band system experiences an effective Zeeman coupling with strength proportional to $`h`$ and direction $`\widehat{M}`$, it is clear that the right hand side of Eq. (16) is negative in sign and that its magnitude increases monotonically with $`h`$, making it easy to solve Eq. (15) numerically. To simplify the calculations presented in subsequent sections, we take advantage of the fact that temperatures of interest are almost always considerably smaller than the itinerant carrier Fermi energy. This allows us to replace $`F_b(\stackrel{}{h})`$ by the ground state energy $`E_b(\stackrel{}{h})`$. Then, using Eq. (14) and Eq. (12), a single calculation of $`E_b(h\widehat{M})`$ over the range from $`h=0`$ to $`h=N_{Mn}J_{pd}J`$ may be used to determine the local-moment magnetization $`M(T)`$ and the free energy $`F(T)=F_{MF}(M(T))`$ at all temperatures. The mean-field theory critical temperature can be identified by linearizing the self-consistent equation at small $`h`$. We find that $$k_BT_c(\widehat{M})=\frac{J(J+1)}{3}\frac{N_{Mn}J_{pd}^2}{V}\frac{d^2F_b(h\widehat{M})}{dh^2}\|_{h=0}.$$ (17) The second derivative of the valence band free-energy with respect to field is proportional to its Pauli spin-susceptibility, which is in turn proportional to the valence band density of states at the Fermi energy, and to $`p^{1/3}`$ at small $`p`$. In the absence of strain, it follows from cubic symmetry that the right-hand-side of Eq. (17) is independent of $`\widehat{M}`$. Below the critical temperature however, the mean-field free energy does depend on $`\widehat{M}`$; this dependence is the magnetic anisotropy energy we wish to calculate. We will see that the dependence of the anisotropy energy on hole density is very different from that of the critical temperature. ## III Magnetic Anisotropy in the strong exchange coupling limit Our mean-field theory simplifies at low temperatures and, for the four-band model, simplifies further when $`h`$ is much larger than the characteristic energy scale of occupied Kohn-Luttinger states. A convenient typical energy scale is the $`h=0`$ hole Fermi energy $`ϵ_{F0}`$. For a given value of $`N_{Mn}J_{pd}`$ the largest value of $`h`$ is reached at $`T=0`$. Then, since $`H_{eff}`$ is always non-zero, $`x_{eff}\mathrm{}`$ and the solution to the mean-field equations is $`M=J`$, implying that $`h=N_{Mn}J_{pd}J`$ for every orientation $`\widehat{M}`$. At $`T=0`$ there is no entropic contribution to the free-energy and $$F_{MF}(T=0)=E_b(\stackrel{}{h}=N_{Mn}J_{pd}J\widehat{M}).$$ (18) We note in passing that the magnetization density at $`T=0`$ has contributions from the localized spins and the itinerant spins: $$M_s(T=0)=2\mu _BJN_{Mn}+\frac{\kappa }{3V}\stackrel{}{S}_{tot}\widehat{M},$$ (19) where $`\kappa `$ is an additional parameter of the Luttinger model. Because of the antiferromagnetic exchange interaction, the two terms here will tend to have opposite signs with the first term, which is independent of the hole density, typically very dominant. For $`hϵ_{F0}`$ two further simplifications occur. When the splitting of the hole bands by the effective Zeeman coupling is sufficiently large, or the hole density $`p`$ is sufficiently small, only the lowest energy hole band will be occupied. Furthermore, the effective Zeeman-term will dominate the mean-field single-particle Hamiltonian and, as we detail below, the envelope function spinor for this occupied hole state has a simple analytic expression. In this section we assume that the spin-orbit splitting energy $`\mathrm{\Delta }_{so}`$ is much larger than all other energies so that we can work with a four band model. More generally, the anisotropy will depend on $`h/\mathrm{\Delta }_{so}`$, even in the limit of small hole densities. To judge whether or not this limit can be achieved in practice, we estimate the Fermi energy of holes using the spherical approximation in which the doubly degenerate $`h=0`$ bulk heavy hole and light hole bands are parabolic with masses $`m_h=m/(\gamma _12\overline{\gamma })0.498m`$ and $`m_l=m/(\gamma _1+2\overline{\gamma })0.086m`$ respectively. An elementary calculation then gives $$ϵ_{F0}=\frac{\mathrm{}^2}{2m}\left(\frac{3\pi ^2p}{2}\right)^{2/3}\overline{\gamma }_0,$$ (20) where $`n=N_h/V`$ is the free-carrier density and $$\overline{\gamma }_0=\left[\frac{(\gamma _12\overline{\gamma })^{3/2}+(\gamma _1+2\overline{\gamma })^{3/2}}{2}\right]^{2/3}.$$ (21) For GaAs $`\overline{\gamma }_03.05`$: the Fermi energy is the same as that of a system with four identical effective mass $`m/\gamma _0`$ bands. Typical high $`T_c`$ $`\mathrm{Ga}_{1x}\mathrm{Mn}_x\mathrm{As}`$ ferromagnetic semiconductor samples have $`p0.1\mathrm{nm}^3`$ and $`N_{Mn}1.0\mathrm{nm}^3`$ ($`x0.05`$), although these parameters can presumably be varied widely. Choosing a $`J_{pd}`$ value in the mid-range of estimated values ($`0.006\mathrm{Ry}\mathrm{nm}^3`$) these parameters imply that $`h0.015\mathrm{Ry}`$ and $`ϵ_{F0}0.01\mathrm{Ry}`$. $`h`$ is neither large compared to $`ϵ_{F0}`$, nor small compared to $`\mathrm{\Delta }_{so}`$. Thus, the simple expressions discussed in this section are not accurate for current high $`T_c`$ systems. As our ability to engineer materials improves it should, however, be possible to grow samples which are in the limit discussed here. Since $`h`$ is comparable to $`ϵ_{F0}`$, we know, even before performing detailed calculations, that valence band quasiparticle spectra in paramagnetic and ferromagnetic states will differ qualitatively. In the large $`h`$ limit the lone occupied spinor at each wavevector $`\stackrel{}{k}`$ will be the member of the $`j=3/2`$ manifold for which the total angular momentum is aligned in the direction $`\widehat{M}`$; the origin of the minus sign here is the antiferromagnetic nature of the interaction between localized spins and hole spins. Explicit expressions for the expansion of such a spin coherent states in terms of the eigenstates of $`j_z`$ are known: $`\|\widehat{M}`$ $`=`$ $`u^3\|3/2+\sqrt{3}[u^2v\|1/2+uv^2\|1/2]`$ (22) $`+`$ $`v^3\|3/2>.`$ (23) In Eq. (23) $`u=i\mathrm{exp}(i\varphi /2)\mathrm{sin}(\theta /2)`$ and $`v=i\mathrm{exp}(i\varphi /2)\mathrm{cos}(\theta /2)`$ where $`\theta `$ and $`\varphi `$ are the spherical coordinates which specify the unit vector $`\widehat{M}`$. In the large $`h/ϵ_{F0}`$ limit, the band term in the single-particle Hamiltonian may be treated using first order perturbation theory. Taking the expectation value of the Kohn-Luttinger Hamiltonian in the spin coherent state we find that $`ϵ(\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{h}{2}}+\widehat{M}\|H_L(\stackrel{}{k})\|\widehat{M}`$ (24) $``$ $`{\displaystyle \frac{h}{2}}+{\displaystyle \frac{\mathrm{}^2k^2}{2m}}\gamma (\widehat{M},\widehat{k}).`$ (25) The first term on the right hand side of Eq. (25) reflects the spin coherent state property, $`\widehat{M}s\|\widehat{M}=\|\widehat{M}/2`$. In Eq. (25) we have noted that for any $`\widehat{M}`$ and any $`\widehat{k}`$, the dependence of hole energy on $`k=\|\stackrel{}{k}\|`$ is quadratic. Using this property, it follows that the Fermi energy $$ϵ_F(\widehat{M})=\frac{\mathrm{}^2}{2m}(6\pi ^2p)^{2/3}\overline{\gamma }(\widehat{M}),$$ (26) and that the ferromagnetic ground state energy density is $$\frac{E_b(\stackrel{}{M})}{V}=\frac{hp}{2}+\frac{3}{5}pϵ_F(\widehat{M}),$$ (27) where $$\overline{\gamma }(\widehat{M})\left[\frac{d\widehat{k}}{4\pi }(\gamma (\widehat{M},\widehat{k}))^{3/2}\right]^{2/3}.$$ (28) In analogy with the $`h=0`$ quantity $`\overline{\gamma }_0`$ defined in Eq. (21), $`\overline{\gamma }(\widehat{M})`$ is an average of the band energy curvature over reciprocal space directions $`\widehat{k}`$, with the smaller values of $`\gamma (\widehat{M},\widehat{k})`$ weighted more heavily. Note that the factor $`3\pi ^2/2`$ in Eq. (20) is replaced by $`6\pi ^2`$ in Eq. (26) because only one band is occupied in this limit, instead of the four which are occupied at $`h=0`$. $`m/\overline{\gamma }(\widehat{M})`$ may be thought of as a spin-orientation dependent effective mass, which is readily evaluated as a function of $`\widehat{M}`$, given the Luttinger parameters of any material. Although the magnetic condensation energy has a term proportional to $`J_{pd}`$, only the band-energy contributes to the $`\widehat{M}`$ dependence of the ferromagnetic ground state energy. The magnetic anisotropy energy in this limit is independent of $`J_{pd}`$ and proportional to the hole density $`p`$ to the 5/3 power. We have evaluated $`\gamma (\widehat{M})`$ as a function of angle for the Luttinger parameters of GaAs and InAs. As discussed in more detail later, we always find that magnetic anisotropy in the absence of strain is well described by a cubic harmonic expansion truncated at sixth order, an approximation commonly used in the literature on magnetic materials. The corresponding cubic harmonic expansion for $`\overline{\gamma }(\widehat{M})`$ is $`\overline{\gamma }(\widehat{M})`$ $`=`$ $`\overline{\gamma }(100)+\gamma _1^{ca}(\widehat{M}_x^2\widehat{M}_y^2+\widehat{M}_y^2\widehat{M}_z^2+\widehat{M}_x^2\widehat{M}_z^2)`$ (29) $`+`$ $`\gamma _2^{ca}\widehat{M}_x^2\widehat{M}_y^2\widehat{M}_z^2.`$ (30) In our calculations we find that $`\overline{\gamma }(\widehat{M})<\gamma _1`$ for all directions $`\widehat{M}`$. This property reflects the fact that small curvature (large Fermi wavevector) directions are weighted more highly in calculating the total hole energy. For both InAs and GaAs we find that the dominant fourth order cubic anisotropy coefficient, $`\gamma _1^{ca}<0`$, indicating nickel-type anisotropy with easy axes along the $`111`$ cube diagonal directions. At a qualitative level, the source of the higher total hole energy when the moment orientation is along a $`001`$ (cube edge) direction is easy to understand. With such a moment orientation, the occupied hole orbital has $`j_z=3/2`$ and energy dispersion given by $`H_{hh}(\stackrel{}{k})`$ in Eq. (77). It follows that $`\gamma (\widehat{M}=100,\widehat{k})`$ has the relatively large value $`\gamma _1+\gamma _2`$ for all orientations of $`\widehat{k}`$ in the $`\widehat{x}\widehat{y}`$ plane. These large values of $`\gamma (\widehat{M},\widehat{k})`$ are important in the average and cause the average over $`\widehat{k}`$ to reach its maximum for this orientation of $`\widehat{M}`$. The cubic magnetic anisotropy energy coefficients in this limit are given by $$K_i^{ca}=+\frac{3}{5}p\frac{\mathrm{}^2}{2m}(6\pi ^2p)^{2/3}\gamma _i^{ca}.$$ (31) Values of $`\gamma _i^{ca}`$ for GaAs and InAs are listed in Table I. We will see later that these expressions apply up to $`p10^{18}\mathrm{cm}^3`$. Inserting this value for the hole density gives coefficients $`2\mathrm{k}\mathrm{J}\mathrm{m}^3`$ for GaAs host material and $`4\mathrm{k}\mathrm{J}\mathrm{m}^3`$ for InAs host materials; magnetic anisotropy is twice as strong in InAs in the strain free case. These anisotropy energy coefficients are not so much smaller than those of the cubic metallic transition metal ferromagnets, despite the much higher carrier densities in the metallic case. We will see below that the scale of the semiconductor magnetic anisotropy energy does not change as the carrier density increases from $`10^{18}\mathrm{cm}^3`$ to $`10^{21}\mathrm{cm}^3`$. The relatively large anisotropy energies occur despite the fact that the saturation moments $`M_s`$ of $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ ferromagnets are more than an order of magnitude smaller than their cubic metal counterparts. It follows from these values that the magnetic hardness parameters of the $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ ferromagnets, $$\kappa \left[\frac{K_1^{ca}}{\mu _0M_s^2}\right]^{1/2},$$ (32) will typically be larger than one. This is unusual in cubic materials and occurs because spin-orbit coupling has a much stronger influence on semiconductor valence bands than on transition metal $`d`$-bands. As we discuss at length in Section V, magnetic anisotropy does not continue to increase rapidly with hole density once two or more bands are occupied in the metallic state. In current high $`T_c`$ samples, we will find that several bands are always occupied, even at zero temperature. The simple limit discussed in this section demonstrates that anisotropy energies, $`T=0`$ saturation moments, and critical temperatures will have radically different dependencies on engineerable parameters. ## IV Strain Dependence of Magnetic Anisotropy: Low Hole Density Limit Because of the low soluability of Mn in III-V semiconductors, $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ materials with $`x`$ large enough to produce cooperative magnetic effects cannot be obtained by equilibrium growth. The MBE growth techniques which have been successfully developed produce $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ films whose lattices are locked to those of their substrates. X-ray diffraction studies have established that the resulting strains are not relaxed by dislocations or other defects, even for thick films. Strains in the $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ film break the cubic symmetry assumed in the previous section. Fortunately, the influence of MBE growth lattice-matching strains on the hole bands of cubic semiconductors is well understood. For the $`001`$ growth direction used to create $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ films, strain generates a purely diagonal contribution to the four band single-particle envelope function Hamiltonian in the representation we use in this paper, adding contributions $`\delta ϵ_h`$ and $`\delta ϵ_l`$ respectively to $`j_z=\pm 3/2`$ heavy hole and $`j_z=\pm 1/2`$ light hole entries. The energy shifts are related to the lattice strains by $$\delta ϵ_h=\frac{e_0}{C_{11}}\left[2a_1(C_{11}C_{12})\frac{a_2}{2}(C_{11}+2C_{12})\right],$$ (33) $$\delta ϵ_l=\frac{e_0}{C_{11}}\left[2a_1(C_{11}C_{12})+\frac{a_2}{2}(C_{11}+2C_{12})\right],$$ (34) where $`e_0`$ is the in-plane strain produced by the substrate-film lattice mismatch: $$e_0=\frac{a_Sa_F}{a_F}.$$ (35) In Eqs. (3335), the $`C_{ij}`$ are the elastic constants of the unstrained $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ film, which we will assume to be identical to those of the host III-V material, $`a_S`$ is the lattice constant of the substrate on which the $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ film is grown, $`a_F`$ is the unstrained lattice constant of bulk $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$, and $`a_1`$ and $`a_2`$ are phenomenological deformation potentials whose values for common III-V semiconductors are known. For the six band model the strain Hamiltonian includes off-diagonal elements and is given up to a constant times the unit matrix by $$H_{\mathrm{strain}}=\mathrm{\Gamma }e_0\left(\begin{array}{cccccc}0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& \frac{1}{\sqrt{2}}\\ 0& 0& 1& 0& \frac{1}{\sqrt{2}}& 0\\ 0& 0& 0& 0& 0& 0\\ 0& 0& \frac{1}{\sqrt{2}}& 0& \frac{1}{2}& 0\\ 0& \frac{1}{\sqrt{2}}& 0& 0& 0& \frac{1}{2}\end{array}\right),$$ (36) where $`\mathrm{\Gamma }=ϵ_lϵ_h=a_2e_0(C_{11}+2C_{12})/C_{11}`$. For GaAs and InAs, $`\mathrm{\Gamma }=0.2382`$ Ry and $`0.2762`$ Ry respectively. As in the previous section, we can derive an explicit expression for the strain contribution to the magnetic anisotropy energy when $`hϵ_{F0}`$ and $`h\mathrm{\Gamma }`$, allowing band and strain terms to be treated as a perturbative correction to the effective Zeeman coupling, and $`h\mathrm{\Delta }_{so}`$, allowing a four-band model to be used. In this way we obtain $$\frac{E_{strain}(\widehat{M})}{V}=\frac{p}{4}[3\delta ϵ_h+\delta ϵ_l]+\frac{3p\mathrm{cos}^2\theta }{4}[\delta ϵ_h\delta ϵ_l].$$ (37) Strain produces a uniaxial contribution to the magnetic anisotropy of $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ films, that favors orientations along the growth direction when strain shifts heavy holes down relative to the light holes and orientations in the plane in the opposite circumstance. Using Eq. (33) and Eq. (34) and deformation potential elastic constant values, we find that a contribution to the energy density given by $`K_{strain}\mathrm{sin}^2(\theta )`$ where the uniaxial anisotropy constant is $`K_{strain}=0.36\mathrm{Ry}e_0p`$ for GaAs and $`K_{strain}=0.41\mathrm{Ry}e_0p`$ for InAs. Since $`a_2<0`$, compressive strain ($`e_0<0`$) lowers the heavy hole energy relative to the light hole energy and favors moment orientations in the growth direction while tensile strain ($`e_0>0`$) favors moment orientations perpendicular to the growth direction. Since the lattice constant of $`\mathrm{Ga}_{1x}\mathrm{Mn}_xAs`$ is larger than that of GaAs, while that of $`\mathrm{In}_{1x}\mathrm{Mn}_xAs`$ is smaller than that of InAs, $`\mathrm{Ga}_{1x}\mathrm{Mn}_xAs`$ on GaAs is under compressive strain and $`\mathrm{In}_{1x}\mathrm{Mn}_xAs`$ on InAs is under tensile strain. Assumming Vergard’s law, $`e_0=.0004`$ and $`e_0=.0028`$ for InAs and GaAs respectively at $`x=0.05`$. For $`p=10^3\mathrm{nm}^3`$, a density for which this low-density result is still reasonably reliable, we find that $`K_{strain}=0.36\mathrm{kJm}^3`$ for InAs and $`K_{strain}=2.2\mathrm{kJm}^3`$ for GaAs. At these densities, strain and cubic band contributions are comparable in GaAs, but the latter contribution is dominant in the InAs case. Strain-induced anisotropy energies can be modified by growing the magnetic film on relaxed buffer layers. This effect has been demonstrated by Ohno et al. who showed that the easy axis for $`\mathrm{Ga}_{1x}\mathrm{Mn}_xAs`$ films changes from in-plane to growth-direction when the substrate is changed from GaAs to relaxed (In,Ga)As buffer layers. For $`15\%`$ In, the magnetic film strain changes from compressive $`e_0=0.0028`$ to tensile $`e_0=0.0077`$ when this change is made. We note that the sense of this change is opposite to that predicted by our large $`h`$, small $`p`$ analytic result which predicts that compressive strains favor growth direction orientations. Similarly, $`\mathrm{In}_{1x}\mathrm{Mn}_xAs`$ on InAs is under a small tensile strain but is observed to have a growth direction easy axis. (Recall that in the large $`h`$ limit the band anisotropy makes the growth direction the hard axis.) This distinction may be taken as an experimental proof that several hole bands are occupied in the magnetic ground state of these materials. The strain anisotropy energy does change sign at smaller values of $`h`$ partly because the first band to be depopulated has primarily heavy-hole character. We will see in the next section that the mean-field does predict the correct sign for the strain anisotropy energy at experimental hole densities. We conclude from the present considerations that strain can play a strong role in band-structure-engineering of $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ ferromagnet magnetic properties. ## V Partially Polarized Hole Band States: Magnetic Anisotropy in the General Case The results for $`F_b(h\widehat{M})`$ discussed in the previous section become accurate when the effective Zeeman coupling $`h`$ is large enough to reduce the number of occupied hole bands to one, and the Fermi energy remains safely smaller than the spin-orbit splitting. The situation is much more complicated at smaller $`h`$ and larger hole densities. Then several bands are occupied, even at $`T=0`$, and these usually give competing contributions to the magnetic anisotropy. It is not true in general that band and strain contributions to the magnetic anisotropy are simply additive. We expect that it will eventually be possible to realize a broad range of material parameters, and hence a broad range of $`h`$ values, in $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ ferromagnets. The range of possibilities is immense and accurate modeling of a particular sample will require accurate values for $`p`$, $`J_{pd}`$, and $`N_{Mn}`$ for that material. In this section we discuss a series of illustrative calculations, starting with ones performed using a four-band model of strain-free $`\mathrm{Ga}_{1x}\mathrm{Mn}_x\mathrm{As}`$ at hole density $`p=0.1\mathrm{nm}^3`$. The valence band energy density is $$\frac{E_b}{V}=\frac{1}{V}\underset{𝐤}{}\underset{j=1}{\overset{N_b}{}}\theta (ϵ_Fϵ_j(𝐤,\stackrel{}{h}))ϵ_j(𝐤,\stackrel{}{h}),$$ (38) where $`ϵ_F`$ is the Fermi energy, $`\stackrel{}{k}`$ is the Bloch wavevector, and the mean-field-theory quasiparticle energies $`ϵ_j(𝐤,\stackrel{}{h})`$ are eigenvalues of the $`N_b\times N_b`$ single-particle Hamiltonian ($`N_b`$ is the number of bands included in the envelope function Hamiltonian) $$H_b=H_L+\stackrel{}{h}\stackrel{}{s}+H_{\mathrm{strain}}.$$ (39) For a strain free model we set $`H_{\mathrm{strain}}=0`$. In Fig. 2 we plot the calculated spin-polarization per hole as a function of $`h`$ for a growth direction field orientation; recall that this quantity can be obtained by differentiating the energy per hole with respect to $`h`$ and that the effective field seen by the localized spins is obtained by multiplying this quantity by $`J_{pd}p`$. The hole spin polarization increases linearly at small $`h`$ with a slope proportional to the valence band Pauli susceptibility. We see that for the hole density of Fig. 2, complete hole spin-polarization is approached only at values of $`h`$ comparable to or larger than the spin-orbit splitting of GaAs, so that the four band model is not physical in this regime. For any given model, and a given moment orientation, a single calculation of this type provides all the microscopic information required to solve the mean-field equations at all temperatures. For a fixed values of $`J_{pd}`$, $`N_{Mn}`$, and $`p`$, $`h`$ must be evaluated as a function of temperature by solving the self-consistent mean-field equation, Eq. (15), and using numerical results like those plotted in Fig. 2. Results for $`h(T)`$ calculated for $`J_{pd}=0.15\mathrm{eVnm}^3`$, $`N_{Mn}=1\mathrm{n}\mathrm{m}^3`$, and $`p=0.1\mathrm{nm}^3`$ are illustrated in Fig. 3 for high symmetry moment orientation directions. For these parameters the critical temperature $`T_c`$ ($`h(T)=0`$ for $`T>T_c`$) is $`100\mathrm{K}`$, in rough agreement with experiment. There are, however, other choices of parameters which are also consistent with the measured critical temperature. In addition, as we discuss further in the next section, it is not clear that this level of theory should always yield accurate results for $`T_c`$. At $`T=0`$, $`h`$ reaches its maximum value, $`J_{pd}N_{Mn}J0.0275\mathrm{Ry}`$. The dependence of energy on magnetization orientation at this value of $`h`$ is illustrated in Fig. 4 and compared with the cubic harmonic expansion truncated at 6th order. The coefficients of this expansion are fixed by energy per volume calculations in $`100`$, $`110`$, and $`111`$ directions. In Fig. 4, and in all other cases we have checked, the truncated cubic harmonic expansion is very accurate. It is therefore sufficient to evaluate the energy per volume in the high-symmetry directions and to use $`K_1^{ca}`$ $`=`$ $`{\displaystyle \frac{4(E_b110E_b100)}{V}}`$ (40) $`K_2^{ca}`$ $`=`$ $`{\displaystyle \frac{27E_b11136E_b110+9E_b100}{V}}.`$ (41) Results for $`K_i^{ca}(h)`$ obtained for the four band model in this way are summarized in Fig. 5. These results can be combined with those in Fig. 3 to obtain the model’s cubic anisotropy coefficients as a function of temperature, hole density and $`h`$. Here we see explicitly the anisotropy reversals which commonly accompany hole band depopulations. We note in Fig. 5 that the analytic results of Section III are recovered only for very large values of $`h`$ at this density. The valence band Fermi surfaces of $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ ferromagnets will be strongly dependent on both temperature and moment direction orientation. Four-band model Fermi surface intersections with the $`k_z=0`$ plane are illustrated in Figs. 68 for $`p=0.1\mathrm{nm}^3`$ at $`h=0`$, $`h=0.01\mathrm{Ry}`$, and $`h=0.0275\mathrm{Ry}`$ respectively. All figures are for moments oriented in the $`100`$ direction. The largest of these three values of $`h`$ is the $`T=0`$ effective field ($`J_{pd}N_{Mn}J`$) for $`x=0.05`$ and $`J_{pd}`$ at the high end of literature estimates. ($`J_{pd}=0.15\mathrm{eVnm}^3`$.) These three figures represent the mean-field-theory Fermi surfaces at three different temperatures. In the spherical approximation, the $`h=0`$ Fermi energy at this hole density is $`ϵ_{F0}0.01\mathrm{Ry}`$, so a strong distortion of the bands is expected in the magnetic state. At $`h=0`$, the hole bands occur in degenerate pairs; we refer to the two less dispersive bands which occupy the larger area in $`\widehat{k}`$ space as heavy-hole bands, although this terminology lacks precise meaning in the general case. As $`h`$ increases, both heavy and light hole bands split. For small $`h`$, the minority-spin heavy-hole band occupation decreases rapidly and all other band occupations increase. The heavy-hole minority spin band is completely depopulated for $`h0.04\mathrm{Ry}`$. Once this band is empty, the light-hole minority-spin Fermi radii begin to shrink rapidly. As seen in Fig. 8, this second band is nearly depopulated at $`h=0.0275\mathrm{Ry}`$. At still stronger fields, the majority-spin light-hole band is depopulated and the single-band limit addressed in preceding sections is achieved. For $`p=0.1\mathrm{nm}^3`$, the single-band limit is achieved only at $`h`$ values for which the four-band model breaks down. Finally we turn to a series of illustrative calculations intended to closely model the ground state of $`\mathrm{Ga}_{.95}\mathrm{Mn}_{.05}\mathrm{As}`$. For this Mn density and the smaller values of $`J_{pd}`$ favored by recent estimates, $`h(T=0)=J_{pd}N_{Mn}J0.01\mathrm{Ry}`$. This value of $`h`$ is not so much smaller than the spin-orbit splitting parameter in GaAs ($`\mathrm{\Delta }_{so}=0.025`$ Ry), so that accurate calculations require a six band model. Even with $`x`$ fixed, our calculations show that the magnetic anisotropy of $`\mathrm{Ga}_{.95}\mathrm{Mn}_{.05}\mathrm{As}`$ ferromagnets is strongly dependent on both hole density and strain. The hole density can be varied by changing growth conditions or by adding other dopants to the material, and strain in a $`\mathrm{Ga}_{.95}\mathrm{Mn}_{.05}\mathrm{As}`$ film can be altered by changing substrates as discussed previously. The cubic anisotropy coefficients (in units of energy per volume) for strain-free material are plotted as a function of hole density in the inset of Fig. 9; the main plot shows the coefficients in units of energy per particle. Over the density range $`p<0.05\mathrm{nm}^3`$, four and six band models are in good agreement. We see from this result that the asymptotic low density region where the anisotropy energy varies as $`p^{5/3}`$ holds only for $`p<0.005\mathrm{nm}^3`$ at this value of $`h`$. The easy axis is nearly always determined by the leading cubic anisotropy coefficient $`K_1^{ca}`$, except near values of $`p`$ where this coefficient vanishes. As a consequence the easy-axis in strain free samples is almost always either along one of the cube edge directions ($`K_1^{ca}>0`$), or along one of the cube diagonal directions ($`K_1^{ca}<0`$). Transitions in which the easy axis moves between these two directions occur twice over the range of hole densities studied. (Similar transitions occur as a function of $`h`$, and therefore temperature, for fixed hole density.) Near the hole density 0.01 nm<sup>-3</sup>, both anisotropy coefficients vanish and a fine-tuned isotropy is achieved. The slopes of the anisotropy coefficient curves vary as the number of occupied bands increases from $`1`$ to $`4`$ with increasing hole density. This behavior is clearly seen from the correlation between oscillations of the anisotropy coefficients and onsets of higher band occupations, plotted in Fig. 10. Six-band model Fermi surfaces are illustrated in Figs. 11-13 by plotting their intersections with the $`k_z=0`$ plane at $`p=0.1\mathrm{nm}^3`$ for the cases of $`100`$, $`110`$, and $`111`$ ordered moment orientations. Comparing Fig. 7 and Fig. 11, which differ only in the band model employed, we see that there is a marked difference between the majority-spin heavy hole bands in four and six-band cases. For the six band model, quasiparticle dispersion is particularly slow, leading to large Fermi radii along the $`110`$ directions. The large and more anisotropic mass is a consequence of mixing with the split-off hole bands. As we see from Figs. 12 and 13, this effect occurs for all ordered moment orientations, although the details of the small minority band Fermi surface projections change markedly. The dependence of quasiparticle band structure on ordered moment orientation, apparent in comparing these figures, should lead to large anisotropic magnetoresistance effects in $`\mathrm{III}_{1x}\mathrm{Mn}_x\mathrm{V}`$ ferromagnets. We also note that in the case of cube edge orientations, the Fermi surfaces of different bands intersect. This property could have important implications for the decay of long-wavelength collective modes . In Fig. 14 we present mean-field theory predictions for the strain-dependence of the anisotropy energy at $`h=0.01\mathrm{Ry}`$ and hole density $`p=0.35\mathrm{nm}^3`$. According to our calculations, the easy axes in the absence of strain are along the cube edges in this case. This calculation is thus for a hole density approximately three times smaller than the Mn density, as indicated by recent experiments. The relevant value of $`e_0`$ depends on the substrate on which the epitaxial $`\mathrm{Ga}_{.95}\mathrm{Mn}_{.05}\mathrm{As}`$ film is grown, as discussed in Section IV. The most important conclusion from Fig. 14 is that strains as small as $`1\%`$ are sufficient to completely alter the magnetic anisotropy energy landscape. For example for (Ga,Mn)As on GaAs, $`e_0=.0028`$ at $`x=0.05`$, the anisotropy has a relatively strong uniaxial contribution even for this relatively modest compressive strain, which favors in-plane moment orientations, in agreement with experiment. A relatively small ($``$1 kJ m<sup>-3</sup>) residual plane-anisotropy remains which favors $`110`$ over $`100`$. For $`x=0.05`$ (Ga,Mn)As on a $`x=0.15`$ (In,Ga)As buffer the strain is tensile, $`e_0=0.0077`$, and we predict a substantial uniaxial contribution to the anisotropy energy which favors growth direction orientations, again in agreement with experiment. For the tensile case, the anisotropy energy changes more dramatically than for compressive strains due to the depopulation of higher subbands, as shown in Fig. 15. At large tensile strains, the sign of the anisotropy changes emphasizing the subtlety of these effects and the latitude which exists for strain-engineering of magnetic properties. ## VI Discussion We first comment on the implications of the considerations described in this paper for the interpretation of current experiments. The hysteretic effects which reflect magnetic anisotropy have been studied most extensively for the highest $`T_c`$ samples currently available. These mean-field-theory predictions depend on three phenomenological parameters, $`N_{Mn}`$ which is sample dependent but accurately known, $`J_{pd}`$ which should be nearly universal for a given III-V host compound but is less accurately known, and the hole density $`p`$ which is sample dependent and not accurately known. Values of $`J_{pd}`$ and $`p`$ must be inferred from experiment, sometimes by comparison with theoretical pictures which are not yet fully developed. The reliability of $`J_{pd}`$ and $`p`$ estimates is improving and presumably will continue to improve. It now seems clear that the values of $`J_{pd}`$ and $`p`$ in current high $`T_c`$ samples are such that several hole bands are partially occupied in the ferromagnetic ground state. In this case, we see from Fig. 5 that our calculations predict cube edge easy axes which include the growth direction. The anisotropy energies are typically $`10^6\mathrm{Rynm}^31\mathrm{k}\mathrm{J}\mathrm{m}^3`$. Similar anisotropy energies are produced by strains as small as $`e_00.001`$ and more typical strains produce larger anisotropy energies. An important conclusion from this work is that strain contributions to the anisotropy will not normally be negligible. We believe that the in-plane easy-axis observed in $`\mathrm{Ga}_{1x}\mathrm{Mn}_x\mathrm{As}`$ films grown on GaAs is a consequence of compressive strain in the magnetic film which dominates the cubic band anisotropy energy. When $`\mathrm{Ga}_{1x}\mathrm{Mn}_x\mathrm{As}`$ is grown on (In,Ga)As, the lattice-matching strain is tensile, reinforcing the growth direction easy-axis anisotropy of strain free samples. (As discussed earlier, the signs of both strain and cubic band contributions to strain change when the light hole bands are depopulated.) We note that in our calculations, the cubic band anisotropy is almost always dominated by the fourth cubic harmonic coefficient. Given this, it follows that only the cube edge easy axes, which includes the growth direction axis, and cube diagonal axes, which are not in the film plane, are possible without strain. Since the local moments are fully polarized in the ferromagnetic ground state, it is easy to estimate the saturation moment $`M_sN_{Mn}g_L\mu _BJ`$, leading to the relatively small numerical value $`\mu _0M_s0.05\mathrm{T}`$. It follows that the growth direction orientation magnetostatic energy $`\mu _0M_s^20.1\mathrm{kJm}^3`$, considerably smaller than the magnetocrystalline anisotropy coefficient. Even when several hole bands are occupied with competing spin-orbit interactions, these materials have relatively large magnetic hardness parameters. Unlike the case of metallic thin film ferromagnets which have much larger saturation moments, the magnetostatic shape anisotropy plays a minor role in the dependence of total energy on moment orientation. The small saturation moment, will also tend to lead to large domain sizes and square easy-axis direction hysterisis loops, as seen in experiment. Coercivities can be estimated from the anisotropy fields defined by $$\mu _0H_a\mu _0M_s\frac{K}{\mu _0M_s^2}.$$ (42) For hard magnetic materials, anisotropy fields are much larger than saturation magnetizations. The itinerant field places an upper bound on and is expected to scale with the coercivity. Our calculations suggest that the coercivity in ferromagnetic samples with a single partially occupied hole band will be immensely larger than the coercivity of current samples. Such samples could be fabricated, for example, by adding donors such as Si to current samples, further compensating the Mn acceptors. According to mean-field-theory, this modification in the sample preparation procedure will lower the ferromagnetic critical temperature, and at the same time, increase the anisotropy energy. Finally we conclude with a few words of caution. This theory of magnetic anisotropy has three principle limitations: i) it is based on a mean-field theory description of the exchange interaction between localized spins and valence band holes, ii) it neglects hole-hole interactions, and iii) it doesn’t account for disorder scattering of the itinerant holes. Confronting these weaknesses would in each case considerably complicate the theory and we feel it is appropriate to seek progress by comparing the present relatively simple theory with experiment. Nevertheless it is worthwhile to speculate on where and how the theory may be expected to fail. Mean-field theory should be reliable when the range of the hole mediated interaction between localized spins is long and the spin-stiffness parameter that characterizes the energy of long-wavelength spin-fluctuations is sufficiently large. Considerations of this type suggest that mean-field theory will fail at high temperatures unless the ratio of the hole density to the localized spin density is small and the field $`h`$ is not too large compared to the Fermi energy. Since $`p`$ is typically smaller than $`N_{Mn}`$ because of anti-site defects in low-temperature MBE growth samples, mean-field theory is likely to be reasonable at least at $`T=0`$ in many (III,Mn)V ferromagnets. Hole-hole interactions will clearly tend to favor the ferromagnetic state by countering the band-energy cost of the spin polarization. Because of strong spin-orbit coupling in the valence band, estimates based on many-body calculations for electron gas systems may be of little use in estimating the importance of this effect more quantitatively. Work is currently in progress which should shed more light on this matter. Nevertheless, it seems likely that these interactions will not have an overriding importance, at least in large hole density samples. Finally we come to disorder. It is clear from experiment, that disorder does not have a qualitative impact on free-carrier mediated ferromagnetism even when those free carriers have been localized by a random disorder potential. It seems likely that disorder will destroy the ferromagnetic state, only when the localization length becomes comparable to the distance between localized spins. On the other hand, since elastic disorder scattering will mix band states with different orientations on the Fermi surface, it also seems clear that a reduction in magnetic anisotropy energy must result. Indeed the coercivities that follow from our anisotropy energy results appear to be larger than what is observed. As far as we are aware, no theory of this effect exists at present. ###### Acknowledgements. We gratefully acknowledge helpful interactions with W.A. Atkinson, T. Dietl, J. Furdyna, J.A. Gaj, J. König, B.-H. Lee, E. Miranda, Hsiu-Hau Lin, and Hideo Ohno. Work at the University of Oklahoma was supported by the NSF under grant No. EPS-9720651 and a grant from the Oklahoma State Regents for Higher Education. Work at Indiana University was performed under NSF grants DMR-9714055 and DGE-9902579. Work at the Institute of Physics ASCR was supported by the Grant Agency of the Czech Republic under grant 202/98/0085. AHM gratefully acknowledges the hospitality of UNICAMP where project work was initiated. ## Appendix In the literature, different representations are used for the four band and six band model Kohn-Luttinger Hamiltonians. In the interest of completeness and clarity, this appendix specifies the expressions on which our detailed calculations are based. Detailed derivations of $`𝐤𝐩`$ perturbation theory for cubic semiconductors can be found elsewhere. The $`k=0`$ states at the top valence band have $`p`$-like character and can be represented by the $`l=1`$ orbital angular momentum eigenstates $`\|m_l`$. In the coordinate representation we can write $`𝐫\|m_l=1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f(r)(x+iy)`$ (43) $`𝐫\|m_l=0`$ $`=`$ $`f(r)z`$ (44) $`𝐫\|m_l=1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}f(r)(xiy).`$ (45) The Kohn-Luttinger Hamiltonian for systems with no spin-orbit coupling, $`_L`$, is written in the representation of the following combinations of $`\|m_l`$ $`\|X`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\|m_l=1\|m_l=1\right)`$ (46) $`\|Y`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2}}}\left(\|m_l=1+\|m_l=1\right)`$ (47) $`\|Z`$ $`=`$ $`\|m_l=0.`$ (48) It reads $$_L=\left(\begin{array}{ccc}Ak_x^2+B(k_y^2+k_z^2)& Ck_xk_y& Ck_xk_z\\ Ck_yk_x& Ak_y^2+B(k_x^2+k_z^2)& Ck_yk_z\\ Ck_zk_x& Ck_zk_y& Ak_z^2+B(k_x^2+k_y^2)\end{array}\right),$$ (49) where $`A`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}(\gamma _1+4\gamma _2),`$ (50) $`B`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}(\gamma _12\gamma _2),`$ (51) $`C`$ $`=`$ $`{\displaystyle \frac{3\mathrm{}^2}{m}}\gamma _3,`$ (52) $`m`$ is the bare electron mass, and $`\gamma _1`$, $`\gamma _2`$, and $`\gamma _3`$ are the phenomenological Luttinger parameters. To include spin-orbit coupling we use the basis formed by total angular momentum eigenstates $`\|j,m_j`$: $`\|1`$ $``$ $`\|j=3/2,m_j=3/2`$ (53) $`\|2`$ $``$ $`\|j=3/2,m_j=1/2`$ (54) $`\|3`$ $``$ $`\|j=3/2,m_j=1/2`$ (55) $`\|4`$ $``$ $`\|j=3/2,m_j=3/2`$ (56) $`\|5`$ $``$ $`\|j=1/2,m_j=1/2`$ (57) $`\|6`$ $``$ $`\|j=1/2,m_j=1/2`$ (58) The basis (58) is related to the orbital angular momentum ($`m_l=1,0,1`$) and spin ($`\sigma =,`$) eigenstates by $`\|1`$ $`=`$ $`\|m_l=1,`$ (59) $`\|2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\|m_l=1,+\sqrt{{\displaystyle \frac{2}{3}}}\|m_l=0,`$ (60) $`\|3`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\|m_l=1,+\sqrt{{\displaystyle \frac{2}{3}}}\|m_l=0,`$ (61) $`\|4`$ $`=`$ $`\|m_l=1,`$ (62) $`\|5`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\|m_l=0,+\sqrt{{\displaystyle \frac{2}{3}}}\|m_l=1,`$ (63) $`\|6`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\|m_l=0,\sqrt{{\displaystyle \frac{2}{3}}}\|m_l=1,`$ (64) or $`\|1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\|X,+i\|Y,\right)`$ (65) $`\|2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}\left(\|X,i\|Y,\right)+\sqrt{{\displaystyle \frac{2}{3}}}\|Z,`$ (66) $`\|3`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}\left(\|X,+i\|Y,\right)+\sqrt{{\displaystyle \frac{2}{3}}}\|Z,`$ (67) $`\|4`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\|X,i\|Y,\right)`$ (68) $`\|5`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\left(\|X,+i\|Y,\right){\displaystyle \frac{1}{\sqrt{3}}}\|Z,`$ (69) $`\|6`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\left(\|X,i\|Y,\right)+{\displaystyle \frac{1}{\sqrt{3}}}\|Z,`$ (70) The six band model Kohn-Luttinger Hamiltonian, $`H_L`$, in the representation of vectors (58) is $$H_L=\left(\begin{array}{cccccc}_{hh}& c& b& 0& \frac{b}{\sqrt{2}}& c\sqrt{2}\\ c^{}& _{lh}& 0& b& \frac{b^{}\sqrt{3}}{\sqrt{2}}& d\\ b^{}& 0& _{lh}& c& d& \frac{b\sqrt{3}}{\sqrt{2}}\\ 0& b^{}& c^{}& _{hh}& c^{}\sqrt{2}& \frac{b^{}}{\sqrt{2}}\\ & & & & & \\ \frac{b^{}}{\sqrt{2}}& \frac{b\sqrt{3}}{\sqrt{2}}& d^{}& c\sqrt{2}& _{so}& 0\\ c^{}\sqrt{2}& d^{}& \frac{b^{}\sqrt{3}}{\sqrt{2}}& \frac{b}{\sqrt{2}}& 0& _{so}\end{array}\right)$$ (71) In the matrix (71) we highlighted the four band model Hamiltonian. The Kohn-Luttinger eigenenergies are measured down from the top of the valence band, i.e., they are hole energies and we use the following notation: $`_{hh}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}[(\gamma _1+\gamma _2)(k_x^2+k_y^2)+(\gamma _12\gamma _2)k_z^2`$ (72) $`_{lh}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}[(\gamma _1\gamma _2)(k_x^2+k_y^2)+(\gamma _1+2\gamma _2)k_z^2`$ (73) $`_{so}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}\gamma _1(k_x^2+k_y^2+k_z^2)+\mathrm{\Delta }_{so}`$ (74) $`b`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}\mathrm{}^2}{m}}\gamma _3k_z(k_xik_y)`$ (75) $`c`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}\mathrm{}^2}{2m}}\left[\gamma _2(k_x^2k_y^2)2i\gamma _3k_xk_y\right]`$ (76) $`d`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\mathrm{}^2}{2m}}\gamma _2\left[2k_z^2(k_x^2+k_y^2)\right].`$ (77) The four band Kohn-Luttinger Hamiltonian can be diagonalized analytically and yields a pair of Kramers doublets with eigenenergies $`\epsilon _𝐤={\displaystyle \frac{H_{hh}+H_{lh}}{2}}\sqrt{{\displaystyle \frac{1}{4}}(H_{hh}H_{lh})^2+\|b\|^2+\|c\|^2}.`$ (78) In the spherical approximation ($`\gamma _2,\gamma _3\overline{\gamma }0.6\gamma _2+0.4\gamma _3`$), the top of the valence band consists of two doubly degenerate parabolic bands with effective masses $`m_h=m/(\gamma _12\overline{\gamma })0.498m`$ and $`m_l=m/(\gamma _1+2\overline{\gamma })0.086m`$. (An approximation in which the light hole bands, which have a much smaller density of states, are ignored is adequate for some purposes.) The four band and six band representations for the hole spin-operator components read $`s_x`$ $`=`$ $`\left(\begin{array}{cccccc}0& 0& \frac{1}{2\sqrt{3}}& 0& \frac{1}{\sqrt{6}}& 0\\ 0& 0& \frac{1}{3}& \frac{1}{2\sqrt{3}}& \frac{1}{3\sqrt{2}}& 0\\ \frac{1}{2\sqrt{3}}& \frac{1}{3}& 0& 0& 0& \frac{1}{3\sqrt{2}}\\ 0& \frac{1}{2\sqrt{3}}& 0& 0& 0& \frac{1}{\sqrt{6}}\\ & & \multicolumn{-1}{c}{}& & & \\ \frac{1}{\sqrt{6}}& \frac{1}{3\sqrt{2}}& 0& 0& 0& \frac{1}{6}\\ 0& 0& \frac{1}{3\sqrt{2}}& \frac{1}{\sqrt{6}}& \frac{1}{6}& 0\end{array}\right)`$ (85) $`s_y`$ $`=`$ $`i\left(\begin{array}{cccccc}0& 0& \frac{1}{2\sqrt{3}}& 0& \frac{1}{\sqrt{6}}& 0\\ 0& 0& \frac{1}{3}& \frac{1}{2\sqrt{3}}& \frac{1}{3\sqrt{2}}& 0\\ \frac{1}{2\sqrt{3}}& \frac{1}{3}& 0& 0& 0& \frac{1}{3\sqrt{2}}\\ 0& \frac{1}{2\sqrt{3}}& 0& 0& 0& \frac{1}{\sqrt{6}}\\ & & \multicolumn{-1}{c}{}& & & \\ \frac{1}{\sqrt{6}}& \frac{1}{3\sqrt{2}}& 0& 0& 0& \frac{1}{6}\\ 0& 0& \frac{1}{3\sqrt{2}}& \frac{1}{\sqrt{6}}& \frac{1}{6}& 0\end{array}\right)`$ (94) $`s_z`$ $`=`$ $`\left(\begin{array}{cccccc}\frac{1}{2}& 0& 0& 0& 0& 0\\ 0& \frac{1}{6}& 0& 0& 0& \frac{\sqrt{2}}{3}\\ 0& 0& \frac{1}{6}& 0& \frac{\sqrt{2}}{3}& 0\\ 0& 0& 0& \frac{1}{2}& 0& 0\\ & & \multicolumn{-1}{c}{}& & & \\ 0& 0& \frac{\sqrt{2}}{3}& 0& \frac{1}{6}& 0\\ 0& \frac{\sqrt{2}}{3}& 0& 0& 0& \frac{1}{6}\end{array}\right)`$ (103)
warning/0006/astro-ph0006182.html	ar5iv	text	# Hot HB stars in globular clusters - Physical parameters and consequences for theory. V. Radiative levitation versus helium mixing ## 1 Introduction The colour-magnitude diagrams of metal-poor globular clusters show a large variety of horizontal-branch (HB) morphologies, including “gaps” along the blue HB and long “blue tails” that extend towards higher effective temperatures. It has been suggested that the gaps might separate stars with different evolutionary origins. Spectroscopic analyses of stars along the blue HB and blue tails in a number of globular clusters (Moehler moeh99 (1999) and references therein) yielded the following results: 1. Most of the stars analysed above and below any gaps are horizontal branch B type (HBB) stars ($`T_{\mathrm{eff}}<20,000`$ K). Their surface gravities are significantly lower (up to more than 0.5 dex, see Fig. 4 of Moehler moeh99 (1999)) than expected from canonical HB evolution theory while their masses are lower than expected by about a factor of 2. For most clusters the problem of the masses may be solved by new globular cluster distances derived from Hipparcos data (see Reid reid99 (1999) and references therein, Heber et al. hemo97 (1997), Moehler moeh99 (1999)). 2. Only in NGC 6752 and M 15 have spectroscopic analyses verified the presence of stars that could be identified with the subdwarf B stars known in the field of the Milky Way ($`T_{\mathrm{eff}}>20,000`$ K, $`\mathrm{log}g`$ $`>`$ 5). In contrast to the cooler HBB stars these stars show gravities and masses that agree well with the expectations of canonical stellar evolution for extreme HB stars (Moehler et al. 1997a , 1997b ). There are currently two scenarios for explaining these apparent contradictions: The dredge-up of nuclearly processed material to the stellar surface of red giant branch (RGB) stars has been invoked to explain the abundance anomalies (in C, N, O, Na, and Al) observed in such stars in many globular clusters (e.g. Kraft kraf94 (1994), Norris & Da Costa 1995a , Kraft et al. krsn97 (1997)). Since substantial production of Al in these low-mass stars only seems to occur inside the hydrogen shell (Langer & Hoffman laho95 (1995), Cavallo et al. casw96 (1996), casw98 (1998)), any mixing process which dredges up Al will also dredge up helium. Possible dredge-up mechanisms include rotationally induced mixing (Sweigart & Mengel swme79 (1979), Zahn zahn92 (1992), Charbonnel char95 (1995)) and hydrogen shell instabilities (Von Rudloff et al. voru88 (1988), Fujimoto et al. fuai99 (1999)). Such dredge-up would increase the helium abundance in the red giant’s hydrogen envelope and thereby increase the luminosity (and the mass loss) along the RGB (Sweigart 1997a , 1997b ). The progeny of these stars on the horizontal branch would then have less massive hydrogen envelopes than unmixed stars. As the temperature of an HB star increases with decreasing mass of the hydrogen envelope, “mixed” HB stars would be hotter than their canonical counterparts. The helium enrichment would also lead to an increased hydrogen burning rate and thus to higher luminosities (compared to canonical HB stars of the same temperature). The luminosities of stars hotter than about 20,000 K are not affected by this mixing process because these stars have only inert hydrogen shells. In this framework the low gravities of hot HB stars would necessarily be connected to abundance anomalies observed on the RGB, thereby explaining both of these puzzles at once. Caloi (calo99 (1999)) and Grundahl et al. (grca99 (1999)) suggested that the low surface gravities of the HBB stars are related to a stellar atmospheres effect caused by the radiative levitation of heavy elements. Such an enrichment in the metal abundance would change the temperature structure of the stellar atmosphere and thereby affect the flux distribution and the line profiles (Leone & Manfrè lema97 (1997)). This scenario would also account for the fact that there is no evidence for deep mixing amongst field red giants (e.g. Hanson et al. hasn98 (1998), Carretta et al. cagr99 (1999)) even though field HBB stars show the same low surface gravities as globular cluster stars (Saffer et al. sake97 (1997), Mitchell et al. misa98 (1998)). Behr et al. (beco99 (1999), 2000b ) have recently reported slightly super-solar iron abundances for HBB stars in M 13 and M 15, in agreement with the radiative levitation scenario. NGC 6752 is an ideal test case for these scenarios: Its distance modulus is very well determined from both white dwarfs (Renzini et al. rebr96 (1996)) and HIPPARCOS parallaxes (Reid reid97 (1997)), and thus any mass discrepancies cannot be explained by a wrong distance modulus. Spectroscopic analyses of the faint blue stars in NGC 6752 showed them to be subdwarf B (sdB) stars. As mentioned above, their mean mass agrees well with the canonical value of 0.5 M. However, almost no stars in the sparsely populated region above the sdB star region have been analysed. If these stars show low surface gravities and canonical masses, then the combination of deep mixing and the long distance scale (for the other globular clusters) would resolve the discrepancies described above. If they show low surface gravities and low masses, diffusion may indeed play a rôle when analysing these stars for effective temperature and surface gravity. Then the low surface gravities found for HBB stars could be artifacts from the use of inappropriate model atmospheres for the analyses. We therefore decided to observe stars in this region of the colour-magnitude diagram and to derive their atmospheric parameters. First results, which strongly support radiative levitation of heavy elements as the explanation, have been discussed by Moehler et al. (1999a ). Here we describe the observations and their reductions, provide the detailed results of the spectroscopic analyses (temperatures, surfaces gravities, helium and partly iron and magnesium abundances, masses), and discuss the consequences of our findings in more detail. ## 2 Observations We selected our targets from the photographic photometry of Buonanno et al. (buca86 (1986), see Table 1 and Fig. 1). For our observations we used the ESO 1.52m telescope with the Boller & Chivens spectrograph and CCD #39 (2048$`\times `$2048 pixels, (15 $`\mu `$m)<sup>2</sup> pixel size, read-out noise 5.4 e<sup>-</sup>, conversion factor 1.2 e<sup>-</sup>/count). We used grating # 33 (65 Åmm<sup>-1</sup>) to cover a wavelength of 3300 Å – 5300 Å. Combined with a slit width of 2″ we thus achieved a spectral resolution of 2.6 Å. The spectra were obtained on July 22-25, 1998. For calibration purposes we observed each night ten bias frames and ten dome flat-fields with a mean exposure level of about 10,000 counts each. Before and after each science observation we took HeAr spectra for wavelength calibration purposes. We observed dark frames of 3600 and 1800 sec duration to measure the dark current of the CCD. As flux standard stars we used LTT 7987 and EG 274. We also analyse data that were obtained as backup targets at the NTT during observing runs dedicated to other programs (60.E-0145, 61.E-0361). The observational set-up and the data reduction are described in Moehler et al. (mola99 (2000), 1999b ). These data have a much lower resolution of 5.4 Å. ## 3 Data Reduction We first averaged the bias and flat field frames separately for each night. As we could not detect any significant change in the mean bias level we computed the median of the bias frames of the four nights and found that the bias level showed a gradient across the image, increasing from the lower left corner to the upper right corner by about 1%. We fitted the bias with a linear approximation along both axes and used this fit as a bias for the further reduction. As no overscan was recorded we could not adjust the bias level. Bias frames taken during the night, however, revealed no significant change in the mean bias level. The mean dark current determined from long dark frames showed no structure and turned out to be negligible (3$`\pm `$3 e<sup>-</sup>/hr/pixel). We determined the spectral energy distribution of the flat field lamp by averaging the mean flat fields of each night along the spatial axis. These one-dimensional “flat field spectra” were then heavily smoothed and used afterwards to normalize the dome flats along the dispersion axis. The normalized flat fields of the first three nights were combined. For the fourth night we used only the flat field obtained during that night as we detected a slight variation in the fringe patterns of the flat fields from the first three nights compared to that of the fourth (below 5%). For the wavelength calibration we fitted 3<sup>rd</sup>-order polynomials to the dispersion relations of the HeAr spectra which resulted in mean residuals of $``$0.1 Å. We rebinned the frames two-dimensionally to constant wavelength steps. Before the sky fit the frames were smoothed along the spatial axis to erase cosmic ray hits in the background. To determine the sky background we had to find regions without any stellar spectra, which were sometimes not close to the place of the object’s spectrum. Nevertheless the flat field correction and wavelength calibration turned out to be good enough that a linear fit to the spatial distribution of the sky light allowed the sky background at the object’s position to be reproduced with sufficient accuracy. This means in our case that after the fitted sky background was subtracted from the unsmoothed frame we do not see any absorption lines caused by the predominantly red stars of the clusters. The sky-subtracted spectra were extracted using Horne’s (horn86 (1986)) algorithm as implemented in MIDAS (Munich Image Data Analysis System). Finally the spectra were corrected for atmospheric extinction using the extinction coefficients for La Silla (Tüg tueg77 (1977)) as implemented in MIDAS. The data for the flux standard stars were taken from Hamuy et al. (hamu92 (1992)) and the response curves were fitted by splines. The flux-calibration is helpful for the later normalization of the spectra as it takes out all large-scale sensitivity variations of the instrumental setup. Absolute photometric accuracy is not an issue here. ## 4 Atmospheric Parameters To derive effective temperatures, surface gravities and helium abundances we fitted the observed Balmer and helium lines with stellar model atmospheres. Beforehand we corrected the spectra for radial velocity shifts, derived from the positions of the Balmer and helium lines. The resulting heliocentric velocities are listed in Table 1. The error of the velocities (as estimated from the scatter of the velocities derived from individual lines) is about 40 km s<sup>-1</sup>. The spectra were then normalized by eye and are plotted in Figs. 2 and 3. To establish the best fit we used the routines developed by Bergeron et al. (besa92 (1992)) and Saffer et al. (sabe94 (1994)), which employ a $`\chi ^2`$ test. The $`\sigma `$ necessary for the calculation of $`\chi ^2`$ is estimated from the noise in the continuum regions of the spectra. The fit program normalizes model spectra and observed spectra using the same points for the continuum definition. We computed model atmospheres using ATLAS9 (Kurucz kuru91 (1991)) and used Lemke’s version<sup>1</sup><sup>1</sup>1For a description see http://a400.sternwarte.uni-erlangen.de/$``$ai26/linfit/linfor.html of the LINFOR program (developed originally by Holweger, Steffen, and Steenbock at Kiel University) to compute a grid of theoretical spectra which include the Balmer lines H<sub>α</sub> to H<sub>22</sub> and He i lines. The grid covered the range 7,000 K $``$ $`T_{\mathrm{eff}}`$ $``$ 35,000 K, 2.5 $``$ $`\mathrm{log}g`$ $``$ 6.0, $`3.0`$ $``$ $`\mathrm{log}\frac{n_{\mathrm{He}}}{n_\mathrm{H}}`$ $``$ $`1.0`$, at a metallicity of \[M/H\] = $`1.5`$. In Table 2 we list the results obtained from fitting the Balmer lines H<sub>β</sub> to H<sub>10</sub> (excluding H<sub>ϵ</sub> to avoid the Ca ii H line) and the He i lines 4026 Å, 4388 Å, 4471 Å, and 4921 Å. The errors given are r.m.s. errors derived from the $`\chi ^2`$ fit (see Moehler et al. 1999b for more details). These errors are obtained under the assumption that the only error source is statistical noise (derived from the continuum of the spectrum). However, errors in the normalization of the spectrum, imperfections of flat field/sky background correction, variations in the resolution (e.g. due to seeing variations when using a rather large slit width) and other effects may produce systematic rather than statistic errors, which are not well represented by the error obtained from the fit routine. Systematic errors can only be quantified by comparing truly independent analyses of the same stars. As this is not possible here we use our experience with the analysis of similar stars and estimate the true errors to be about 10% in $`T_{\mathrm{eff}}`$ and 0.15 dex in $`\mathrm{log}g`$ (cf. Moehler et al. 1997b , mola98 (1998)). Two stars show $`BV`$ colours that are significantly redder than expected from their effective temperatures (B 2697: $`BV`$ = $`+0\stackrel{\mathrm{m}}{.}08`$, $`T_{\mathrm{eff}}`$ = 15,700 K; B 3006: $`BV`$ = $`0\stackrel{\mathrm{m}}{.}10`$, $`T_{\mathrm{eff}}`$ = 30,000 K), possibly indicating that the colours are affected by binarity or photometric blending with a cool star. While the spectra look quite normal, we will not include these stars in any statistical discussion below. To increase our data sample we reanalysed the NTT spectra described and analysed by Moehler et al. (1997b ). We did not reanalyse the EFOSC1 data published in the same paper as they are of worse quality. We find that the atmospheric parameters determined by line profile fitting agree rather well with those published by Moehler et al. (1997b ). The temperatures and gravities obtained from these metal-poor atmospheres are compared with the values predicted by canonical HB tracks in Fig. 4 (top panel). These tracks, which were computed for a main sequence mass of 0.805 M, an initial helium abundance Y of 0.23 and a scaled-solar metallicity \[M/H\] of $``$1.54, define the locus of canonical HB models which lose varying amounts of mass during the RGB phase. According to the Reimers mass-loss formulation the value of the mass-loss parameter $`\eta _R`$ would vary from $``$0.4 at the red end of the observed HB in NGC 6752 to $``$0.7 for the sdB stars, given the present composition parameters. One can see from Fig. 4 (top panel) that the HBB stars in NGC 6752 show the same effect as seen in other globular clusters, namely, an offset from the zero-age horizontal branch (ZAHB) towards lower surface gravities over the temperature range 4.05 $`<`$ $`\mathrm{log}T_{\mathrm{eff}}`$ $`<`$ 4.30 (11,200 K $`<`$ $`T_{\mathrm{eff}}`$ $`<`$ 20,000 K). At lower or higher temperatures the gravities agree with the locus of the canonical HB tracks. ### 4.1 Radiative levitation of heavy elements As described in Moehler et al. (1999a , see also Fig. 5), we found evidence for iron enrichment in the spectra of the HBB stars obtained at the ESO 1.52m telescope, whereas the magnesium abundance appeared consistent with the cluster magnesium abundance. The actual iron abundances derived for these stars by fitting the iron lines in the ESO 1.52m spectra are listed in Table 4. The mean iron abundance turns out to be \[Fe/H\] = $`+0.12\pm 0.40`$ (internal errors only, $`\mathrm{log}ϵ_{Fe}=7.58`$) for stars hotter than about 11,500 K – in good agreement with the findings of Behr et al. (beco99 (1999), 2000b ) for HBB/HBA (horizontal branch A type) stars in M 13 and M 15 and Glaspey et al. (glmi89 (1989)) for two HBB/HBA stars in NGC 6752. This iron abundance is a factor of 50 greater than that of the cluster, but still a factor of 3 smaller than that required to explain the Strömgren $`u`$-jump discussed by Grundahl et al. (grca99 (1999), $`\mathrm{log}ϵ_{Fe}=8.1`$). The mean magnesium abundance for the same stars is \[Mg/H\] = $`1.13\pm 0.29`$ (internal errors only), corresponding to \[Mg/Fe\] = $`+`$0.4 for \[Fe/H\] = $``$1.54. This value agrees well with the abundance \[Mg/Fe\] = $`+0.4`$ found by Norris & da Costa (1995b ) for red giants in NGC 6752. The abundances are plotted versus temperature in Fig. 5. The trend of decreasing helium abundance with increasing temperature seen in the ESO 1.52m data (and also reported by Behr et al. beco99 (1999) for HB stars in M 13) is not supported towards higher temperatures by the NTT data. This could be due to the lower resolution of the NTT data which may tend to overestimate abundances (Glaspey et al. glmi89 (1989)). As iron is very important for the temperature stratification of stellar atmospheres we tried to take the increased iron abundance into account by computing model atmospheres for \[M/H\] = 0. Indeed a backwarming effect of 2–4% on the temperature structure was found in the formation region of the Balmer lines, when comparing solar composition models with the metal-poor models. We then repeated the fit to derive $`T_{\mathrm{eff}}`$, $`\mathrm{log}g`$, and $`\mathrm{log}\frac{n_{\mathrm{He}}}{n_\mathrm{H}}`$ with these enriched model atmospheres. The resulting effective temperatures and gravities changed as displayed in Fig. 6. The results are listed in Table 3 and plotted in Fig. 4 (central panel). From Fig. 4 (central panel) it is clear that the use of solar-metallicity model atmospheres moves most stars closer to the canonical zero-age horizontal branch (ZAHB) due to a combination of lower $`T_{\mathrm{eff}}`$ and/or higher $`\mathrm{log}g`$. The three stars between 10,000 K and 12,000 K, however, fall below the canonical ZAHB when fitted with enriched model atmospheres. This is plausible as the radiative levitation is supposed to start around 11,000 – 12,000 K (Grundahl et al. grca99 (1999)) and the cooler stars therefore should have metal-poor atmospheres (see also Fig. 5, where the coolest analysed star shows no evidence of iron enrichment). This assumption is also supported by the results of Glaspey et al. (glde85 (1985), NGC 6397; glmi89 (1989), NGC 6752) and Behr et al. (beco99 (1999), M 13; 2000b , M 15). Now the stars below 15,300 K scatter around the locus defined by the canonical HB tracks. The stars between 15,500 K and 19,000 K, however, still show offsets from the canonical locus while for the sdB stars not much is changed. Interestingly, 15,500 K is roughly the temperature<sup>2</sup><sup>2</sup>2We determined this temperature by comparing the $`(uy)_0`$ value, at which the stars return to the ZAHB ($`(uy)_0+0.4`$) to theoretical colours from Kurucz (kuru92 (1992)) for \[M/H\] = $`+0.5`$, which is the metallicity required to explain the $`u`$-jump. Assuming $`\mathrm{log}g`$ = 4.0 this comparison results in $`T_{\mathrm{eff}}`$ $``$ 15,000 K. at which the stars in NGC 6752 return to the ZAHB in ($`uy`$, $`u`$) of Grundahl et al. (grca99 (1999)). Grundahl et al. caution, however, that their faint photometry for NGC 6752 might be affected by poor seeing, and that in the Strömgren CMD of the better observed cluster, M13, the stars do not return to the ZAHB until a temperature of about 20,000 K. We next repeated the Balmer line profile fits by increasing the metal abundance of the model atmospheres to \[M/H\]=$`+`$0.5 (see Fig. 4, bottom panel, and Table 5), which did not significantly change the resulting values for $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$. In particular, note that especially the “deviant” stars (now between 15,300 K and 19,000 K) remain offset from the canonical ZAHB. ### 4.2 Helium mixing As outlined in Sect. 1, helium mixing during the RGB phase may also be able to explain the low gravities of the HBB stars. Under this scenario the mixing currents within the radiative zone below the base of the convective envelope of a red giant star are assumed to penetrate into the top of the hydrogen shell where helium is being produced by the hydrogen burning reactions. Ordinarily one would expect the gradient in the mean molecular weight $`\mu `$ to prevent any penetration of the mixing currents into the shell. If, however, the timescale for mixing were shorter than the timescale for nuclear burning, then the helium being produced at the top of the shell might be mixed outward into the envelope before a $`\mu `$ gradient is established. Under these circumstances a $`\mu `$ gradient would not inhibit deep mixing simply because such a gradient would not exist within the mixed region. Since deep mixing is presumably driven by rotation, one would expect a more rapidly rotating red giant to show a larger increase in the envelope helium abundance. This, in turn, would lead to a brighter RGB tip luminosity and hence to greater mass loss. The progeny of the more rapidly rotating giants should therefore lie at higher effective temperatures along the HB than the progeny of the more slowly rotating giants. This predicted increase in the stellar rotational velocity with effective temperature along the HB has not, however, been confirmed by the recent observations of M13 by Behr et al. (2000a ). These observations show that HB stars in M13 hotter than 11,000 K are, in fact, rotating slowly with $`v\mathrm{sin}i<`$ 10 km s<sup>-1</sup> in contrast to the cooler HB stars where rotational velocities as high as 40 km s<sup>-1</sup> are found (see also Peterson et al. pero95 (1995)). There are a couple of possible explanations for this apparent discrepancy. One possibility is that the greater mass loss suffered by the HBB stars might carry away so much angular momentum that the surface layers are spun down even though the core is still rotating rapidly. Alternatively Sills & Pinsonneault (sipi00 (2000)) have suggested that the observed gravitational settling of helium in HBB stars might set up a $`\mu `$ gradient in the outer layers which inhibits the transfer of angular momentum from the rapidly rotating interior to the surface. Thus the surface rotational velocities may not necessarily be indicative of the interior rotation. In order to explore the consequences of helium mixing for the HBB stars quantitatively, we evolved a set of 13 sequences up the RGB to the helium flash for varying amounts of helium mixing using the approach of Sweigart (1997a , 1997b ). As in the case of the canonical models discussed previously, all of these mixed sequences had an initial helium abundance Y of 0.23 and a scaled-solar metallicity \[M/H\] of $``$1.54. The main-sequence mass was taken to be 0.805 M, corresponding to an age at the tip of the RGB of 15 Gyr. The mixing depth, as defined by the parameter $`\mathrm{\Delta }X_{mix}`$ of Sweigart (1997a , 1997b ), ranged from 0.0 (canonical, unmixed case) to 0.24 in increments of 0.02. Mass loss via the Reimers formulation was included in the calculations with the mass-loss parameter $`\eta _R`$ set equal to 0.40. This value for $`\eta _R`$ was chosen so that a canonical, unmixed model would lie near the red end of the observed blue HB in NGC 6752. Both the mixing and mass loss were turned off once the models reached the core He flash at the tip of the RGB, and the subsequent evolution was then followed through the helium flash to the end of the HB phase using standard techniques. We did not investigate the changes in the surface abundances of CNO, Na and Al caused by the helium mixing, since such a study was beyond the scope of the present paper. Rather, our objective was to determine how the mixing affected those quantities which impact on the HB evolution, i.e., envelope helium abundance and mass. We do note that the mixing in the more deeply mixed RGB models would have penetrated into regions of substantial Na and Al production according to the calculations of Cavallo et al. (casw96 (1996), casw98 (1998)). However, the resulting changes in the surface Na and Al abundances will depend on the assumed initial Ne and Mg isotopic abundances and on the adopted nuclear reaction rates, which in some cases are quite uncertain. The locus of the above helium-mixed sequences in the $`\mathrm{log}g`$ \- $`\mathrm{log}T_{\mathrm{eff}}`$ plane is indicated by the dashed lines in the top panel of Fig. 4. The red end of the mixed ZAHB in this panel, located at $`\mathrm{log}T_{\mathrm{eff}}`$ = 3.93, is set by the canonical, unmixed sequence for the present set of model parameters. Since mixing increases the RGB mass loss, a mixed HB model will have a higher effective temperature than the corresponding canonical model. At the same time mixing increases the envelope helium abundance in the HB model, which, in turn, increases both the hydrogen-burning and surface luminosities. The net effect is to shift the mixed locus in Fig. 4 towards lower gravities with increasing $`T_{\mathrm{eff}}`$ compared to the canonical locus, until a maximum offset is reached for 15,500 K $`<`$ $`T_{\mathrm{eff}}`$ $`<`$ 19,000 K. At higher temperatures the mixed locus shifts back towards the canonical locus, as the contribution of the hydrogen shell to the surface luminosity declines due to the decreasing envelope mass. The predicted locus along the extreme HB (EHB) does not depend strongly on the extent of the mixing, since the luminosities and gravities of the EHB stars are primarily determined by the mass of the helium core, which is nearly the same for the mixed and canonical models. Overall the variation of $`\mathrm{log}g`$ with $`T_{\mathrm{eff}}`$ along the mixed locus in the top panel of Fig. 4 mimics the observed variation. The results presented in Sect. 4.1 demonstrate that radiative levitation of heavy elements can account for a considerable fraction of the gravity offset along the HBB, especially for temperatures cooler than 15,100 K. Consequently the amount of helium mixing required to explain the remaining offset between 15,300 K and 19,000 K is much less than the amount required to explain the offsets found without accounting for radiative levitation (top panel of Fig. 4). In order to compare the gravities predicted by the helium-mixing scenario with those derived from the metal-enhanced atmospheres, we computed a second set of mixed sequences using the same approach as above but with a larger value of the mass-loss parameter $`\eta _R`$, i.e., $`\eta _R`$ = 0.45. The red end of the mixed ZAHB for these sequences is located at $`\mathrm{log}T_{\mathrm{eff}}`$ = 4.01 and is therefore hotter than the red end of the mixed ZAHB for the sequences with $`\eta _R`$ = 0.40. The HB stars cooler than this temperature in NGC 6752 would then be identified with unmixed stars which lost less mass along the RGB. By increasing the mass loss efficiency we reduce the amount of mixing needed to populate the temperature range 15,300 K $`<`$ $`T_{\mathrm{eff}}`$ $`<`$ 19,000 K and therefore the size of the resulting gravity offset. The locus of the mixed sequences with $`\eta _R`$ = 0.45 is indicated by the dashed lines in the central panel of Fig. 4. The gravity offsets along this mixed locus seem to provide a reasonable fit to the gravities given by the model atmospheres with solar metallicity. Finally we computed a third set of mixed sequences with the mass-loss parameter increased further to $`\eta _R`$ = 0.50 for comparison with the gravities obtained from the atmospheres with super-solar metallicity in the bottom panel of Fig. 4. As expected, these mixed sequences show a smaller gravity offset in the temperature range 15,500 K $`<`$ $`T_{\mathrm{eff}}`$ $`<`$ 19,000 K. Moreover, the red end of the mixed ZAHB shifts blueward to $`\mathrm{log}T_{\mathrm{eff}}`$ = 4.08. ## 5 Masses We calculated masses for the programme stars in NGC 6752 from their values of $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ using the equation: $`\mathrm{log}\frac{\mathrm{M}}{\mathrm{M}_{}}=const.+\mathrm{log}g+0.4((mM)_VV+V_{th})`$ where $`V_{th}`$ denotes the theoretical brightness at the stellar surface as given by Kurucz (kuru92 (1992)). We decided to use the photometry of Buonanno et al. (buca86 (1986)) to derive masses. As can be seen from Table 1 the photometry of Thompson et al. (thka99 (1999)) yields in general fainter visual magnitudes than the photometry of Buonanno et al. The effect on the masses, however, is small: On average, the masses derived from the Thompson et al. photometry are 5% lower than those derived from the Buonanno et al. photometry. We adopted $`(mM)_0=13.17`$ and $`E_{BV}`$ = 0.04 for the distance modulus and reddening. These are mean values derived from the determinations of Renzini et al. (rebr96 (1996)), Reid (reid97 (1997), reid98 (1998)), and Gratton et al. (grfu97 (1997)). The errors in $`\mathrm{log}\mathrm{M}`$ are estimated to be about the same as obtained for the older NGC 6752 data described by Moehler et al. (1997b ): 0.15 dex for stars above the gap, 0.17 dex for stars within the gap region and 0.22 dex for stars below the gap. The masses derived from the analysis using metal-poor model atmospheres are plotted in Fig. 7 (top panel). The sdB stars hotter than 20,000 K ($`\mathrm{log}T_{\mathrm{eff}}`$ = 4.3) scatter around the canonical ZAHB, whereas the stars below 16,000 K ($`\mathrm{log}T_{\mathrm{eff}}`$ = 4.2) lie mainly below the canonical ZAHB. Even stronger deviations towards low masses are found between 16,000 K and 20,000 K. Comparing the masses to those predicted by the mixed ZAHB ($`\eta _\mathrm{R}=`$ 0.40) we obtain similar results. To quantify the offsets we compare the masses of the stars to those they would have on the theoretical ZAHB at the same $`T_{\mathrm{eff}}`$ (see Table 6). We divide the stars into three groups for the further discussion (excluding stars below 11,500 K, for which diffusion should play no rôle, as well as B 2697 and B 3006): cool HBB stars ($`T_{\mathrm{eff}}`$ $`<`$ 16,000 K, 16 stars), hot HBB stars (16,000 K $``$ $`T_{\mathrm{eff}}`$ $`<`$ 20,000 K, 9 stars), and sdB stars ($`T_{\mathrm{eff}}`$ $``$ 20,000 K, 12 stars). The effective temperatures here are those derived from metal-poor model atmospheres. The results of the analyses using solar-metallicity model atmospheres are plotted in the central panel of Fig. 7 (see Table 6). The effect on the masses is similar to that on the temperatures/gravities (Fig. 4, central panel) – below 15,300 K (cool HBB stars) and above 20,000 K (sdB stars) the masses basically scatter around the canonical ZAHB, but the hot HBB stars between these two groups still show too low masses. Comparing the masses to those predicted by the mixed ZAHB ($`\eta _\mathrm{R}`$ = 0.45) gives similar results. As the stars become cooler when analysed with more metal-rich atmospheres the temperature boundaries were shifted to include the same stars as for the comparison made above. Even the use of metal-rich model atmospheres for the analyses does not change much (see Table 6 and the bottom panel of Fig. 7). Obviously the hot HBB stars in the intermediate temperature range still show low masses, despite the use of metal-rich model atmospheres. Thus the problem of the stars in this temperature range(15,300 K$`<`$ $`T_{\mathrm{eff}}`$ $`<`$19,000 K) cannot be completely solved by the scaled-solar metal-rich atmospheres used here. ## 6 Discussion We find that the atmospheres of HBB stars in NGC 6752 with $`T_{\mathrm{eff}}`$ $`>`$ 11,500 K are enriched in iron (\[Fe/H\] $`+`$0.1) whereas their magnesium abundances are the same as found in cluster giants. Our results are consistent with those of Behr et al. (beco99 (1999), 2000b ) for HBB stars in M 13 and M 15. Using model atmospheres that try to take into account this enrichment in iron (and presumably other heavy elements) reconciles the atmospheric parameters of stars with 11,500 K $``$ $`T_{\mathrm{eff}}`$ $`<`$ 15,100 K with canonical expectations as suggested by Grundahl et al. (grca99 (1999)). Also the masses derived from these analyses are in good agreement with canonical predictions within this temperature range. However, we found that even with model atmospheres as metal rich as \[M/H\] = $`+`$0.5 the atmospheric parameters of the hot HBB stars (15,300 K $`<`$ $`T_{\mathrm{eff}}`$ $`<`$ 19,000 K) in NGC 6752 cannot be reconciled with the canonical ZAHB. Both the gravities and masses of these hot HBB stars remain too low. In addition, the masses for the stars below 15,100 K are slightly too high for the super-solar metallicity (the EHB stars are hardly affected at all by changes in the metallicity of the model atmospheres). Michaud et al. (miva83 (1983)) noted that diffusion will not necessarily enhance all heavy elements by the same amount and that the effects of diffusion vary with effective temperature. Elements that were originally very rare may be enhanced even more strongly than iron (see also Behr et al. beco99 (1999), where P and Cr are enhanced to \[M/H\] $`+1`$). The question of whether diffusion is the (one and only) solution to the “low gravity” problem cannot be answered without detailed abundance analyses to determine the actual abundances and the use of model atmospheres that allow the use of non-scaled solar abundances (like ATLAS12). We can, however, state that those model atmospheres, which reproduce the $`u`$-jump discussed by Grundahl et al. (grca99 (1999)) cannot completely reconcile the atmospheric parameters of hot HB stars with canonical theory. Model atmospheres with abundance distributions that may solve the discrepancy between theoretically predicted and observed atmospheric parameters of hot HB stars may then, in turn, not reproduce the Strömgren $`u`$-jump. It is intriguing that the temperature, at which the stars in $`u,uy`$ seem to return to the ZAHB, is roughly the same at which they start to deviate again from the canonical ZAHB in $`\mathrm{log}g`$, $`T_{\mathrm{eff}}`$ when analysed with metal-rich atmospheres. The stars between 15,300 K and 19,000 K (when analysed with metal-rich atmospheres) are currently best fit by a moderately mixed ZAHB. However, the fact that their masses are too low cautions against identifying He mixing as the only cause for these low gravities - because in this case the luminosities of the stars would be increased and canonical masses would result. ###### Acknowledgements. We want to thank the staff of the ESO La Silla observatory for their support during our observations. We would also like to thank T. Lanz for very helpful discussions and the referees J. Cohen and B. Behr for useful remarks. S.M. acknowledges financial support from the DARA under grant 50 OR 96029-ZA. A.V.S. acknowledges financial support from NASA Astrophysics Theory Program proposal NRA-99-01-ATP-039.
warning/0006/nlin0006007.html	ar5iv	text	# Exploring phase space localization of chaotic eigenstates via parametric variation ## I Introduction The two general motivations for our investigation are understanding better the nature of eigenstates of bounded quantum systems possessing ‘simple’ classical analogs, and exploring new features of such systems’ behavior as a system parameter is smoothly varied. Simple in this context refers to few degrees of freedom and a compact Hamiltonian. Nevertheless, the classical dynamics may display a rich variety of features from regular to strongly chaotic motion. We focus on the strongly chaotic limit for which semiclassical quantization of individual chaotic eigenstates does not hold, and the correspondence principle is less well understood . Even though there has been some recent progress , it turns out that with a detailed understanding of chaotic systems a statistical theory provides a well-developed, alternative approach to these difficulties. Twenty years ago, Berry conjectured and Voros discussed that in this case as $`\mathrm{}0`$ the eigenstates should respect the ergodic hypothesis in phase space, $`\delta (EH(𝐩,𝐪))`$, as it applies to wavefunctions. In essence, the eigenmodes should appear as Gaussian random wavefunctions locally in configuration space with their wavevector constrained by the ergodic measure of the energy surface. Discussion of the properties of random waves and recent supporting numerical evidence can be found in refs. . The second general motivation relates to a long recognized class of problems, i.e. a system’s response to parametric variation. Our interest here is restricted to external, controllable parameters such as electro-magnetic fields, temperatures, applied stresses, changing boundary conditions, etc…, through whose variation one can extract new information about a system not available by other means. A multitude of examples can be found in the literature . A recent concern has been universalities in the response of chaotic or disordered systems and statistical approaches to measuring the response . Universal parametric correlations have been derived via field theoretic or random matrix methods for quantities involving level slopes (loosely termed velocities in this paper), level curvatures, and eigenfunction amplitudes . In contrast, our motivation is not the universal features per se for they cannot tell us anything specific about the system other than it is, in fact, chaotic and/or symmetry is present. Rather we are interested in what system specific information can be extracted in the case that the system’s response deviates from universal statistical laws. The specific application discussed in this paper shows how one can decipher phase space localization features of the eigenstates. The theory naturally divides into a two-step process. One must first understand any implied limiting universal response of chaotic systems. Next, one must develop a theory which gives a correct interpretation of any deviations seen from the universal response. The necessarily close interplay between theory and observation required to deduce new information forms part of the attractiveness of investigating parametric response. Taking up the first step of understanding universal response, an expected but rarely discussed property is the independence of eigenvalue and eigenfunction fluctuation measures which is found in the random matrix theories anticipated to describe the statistical properties of quantum systems with chaotic classical analogs . Coupled with Berry’s conjecture mentioned above, these properties imply a ‘democratic’ response to parametric variation for an ergodically behaving quantum system. The perturbation connects one state to all other states locally with equal probability. The variation of any one eigenstate or eigenvalue over a large enough parameter range will be statistically equivalent to their respective neighboring states or levels. In a pioneering work on the ergodic hypothesis using the stadium billiard, now a paradigm of chaos studies, McDonald noticed larger than average intensities of the eigenstates in certain regions . In his thesis he states that “a small class of modes (bouncing ball, whispering gallery, etc.) seem to correspond naively to a definite set of ‘special’ ray orbits.” Heller initiated a theory concerning these large intensities when he modified the random wavefunction picture with his prediction and numerical observations of eigenstate scarring . He derived a criterion for eigenstate intensity in excess of the ergodic predictions along the shorter, less unstable periodic orbits. Scarring is thus one possible phase space ‘localization’ property of a chaotic eigenstate. Other possibilities result from time scales not related directly to the Lyapunov instability such as transport barriers in the form of broken separatrices , and cantori , or diffusive motion . In the context of this paper, we take localization to mean some deviation from the ergodic expectation beyond the inherent quantum fluctuations, and it creates the possibility of a non-democratic response to parametric variation. A perturbation could preferentially connect certain states or classes of states, thus leading to additional short-range avoided crossings or like level movements within a particular class, etc. Debate ensued Heller’s work on eigenstate scarring, in part, because of the difficulty in quantitatively characterizing and predicting its extent in either a particular eigenstate or even collective groups of eigenstates. Judging from the earlier literature, it was easier to graph eigenstates in order to see the scarring by eye than define precisely what it means or what its physical significance is. Furthermore, he linearized the semiclassical theory which was insufficient for a full description of scarring. We remark that recent work suggests the opposite, i.e. the linearized theory is sufficient assuming $`\mathrm{}`$ is smaller than some system specific value which is ‘small enough’ . However, many of the experimental and numerical investigations are far from this regime and the nonlinear dynamical contributions are essential for understanding most of the work being done. The theory incorporating nonlinear dynamical contributions was developed much later than Heller’s introduction of scarring. It is based on heteroclinic orbit expansions for wave packet propagation and strength functions. Ahead, we make extensive use of these forms to derive a semiclassical theory applicable to problems involving parametric variation. In a previous Letter , one of us (ST) introduced a measure that very sensitively probes phase space localization for systems having continuously tunable parameters in their Hamiltonians. It correlates level motions under perturbation with overlap intensities between eigenstates and optimally localized wave packet states. The basic idea is that the wave packet overlap intensities select eigenstates that potentially have excess support in the neighborhood of the phase point at the wave packet’s position and momentum centroids. The perturbation will push these levels somewhat in the same direction depending on how it is distorting the energy surface near that particular phase point. If the level velocities associated with those states have similar enough values, then significant non-zero correlations will result that reveal the localization. The measure can be used in a forward or reverse direction. If phase space localization is present in a system of interest, then it predicts experimentally verifiable manifestations of that localization. Conversely, one can first experimentally determine the level velocity - overlap intensity measure in that system for the purpose of inferring the existence and extent of localization. Our purpose in this paper is to give a complete account of that Letter, develop further the semiclassical theory, and explore the full phase space and $`\mathrm{}`$ behavior of the stadium billiard, a continuous time system. In a companion paper immediately following this one, we give the theory for quantized maps (discretized time) . The next section introduces strength functions and a new class of correlation coefficients. Section III utilizes ergodicity and random wave properties to motivate the introduction of random matrix ensembles. The ensembles describe the statistical properties of chaotic systems in the $`\mathrm{}0`$ limit. The correlation measures vanish for these ensembles indicating the absence of localization and universal response to perturbation (i.e. parameter variation). Section IV gives the semiclassical theories of level velocities, strength functions, and overlap intensity-level velocity correlation coefficients. We finish with a full treatment of the stadium billiard and concluding remarks. ## II Preliminaries Consider a quantum system governed by a smoothly parameter-dependent Hamiltonian, $`\widehat{H}(\lambda )`$ with classical analog $`H(𝐩,𝐪;\lambda )`$. We suppose that the dynamics are chaotic for all values of the $`\lambda `$ range of interest, and suppose the absence of symmetry breaking. Then the expectation is that all statistical properties are stationary with respect to $`\lambda `$. Without loss of generality, we also assume the phase space volume of the energy surface is constant as a function of $`\lambda `$. This ensures that the eigenvalues do not collectively drift in some direction in energy, but rather wander locally. We use the same strength function Heller employed in his prediction of scarring except slightly generalized to include parametric behavior; $`S_\alpha (E,\lambda )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dte^{iEt/\mathrm{}}\alpha \|e^{i\widehat{H}(\lambda )t/\mathrm{}}\|\alpha `$ (1) $`=`$ $`Tr[\widehat{p}_\alpha \delta (E\widehat{H}(\lambda ))]`$ (2) $`=`$ $`{\displaystyle \underset{n}{}}p_{\alpha n}(\lambda )\delta (EE_n(\lambda ));`$ (3) $`p_{\alpha n}(\lambda )`$ $`=`$ $`\|\alpha \|E_n(\lambda )\|^2`$ (4) where $`\widehat{p}_\alpha =\|\alpha \alpha \|`$. $`S_\alpha (E,\lambda )`$ is the Fourier transform of the autocorrelation function of a special initial state $`\|\alpha `$ of interest. Ahead $`\overline{S}_\alpha (E,\lambda )`$ will denote the smooth part resulting from the Fourier transform of just the extremely rapid initial decay due to the shortest time scale of the dynamics (zero-length trajectories). We will take $`\|\alpha `$ to be a Gaussian wave packet because of its ability to probe “quantum phase space,” but other choices are possible. Say momentum space localization were the main interest, the natural choice would be a momentum eigenstate. $`\|\alpha `$ can be associated with a phase space image $`\rho _\alpha (𝐩,𝐪)`$ of Gaussian functional form using Wigner transforms or related techniques. $`\rho _\alpha (𝐩,𝐪)`$ turns out to be positive definite and maximally localized in phase space, i.e. it occupies a volume of $`h^d`$. For a fixed value of the parameter, an example strength function is shown in Fig. (1). If the wave packet is centered somewhere on a short periodic orbit, large amplitudes necessarily indicate significant wave intensity all along the orbit as seen in the inset eigenstates. This behavior cannot, a priori, be stated to be obviously in violation of the quantum statistical fluctuation laws even if it appears so. That remains to be determined. With the inclusion of parametric variation, the eigenvalues of a chaotic system are supposed to move along smoothly varying curves of the type shown in the upper square of Fig. (2). Many of the previous studies of parametric variation focussed on the properties of such level curves. A great deal is known about the distribution of level velocities , the decay of correlations in parametric statistics , the distribution of level curvatures , and the statistics of the occurrences of avoided crossings . We now superpose the strength function overlap intensity information on Fig. (2) in the lower square as vertical lines centered on the levels; the lengths are scaled by the intensities (3-D versions of this figure turned out not to be very helpful). By considering the full strength function and not just the level curves (i.e. density of states), the eigenstate properties can be more directly probed. A new class of statistical measures can be defined that cross correlate intensities with levels. The most evident examples are the four correlation coefficients involving both level curves and eigenstate amplitudes that can be defined from the following quantities: (i) the level velocities, $`E_n(\lambda )/\lambda `$, (ii) level curvatures, $`^2E_n(\lambda )/\lambda ^2`$, (iii) overlaps, $`p_{\alpha n}`$, and (iv) overlap changes, $`p_{\alpha n}/\lambda `$. The most important is the overlap intensity-level velocity correlation coefficient, $`𝒞_\alpha (\lambda )`$, which is defined as $$𝒞_\alpha (\lambda )=\frac{p_{\alpha n}\frac{E_n(\lambda )}{\lambda }_E}{\sigma _\alpha \sigma _E}$$ (5) where $`\sigma _\alpha ^2`$ and $`\sigma _E^2`$ are the local variances of the overlaps and level velocities, respectively. The brackets denote a local energy average in the neighborhood of $`E`$. It weights most the level velocities whose associated eigenstates possibly share common localization characteristics and measures the tendency of these levels to move in a common direction. In this expression, the phase space volume remains constant so that the level velocities are zero-centered (otherwise the mean must be subtracted), and $`𝒞_\alpha (\lambda )`$ is rescaled to a unitless quantity with unit variance making it a true correlation coefficient. The set of states included in the local energy averaging can be left flexible except for a few constraints. Only energies where $`\overline{S}_\alpha (E,\lambda )`$ is roughly constant can be used or some intensity unfolding must be applied. Also, the energy range must be small so that the classical dynamics are essentially the same throughout the range, but it must also be broad enough to include several eigenstates. $`𝒞_\alpha (\lambda )`$ thus has a simple form and the additional advantage of involving quantities of direct physical interest. Level velocities (curvatures also) arise in thermodynamic properties of mesoscopic systems , and overlap intensities often arise in the manner used to couple into the system . It is the most sensitive measure of the four possible combinations, the others being the intensity-curvature, intensity change - curvature and intensity change - level velocity correlation coefficients. The first two are far less sensitive measures of eigenstate localization effects, even though curvature distributions are affected by localization because of the relative rareness of being near avoided crossings where curvatures are large. The last shows no effect since intensities will change whether the level is moving up or down. These three measures will not be considered further in this paper, but we did calculate them to verify their lack of sensitivity. ## III Ergodicity, Random Waves, and Random Matrix Theory Semiclassical expressions for wavefunctions have the form $$\mathrm{\Psi }(𝐱)=\underset{n}{}A_n(𝐱)\mathrm{exp}\left(iS_n(𝐱)/\mathrm{}i\nu _n\pi /2\right)$$ (6) where $`S_n(𝐱)`$ is a classical action, $`\nu _n`$ is a phase index, and $`A_n(𝐱)`$ is a slowly varying function given by the square root of a classical probability. The classical trajectory underlying each term arrives at the point $`𝐱`$ with momentum, $`𝐩_n=S_n(𝐱)`$. For a chaotic system, a complete theory leading to an equation of the form of Eq. (6) does not exist . Nevertheless, Berry conjectured that for the purposes of understanding the statistical properties of chaotic eigenfunctions, the ergodic hypothesis implies that the true eigenfunction will appear statistically equivalent to a large sum of these terms each arriving with a random phase (since each wave contribution extends over a complicated, chaotic path). For systems whose Hamiltonian is a sum of kinetic and potential energies, the energy surface constraint $`\delta (EH(𝐩,𝐪))`$ fixes only the magnitude of the wavevector. The eigenfunctions therefore appear locally as a sum of randomly phased plane waves pointing in arbitrary directions with fixed wavevector $`k`$. The central limit theorem asserts such waves are Gaussian random. An example is shown in Fig. (3) for a two-degree-of-freedom system where the spatial correlations fall off as a Bessel function, $`J_0(kr)`$. If the eigenstates truly possessed these characteristics, then a perturbation of the Hamiltonian would have matrix elements that behaved as Gaussian random variables whose variance depended only on the energy separation of the two eigenstates, i.e. an energy-ordered, banded random matrix. The energy ordering separates the weakly interacting states, and therefore only the local structure is of importance here. The range of the averaging carried out in the correlation function is taken to be much less than the bandwidth of such a random matrix. The ultimate statistical expression of this structure is embodied in one of the standard Gaussian ensembles (GE). We construct a parametrically varying ensemble $`\{\widehat{H}(\lambda )\}`$ as $$\widehat{H}(\lambda )=\widehat{H}_0+\lambda \widehat{H}_1$$ (7) where $`\widehat{H}_0`$ and $`\widehat{H}_1`$ are independently chosen GE matrices. Note that the sum of two GE matrices is also a GE matrix which thus satisfies our desire to consider stationary statistical properties as $`\lambda `$ varies. It is unnecessary to specify the abstract vector space of $`\{\widehat{H}(\lambda )\}`$ (only the dimensionality of the space) in the definition of the ensemble. However, $`\|\alpha `$ has to be overlapped with the eigenstates, and thus a localized wave packet seemingly must be specified. In fact, the specific choice is completely irrelevant because the GEs are invariant under the set of transformations that diagonalize them. $`\|\alpha `$ can be taken as any fixed vector in the space by invariance. The overlaps and level velocities turn out to be independent over the ensemble since diagonalizing $`\{\widehat{H}_0\}`$ leaves $`\{\widehat{H}_1\}`$ invariant and the level velocities are equal to the diagonal matrix elements of $`\widehat{H}_1`$. With the overbar denoting ensemble averaging, $$\overline{𝒞_\alpha (\lambda )}=\frac{\overline{p_{\alpha n}\frac{E_n}{\lambda }}_E}{\sigma _\alpha \sigma _E}=\frac{\overline{p_{\alpha n}}_E\overline{\frac{E_n}{\lambda }}_E}{\sigma _\alpha \sigma _E}=0$$ (8) In fact, it is essential to keep in mind that every choice of $`\|\alpha `$ gives zero correlations within the random matrix framework. The existence of even a single $`\|\alpha `$ in a particular system that leads to nonzero correlations violates ergodicity. It is straightforward to go further and consider the mean square fluctuations of $`𝒞_\alpha (\lambda )`$, $`\overline{𝒞_\alpha (\lambda )^2}`$ $`=`$ $`{\displaystyle \frac{\overline{\left(p_{\alpha i}\frac{E_i}{\lambda }_E\right)^2}}{\left(\sigma _\alpha \sigma _E\right)^2}}`$ (9) $`=`$ $`{\displaystyle \frac{1}{\left(N\sigma _\alpha \sigma _E\right)^2}}{\displaystyle \underset{i}{\overset{N}{}}}{\displaystyle \underset{j}{\overset{N}{}}}\overline{p_{\alpha i}{\displaystyle \frac{E_i}{\lambda }}p_{\alpha j}{\displaystyle \frac{E_j}{\lambda }}}`$ (10) $`=`$ $`{\displaystyle \frac{1}{\left(N\sigma _\alpha \sigma _E\right)^2}}{\displaystyle \underset{i}{\overset{N}{}}}{\displaystyle \underset{j}{\overset{N}{}}}\overline{p_{\alpha i}p_{\alpha j}}\overline{{\displaystyle \frac{E_i(\lambda )}{\lambda }}{\displaystyle \frac{E_j(\lambda )}{\lambda }}}`$ (11) $`=`$ $`{\displaystyle \frac{1}{\left(N\sigma _\alpha \sigma _E\right)^2}}{\displaystyle \underset{i}{\overset{N}{}}}\overline{p_{\alpha i}^2}\overline{\left({\displaystyle \frac{E_i(\lambda )}{\lambda }}\right)^2}={\displaystyle \frac{1}{N}}`$ (12) where $`N`$ is the effective number of states used in the energy averaging. Again the level velocities are independent of the eigenvector components. The $`E_j(\lambda )/\lambda =j\|\widehat{H}_1\|j`$ and thus the $`ij`$ terms vanish due to the independence of the diagonal elements of the perturbation leaving only the diagonal terms that involve the quantities that respectively enter the variance of the eigenvector components and the mean square level velocity. The final result reflects the equivalence of ensemble and spectral averaging in the large-$`N`$ limit. Therefore, in ergodically behaving systems, $`𝒞_\alpha (\lambda )=0\pm N^{1/2}`$ for every choice of $`\|\alpha `$. Fig. (2) was made using the orthogonal GE. It illustrates a manifestation of ergodicity, i.e. universal response of the quantum levels with respect to $`\lambda `$ and democratic behavior of the overlap intensities. ## IV Semiclassical Dynamics We develop a theory based upon semiclassical dynamics which explains how nonzero overlap correlation coefficients arise out of the localization properties of the system. The theory simply reflects the quantum manifestations of finite time correlations in the classical dynamics. In a chaotic system, the classical propagation of $`\rho _\alpha (𝐩,𝐪)`$ will relax to an ergodic long time average. However, wave packet revivals in the corresponding quantum system earlier than this relaxation time can occur . In Heller’s original treatment of scars , he uses arguments based upon these recurrences which occur at finite times to infer localization in the eigenstates. In the correlation function, the intensities, $`p_{\alpha n}`$, weight most heavily the level motions of the group of eigenstates localized near $`\rho _\alpha (𝐩,𝐪)`$, if indeed such eigenstates exist. If we construct the Hamiltonian as in Eq. (7) where $`\widehat{H}_0`$ is the unperturbed part, then by first-order perturbation theory, the level velocities are the diagonal matrix elements of $`\widehat{H}_1`$ just as in random matrix theory. We showed in the previous section that in random matrix theory these elements weighted with the intensities are zero centered. For a general quantum system the equivalent expectation would be fluctuations about the corresponding classical average of the perturbation over the microcanonical energy surface, $`\delta (EH(𝐩,𝐪))`$. In this case, $`𝒞_\alpha (\lambda )0`$ for all $`\|\alpha `$. On the other hand, the quantum system will fluctuate differently if there is localization in the eigenstates. Note that this means some choices of $`\|\alpha `$ will still lead to zero correlations. It only takes one statistically significant nonzero result to demonstrate localization conclusively, but to obtain a complete picture, it is necessary to consider many $`\|\alpha `$ covering the full energy surface. We begin by examining the individual components of the overlap correlation coefficient, the level velocities and intensities. Their $`\mathrm{}`$-dependences are derived and also they are shown to be consistent with random matrix theory as $`\mathrm{}0`$. Finally, the weighted level velocities are discussed. We give an estimate based upon a semiclassical theory involving homoclinic orbits for the slope of the large intensities. ### A Level velocities In random matrix theory (RMT) level velocities are Gaussian distributed as would also be expected of a highly chaotic system in the small $`\mathrm{}`$ limit. Thus, the mean and variance, $`\sigma _E^2`$, give a complete statistical description in the limiting case and are the most important quantities more generally. Since the purpose of this section is to derive their scaling properties, it is better to work with dimensionless quantities. Thus, the dimensionless variance is defined as $`\stackrel{~}{\sigma }_E^2\overline{d}^2(E,\lambda )\sigma _E^2`$ where $`\overline{d}(E,\lambda )`$ is the mean level density which is the reciprocal of the mean level spacing. We begin by following arguments originally employed by Berry and Keating in which they investigated the level velocities normalized by the mean level spacing for classically chaotic systems with the topology of a ring threaded by quantum flux. In order to make the discussion self contained we will summarize their basic ideas using their notation and then extend their results to include level velocities for any classically chaotic system. More recently, Leboeuf and Sieber studied the non-universal scaling of the level velocities using a similar semiclassical theory. The $`\mathrm{}`$-dependence of the average and root mean square level velocities for an arbitrary parameter change is derived and is consistent with the previous works. The smoothed spectral staircase is $$N_ϵ(E,\lambda )=\underset{n}{}\theta _ϵ(EE_n(\lambda ))$$ (13) and taking the derivative with respect to the parameter, we obtain $$\frac{N_ϵ(E,\lambda )}{\lambda }=\underset{n}{}\delta _ϵ(EE_n(\lambda ))\frac{E_n(\lambda )}{\lambda }$$ (14) The quantity $`ϵ`$ is an energy smoothing term which will be taken smaller than the mean level spacing. Our calculations will use Lorentzian smoothing where $$\delta _ϵ(x)=\frac{ϵ}{\pi \left(x^2+ϵ^2\right)}$$ (15) The energy averaging of Eq. (14) yields $$\frac{N_ϵ(E,\lambda )}{\lambda }_E=\overline{d}(E,\lambda )\frac{E_n(\lambda )}{\lambda }_n$$ (16) Thus, in order to obtain information about the level velocities, we will evaluate the spectral staircase. The semiclassical construction of the spectral staircase is broken into an average part and an oscillating part $`N_ϵ(E,\lambda )`$ $`=`$ $`\overline{N}(E,\lambda )+{\displaystyle \underset{p}{}}B_p(E,\lambda )\mathrm{exp}\left\{i\left[{\displaystyle \frac{S_p(E,\lambda )}{\mathrm{}}}\right]\right\}`$ (18) $`\times \mathrm{exp}\left\{{\displaystyle \frac{ϵT_p(E,\lambda )}{\mathrm{}}}\right\}`$ The average staircase $`\overline{N}(E,\lambda )`$ is the Weyl term and to leading order in $`\mathrm{}`$ is given by $$\overline{N}(E,\lambda )=\frac{1}{h^d}\theta (EH(𝐩,𝐪;\lambda ))𝑑𝐩𝑑𝐪$$ (19) This simply states that each energy level occupies a volume $`h^d`$ in phase space. A change in the phase space volume will produce level velocities due to the rescaling. We wish to study level velocities created by a change in the dynamics, not the rescaling. Hence, without loss of generality we will require the phase space volume to remain unchanged, so $`\overline{N}(E,\lambda )/\lambda =0`$. The oscillating part of the spectral staircase is a sum over periodic orbits. In general, a perturbation will alter the value of the classical actions, $`S_p`$, the periods, $`T_p`$, and the amplitudes, $$B_p=\frac{\mathrm{exp}(i\nu _p)}{2\pi \sqrt{det(M_p1)}}$$ (20) where $`M_p`$ is the stability matrix and $`\nu _p`$ is the Maslov phase index. The summation is most sensitive to the changing actions and periods because of the associated rapidly oscillating phases, i.e. the division by $`\mathrm{}`$ in the exponential. Since the energy smoothing term, $`ϵ`$, is taken smaller than a mean level spacing, it scales at least by $`\mathrm{}^d`$ and the derivatives of the period vanish as $`\mathrm{}0`$. Thus, only the derivatives of the actions are considered, and the oscillating part of the staircase yields $`{\displaystyle \frac{N_{osc}(E,\lambda )}{\lambda }}_E`$ $`=`$ $`{\displaystyle \underset{p}{}}B_p\left[{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle \frac{S_p(E,\lambda )}{\lambda }}\right]`$ (23) $`\times \mathrm{exp}\left\{i\left[{\displaystyle \frac{S_p(E,\lambda )}{\mathrm{}}}\right]\right\}`$ $`\times \mathrm{exp}\left\{{\displaystyle \frac{ϵT_p(E,\lambda )}{\mathrm{}}}\right\}_E`$ It has been shown that the change in the action for a periodic orbit is $$\frac{S_p}{\lambda }=_0^{T_p}\frac{H(𝐩,𝐪;\lambda )}{\lambda }𝑑t$$ (24) The above integral is over the path of the unperturbed orbit and the Hamiltonian can have the form of Eq. (7) where $`H_0`$ is the unperturbed part. Eq. (23) can be solved without the explicit knowledge of the periodic orbits in the $`\mathrm{}0`$ limit. The quantity $`S_p/\lambda `$ is replaced by its average. By the principle of uniformity , the collection of every periodic orbit covers all of phase space with a uniform distribution. Thus, the time integral can be replaced by an integral over phase space upon taking the average, $`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}{\displaystyle \frac{S_p}{\lambda }}_p`$ $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \frac{H(𝐩,𝐪;\lambda )}{\lambda }}`$ (26) $`\times \delta (EH(𝐩,𝐪;\lambda ))d𝐩d𝐪`$ where $`V`$ is the phase space volume of the energy surface. The above treatment of the average is only valid for the long orbits, but we may ignore the finite set of short orbits in the sum for small enough $`\mathrm{}`$. $`H(𝐩,𝐪;\lambda )/\lambda `$ is the perturbation of the system that distorts the energy surface. Since the phase space is assumed to remain constant, then the average change in the actions of the periodic orbits is zero in the limit of summing over all the orbits. If only a finite number of orbits are considered, corresponding to a finite $`\mathrm{}`$, then there might be some residual effect of the oscillating part which will cause a deviation from RMT. Continuing to follow Berry and Keating, the mean square of the counting function derivatives can be expressed in terms of the level velocities $`\left({\displaystyle \frac{N_ϵ}{\lambda }}(E,\lambda )\right)^2_E`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \frac{E_n}{\lambda }}(\lambda ){\displaystyle \frac{E_m}{\lambda }}(\lambda )`$ (29) $`\times \delta _ϵ(EE_n(\lambda ))`$ $`\times \delta _ϵ(EE_m(\lambda ))_E`$ For a non-degenerate spectrum, the summation is non-zero only if $`n=m`$ because of the product of the two delta functions. Since Lorentzian smoothing is applied, then $$\delta _ϵ^2(x)\frac{1}{2\pi ϵ}\delta _{ϵ/2}(x)$$ (30) for $`ϵ\overline{d}^1`$. Thus we have $$\left(\frac{N_ϵ}{\lambda }(E,\lambda )\right)^2_E=\frac{\overline{d}}{2\pi ϵ}\left(\frac{E_n}{\lambda }(\lambda )\right)^2_n$$ (31) The final result will be independent of $`ϵ`$ and the type of smoothing, i.e. Lorentzian or Gaussian. Using the $`\lambda `$ derivative of Eq. (18), the dimensionless level velocities are $`\stackrel{~}{\sigma }_E^2`$ $`=`$ $`{\displaystyle \frac{2\pi ϵ\overline{d}}{\mathrm{}^2}}{\displaystyle \underset{p}{}}{\displaystyle \underset{p^{}}{}}\|B_pB_p^{}\|{\displaystyle \frac{S_p}{\lambda }}{\displaystyle \frac{S_p^{}}{\lambda }}`$ (32) $`\times `$ $`\mathrm{exp}\left\{i\left[{\displaystyle \frac{S_pS_p^{}}{\mathrm{}}}\right]\right\}\mathrm{exp}\left\{{\displaystyle \frac{ϵ}{\mathrm{}}}[T_p+T_p^{}]\right\}_E`$ (33) The diagonal and off-diagonal contributions are separated, so $$\stackrel{~}{\sigma }_E^2=\stackrel{~}{\sigma }_{E,diag}^2+\stackrel{~}{\sigma }_{E,off}^2$$ (34) As $`\mathrm{}0`$, the phase of the exponential oscillates rapidly and averages out to be zero unless $`S_p=S_p^{}`$. We will assume that this occurs rarely except when $`p=p^{}`$. The product $`(S_p/\lambda )(S_p^{}/\lambda )`$ can take on both positive and negative values. This also helps to reduce the contributions of the off-diagonal terms. For a more complete discussion of the diagonal vs. off-diagonal terms see . We will only present the results for the diagonal terms, since the correlations between the actions of different orbits is not known but should not alter the leading $`\mathrm{}`$-dependence. The diagonal contribution is $$\stackrel{~}{\sigma }_{E,diag}^2=\frac{2\pi ϵ\overline{d}g}{\mathrm{}^2}\underset{p}{}\|B_p\|^2\left(\frac{S_p}{\lambda }\right)^2\mathrm{exp}\left\{\frac{2ϵT_p}{\mathrm{}}\right\}_E$$ (35) The factor $`g`$ depends on the symmetries of the system. For systems with time-reversal invariance $`g=2`$ and without time-reversal symmetry $`g=1`$. The precise values of $`S_p/\lambda `$ are specific to each periodic orbit rendering the sum difficult to evaluate precisely. A statistical approach is possible though which generates a relationship between the sum and certain correlation decays. Hence, the quantity $`(S_p/\lambda )^2`$ in Eq. (35) is replaced by its average, $`\left({\displaystyle \frac{S_p}{\lambda }}\right)^2_p`$ $`=`$ $`{\displaystyle _0^T}{\displaystyle _0^T}{\displaystyle \frac{H(𝐩(t),𝐪(t);\lambda )}{\lambda }}`$ (36) $`\times `$ $`{\displaystyle \frac{H(𝐩(t^{}),𝐪(t^{});\lambda )}{\lambda }}_pdt^{}dt`$ (37) $`=`$ $`2{\displaystyle _0^T}{\displaystyle _t^T}{\displaystyle \frac{H(𝐩(t),𝐪(t);\lambda )}{\lambda }}`$ (38) $`\times `$ $`{\displaystyle \frac{H(𝐩(t^{}+t),𝐪(t^{}+t);\lambda )}{\lambda }}_pdt^{}dt`$ (39) Long orbits increasingly explore the available phase space on an ever finer scale. As the time between two points in a chaotic system goes to infinity, then they become uncorrelated from each other. This is a consequence of the mixing property, $$f(0)f(t)_p0$$ (40) This property is independent of the placement of the two points, i.e. the two points can lie on the same orbit as long as the time between the points increases to infinity. Thus, by the central limit theorem, $`S_p/\lambda `$ will be Gaussian distributed for the sufficiently long periodic orbits. The time dependence of Eq. (36) is approximated by a method discussed by Bohigas et al. . They define $$K(E)=_0^{\mathrm{}}\frac{H(𝐩(0),𝐪(0);\lambda )}{\lambda }\frac{H(𝐩(t),𝐪(t);\lambda )}{\lambda }_p𝑑t$$ (41) which can be evaluated in terms of properties of the perturbation. The variance of the actions in the limit of long periods becomes $$\left(\frac{S_p}{\lambda }\right)^2_p2K(E)T$$ (42) Applying the Hannay and Ozorio de Almeida sum rule , the following substitution is made $$\underset{p}{}\|B_p\|^2\mathrm{}\frac{1}{2\pi ^2}_0^{\mathrm{}}\frac{dT}{T}\mathrm{}$$ (43) Hence, the diagonal contribution is $`\stackrel{~}{\sigma }_{E,diag}^2`$ $``$ $`{\displaystyle \frac{ϵ\overline{d}g}{\pi \mathrm{}^2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{T}}\left(2K(E)T\right)\mathrm{exp}\left\{{\displaystyle \frac{2ϵT}{\mathrm{}}}\right\}𝑑T`$ (44) $``$ $`{\displaystyle \frac{gK(E)\overline{d}}{\pi \mathrm{}}}`$ (45) $``$ $`\mathrm{}^{(d+1)}`$ (46) The variance of the level velocities on the scale of a mean spacing grows $`\mathrm{}^1`$ faster than the density of states as the semiclassical limit $`(\mathrm{}0)`$ is approached; see numerical tests performed on the stadium in the next section. The exact level velocities are perturbation dependent and cannot be determined without specific knowledge of the system (i.e. the evaluation of $`K(E)`$). $`K(E)`$ is a classical quantity that contains dynamical information about the periodic orbits. It should scale as the reciprocal of the Lyapunov exponent . Leboeuf and Sieber derived $`K(E)`$ for billiards where the perturbation is a moving boundary. In this case $`K(E)`$ depends upon the autocorrelation function and the fluctuations of the number of bounces. For maps $`K(E)`$ is an action velocity diffusion coefficient . $`\{S_p/\lambda \}`$ being Gaussian distributed is linked to the level velocities being Gaussian distributed as in RMT. If the $`\{S_p/\lambda \}`$ are not Gaussian distributed by the Heisenberg time, then one should not expect the level velocities to be consistent with RMT; again see the stadium results ahead. ### B Overlap intensities Now we investigate the overlap intensities and derive a semiclassical expression for the $`\mathrm{}`$-scaling of the root mean square. Eckhardt et al. developed a semiclassical theory based on periodic orbits to obtain the matrix elements of a sufficiently smooth operator. However, the projection operator of interest here, $`\|\alpha \alpha \|`$, is not smooth on the scale of $`\mathrm{}`$ for Gaussian wave packets. Thus, their stationary phase approximations do not apply, in principle, to the oscillating part of the strength function. In Berry’s work on scars , he used Gaussian smoothing of the Wigner transform of the eigenstates to obtain a semiclassical expression for the strength of the scars. His approach led to a sum over periodic orbits. We will use the energy Green’s function similar to Tomsovic and Heller in where they derived the autocorrelation function using the time Green’s function and gave results for the strength functions as well. This technique results in a connection between the overlap intensities and the return dynamics, namely the homoclinic orbits. For completeness, we present the smooth part of the strength function which is easily obtained from the zero-length trajectories, $$\overline{S}_\alpha (E,\lambda )=\frac{1}{h^d}A(𝐪,𝐩)\delta (EH(𝐪,𝐩))𝑑𝐪𝑑𝐩$$ (47) $`A(𝐪,𝐩)`$ is the Wigner transform of the Gaussian wave packet and is given by $$A(𝐪,𝐩)=2^d\mathrm{exp}\{(𝐩𝐩_\alpha )^2\sigma ^2/\mathrm{}^2(𝐪𝐪_\alpha )^2/\sigma ^2\}$$ (48) The above results were previously used by Heller in the derivation the envelope of the strength function and does not contain any information about the dynamics of the system. The oscillating part of the strength function, on the other hand, includes dynamical information, $$S_{\alpha ,osc}(E,\lambda )=\frac{1}{\pi }\mathrm{I}m\alpha \|𝐪G(𝐪,𝐪^{};E)𝐪^{}\|\alpha 𝑑𝐪𝑑𝐪^{}$$ (49) where $`G(𝐪,𝐪^{};E)`$ $`=`$ $`{\displaystyle \frac{1}{i\mathrm{}(2\pi i\mathrm{})^{(d1)/2}}}{\displaystyle \underset{j}{}}\|D_s\|^{1/2}`$ (51) $`\times e^{i[S_j(𝐪,𝐪^{};E)/\mathrm{}\nu _j^{}\pi /2]}`$ is the semiclassical energy Green’s function. The above sum is over all paths that connect $`𝐪`$ to $`𝐪^{}`$ on a given energy surface $`E`$. The action is quadratically expanded about each reference trajectory; see Appendix A for details. The initial and final points ($`𝐪_i`$ and $`𝐪_f`$) of the reference trajectories are obtained by considering the evolution of the wave packet. Nearby points will behave similarly for short times. Thus, the phase space can be partitioned into connecting areas. As the time is increased the number of partitions grow and the size of their area shrinks. The reference trajectories are the paths that connect the partitions. The autocorrelation function in has the same form as Appendix A where the paths that contribute to the saddle points are the orbits homoclinic to the centroid of the Gaussian wave packet so that $`𝐪_i`$ and $`𝐪_f`$ lie on the intersections of the stable and unstable manifolds. The result from Appendix A is $`S_{\alpha ,osc}(E)`$ $`=`$ $`{\displaystyle \frac{\sigma }{\pi ^{1/2}\mathrm{}}}\mathrm{R}e{\displaystyle \underset{j}{}}\left({\displaystyle \frac{det\stackrel{~}{𝐀}^{21}}{det𝐀}}\right)^{1/2}`$ (52) $`\times `$ $`\left({\displaystyle \frac{1}{\|\dot{q}^{(N)}\|\|\dot{q}^{(N)}\|}}\right)^{1/2}f_j(𝐪_f,𝐪_i)e^{iS_j(𝐪_f,𝐪_i;E)/\mathrm{}}`$ (53) where $`f_j(𝐪_f,𝐪_i)`$ $`=`$ $`\mathrm{exp}\{{\displaystyle \frac{1}{4}}𝐛𝐀^1𝐛{\displaystyle \frac{i}{\mathrm{}}}𝐩_\alpha (𝐪_f𝐪_i)`$ (55) $`\times {\displaystyle \frac{(𝐪_f𝐪_\alpha )^2}{2\sigma ^2}}{\displaystyle \frac{(𝐪_i𝐪_\alpha )^2}{2\sigma ^2}}{\displaystyle \frac{i\nu _j^{}\pi }{2}}\}`$ The function $`f_j(𝐪_f,𝐪_i)`$ in the above equation is a damping term which depends on the end points of the homoclinic orbits. Only orbits which approach the center of the Gaussian wave packet in phase space will contribute to the sum. The time derivatives of the parallel coordinates are evaluated at the saddle points which are near the centroid of the Gaussian, so we may set $`\|\dot{q}^{(N)}\|\|\dot{q}^{(N)}\|\|p_\alpha \|/m`$. The sum over homoclinic orbits used for the autocorrelation function in converged well to the discrete quantum strength function when only those orbits whose period did not exceed the Heisenberg time, $`(\tau _H=2\pi \mathrm{}\overline{d}(E,\lambda ))`$, were included. As happened with the periodic orbits and the level velocities, in order to evaluate Eq. (52) the homoclinic orbits and their stabilities must be computed rendering the sum tedious to evaluate precisely as done in . By taking a statistical approach we can gain some insight into the workings of this summation. The variance of the intensities are obtained by a similar fashion as the level velocities. Using Eq. (1) and Eq. (30), we have $$S_{\alpha ,osc}^2(E,\lambda )_E=\frac{\overline{d}}{2\pi ϵ}\sigma _\alpha ^2$$ (56) Since the square of the strength function is a product of two delta functions, an energy smoothing term is required. After making the diagonal approximation, we obtain $`\sigma _{\alpha ,diag}^2`$ $`=`$ $`{\displaystyle \frac{2\pi ϵg}{\overline{d}}}{\displaystyle \underset{j}{}}{\displaystyle \frac{m^2\sigma ^2}{\pi \mathrm{}^2}}\left\|{\displaystyle \frac{det\stackrel{~}{𝐀}^{21}}{det𝐀}}\right\|`$ (58) $`\times {\displaystyle \frac{\|f_j(𝐪_f,𝐪_i)\|^2}{\|p_\alpha \|^2}}e^{2ϵT_j/\mathrm{}}_E`$ A classical sum rule is applied to the above sum for special cases including two-dimensional systems; see Appendix B for the details. Thus, $`\sigma _{\alpha ,diag}^2`$ $``$ $`{\displaystyle \frac{2ϵgm^2\sigma ^2}{\mathrm{}^2\overline{d}\|p_\alpha \|^2}}{\displaystyle \mathrm{exp}\{2ϵT/\mathrm{}\}𝑑T}`$ (59) $``$ $`{\displaystyle \frac{gm^2\sigma ^2}{\mathrm{}\overline{d}\|p_\alpha \|^2}}`$ (60) Setting $`\sigma \mathrm{}^{1/2}`$ which shrinks the momentum and position uncertainties similarly, the $`\mathrm{}`$-scaling of $`\sigma _{\alpha ,diag}^2`$ is $`\mathrm{}^d`$; see numerical tests of the stadium in the next section. Assuming that the amplitudes of the wavefunctions are Gaussian random, then the RMT result for strength functions is a Porter-Thomas distribution which has a variance that is proportional to the square of its average. The average strength function, Eq. (47), scales as $`(\sigma /\mathrm{})^d`$ for Hamiltonians which can be locally expanded as a quadratic. Therefore, with $`\sigma \mathrm{}^{1/2}`$ the variance of the strength function, Eq. (59), scales as the square of the average and is consistent with RMT. ### C Weighted level velocities A semiclassical treatment of the overlap correlation coefficient defined in Eq. (5) is now developed. As stated in the introduction, the companion paper presents the semiclassical theory for maps. We stress that in the preceding subsections and in what follows is for conservative Hamiltonian systems. Here, the $`\mathrm{}`$-dependence of the average overlap correlation coefficient is established and a semiclassical argument for the existence of nonzero correlations is presented. #### 1 Actions of homoclinic orbits To calculate the overlap correlation coefficient, the rate of change of the actions for homoclinic orbits will be necessary. As discussed earlier, this was accomplished for periodic orbits . We extend these results to include the actions of homoclinic orbits. Homoclinic orbits have infinite periods causing their actions to become infinite. We are interested in the limiting difference of the action, $`𝒮_j^{(p)}`$, between the $`j^{th}`$ homoclinic orbit and repetitions of its corresponding periodic orbit $`p`$. The difference is finite and is equal to the area bounded by the stable and unstable manifolds with intersection at the $`j^{th}`$ homoclinic point in a Poincaré map. $`𝒮_j^{(p)}`$ provides information about the additional phase gathered by the homoclinic orbit. The action of the $`j^{th}`$ homoclinic orbit as $`n\mathrm{}`$ in the time interval $`(nT_p,nT_p)`$ is $$S_{n,j}^{(p)}2nS_p+𝒮_j^{(p)}$$ (61) where $`T_p`$ and $`S_p`$ are the period and action of the periodic orbit, respectively. As a consequence of the Birkhoff-Moser theorem , if the Poincaré map is invertible and analytic, then there exist infinite families of periodic orbits that accumulate on a homoclinic orbit. It is thus possible to estimate the action of the homoclinic orbit by these periodic orbits whose action is given by $$\alpha _{n,j}^{(p)}=nS_p+𝒮_j^{(p)}s_{n,j}^{(p)}$$ (62) where $`s_{n,j}^{(p)}`$ is the difference in action between a path defined by $`𝒮_j^{(p)}`$ along the stable and unstable manifolds and the path of the new periodic orbit in a Poincaré map. $`s_{n,j}^{(p)}`$ depends exponentially on $`n`$, so as $`n\mathrm{}`$, $`\alpha _{2n,j}^{(p)}`$ approaches the action of the homoclinic orbit. Thus, in the limit of large $`n`$, $`𝒮_j^{(p)}`$ is approximated by the difference between two periodic orbits (i.e. $`𝒮_j^{(p)}\alpha _{n,j}^{(p)}nS_p`$). Hence, the change in $`𝒮_j^{(p)}`$ due to a small perturbation is calculated as in , $`\mathrm{\Delta }𝒮_j^{(p)}`$ $`=`$ $`\mathrm{\Delta }\lambda {\displaystyle _{\alpha _{n,j}^{(p)}}}{\displaystyle \frac{H(𝐩,𝐪;\lambda )}{\lambda }}𝑑t`$ (64) $`+n\mathrm{\Delta }\lambda {\displaystyle _{S_p}}{\displaystyle \frac{H(𝐩,𝐪;\lambda )}{\lambda }}𝑑t+O(\lambda ^2)`$ where the integrals are over the unperturbed periodic orbits. The differences, $`s_{n,j}^{(p)}`$, can be made smaller than the second order term in Eq. (64) by taking $`n`$ large enough. Interchanging the order of integration and differentiation, the integrals reduce to the unperturbed energy times the derivative of the orbit period with respect to the parameter, $$\mathrm{\Delta }𝒮_j^{(p)}\mathrm{\Delta }\lambda E_p\frac{}{\lambda }\left(T_{\alpha _{n,j}^{(p)}}nT_p\right)$$ (65) The orbit period $`T`$ can be expressed as $`S/E`$. Thus, the difference of the two periods as $`n\mathrm{}`$ is $$T_{\alpha _{n,j}^{(p)}}nT_p=\frac{\alpha _{n,j}^{(p)}}{E}n\frac{S_p}{E}\frac{𝒮_j^{(p)}}{E}$$ (66) Hence, $`\mathrm{\Delta }𝒮_j^{(p)}`$ $``$ $`\mathrm{\Delta }\lambda {\displaystyle \frac{}{\lambda }}\left(E_p{\displaystyle \frac{𝒮_j^{(p)}}{E}}\right)`$ (67) $``$ $`\mathrm{\Delta }\lambda {\displaystyle _{𝒮_j^{(p)}}}{\displaystyle \frac{H(𝐩,𝐪;\lambda )}{\lambda }}𝑑t`$ (68) Note that the integral is over the unperturbed path along the stable and unstable manifolds. For zero correlations, $`\mathrm{\Delta }𝒮_j^{(p)}`$ must be “randomly” distributed about zero. If enough time is allowed, then for ergodic systems the set of all homoclinic orbits for a given energy will come arbitrarily close to any point in phase space on that energy surface. Thus, an integral over phase space on the original energy surface can be substituted for the time integral, $`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{T}}\mathrm{\Delta }𝒮_j^{(p)}_j`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }\lambda }{V}}{\displaystyle \delta \left(EH(𝐩,𝐪;0)\right)}`$ (70) $`\times {\displaystyle \frac{H(𝐩,𝐪;\lambda )}{\lambda }}d𝐩d𝐪`$ Since the density of states are kept constant, the perturbations fluctuate about zero and the integral vanishes. Hence, the average change in the actions will be zero. The mean square fluctuations of the actions, $`\left(\mathrm{\Delta }𝒮_j^{(p)}\right)^2_j`$ $`=`$ $`\left(\mathrm{\Delta }\lambda \right)^2{\displaystyle _0^{\tau _H}}{\displaystyle _0^{\tau _H}}{\displaystyle \frac{H(𝐩,𝐪;\lambda )}{\lambda }}`$ (72) $`\times {\displaystyle \frac{H(𝐩^{},𝐪^{};\lambda )}{\lambda }}_jdtdt^{}`$ approaches a Gaussian distribution by the central limit theorem via the same reasoning as that for periodic orbits. Again we can define $`K_{hom}(E)`$ as in Eq. (41), except now the average is over homoclinic orbits and instead of integrating to infinity we only integrate to the Heisenberg time to be consistent with the range of the sum in Eq. (52). It is the short time dynamics that dominate. Long time correlations will average to zero by the mixing property (Eq. (40)). Thus, the variance of the actions becomes $$\left(\frac{𝒮_j^{(p)}}{\lambda }\right)^2_j2K_{hom}(E)T$$ (73) $`K_{hom}(E)`$ will approach $`K(E)`$ in the semiclassical limit ($`\tau _H\mathrm{}`$). #### 2 Overlap intensity-level velocity correlation coefficient In the previous two subsections, we have examined the pieces that constitute the overlap correlation coefficient. The semiclassical theories of the level velocities and the intensities are now combined to construct a semiclassical theory for the weighted level velocities. The numerator of the overlap correlation coefficient is proportional to the energy averaged product of the intensities and level velocities, $`S_\alpha (E,\lambda ){\displaystyle \frac{N(E,\lambda )}{\lambda }}_E`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}p_{\alpha n}(\lambda ){\displaystyle \frac{E_m}{\lambda }}`$ (75) $`\times \delta _ϵ(EE_n)\delta _ϵ(EE_m)_E`$ $`=`$ $`{\displaystyle \frac{\overline{d}}{2\pi ϵ}}p_{\alpha n}(\lambda ){\displaystyle \frac{E_n}{\lambda }}_n`$ (76) $`=`$ $`{\displaystyle \frac{\overline{d}}{2\pi ϵ}}\stackrel{~}{𝒞}_\alpha (\lambda )`$ (77) Lorentzian smoothing was again employed, Eq. (30), and we’ve defined $`\stackrel{~}{𝒞}_\alpha (\lambda )`$ to be the numerator of the overlap correlation coefficient (without the division of the rms level velocities and intensities). By the definition of the overlap correlation coefficient only the oscillating part of the level velocities and the intensities are considered. Using the derivative with respect to lambda of Eq. (18) and Eq. (52) the numerator becomes $`\stackrel{~}{𝒞}_\alpha (\lambda )`$ $`=`$ $`{\displaystyle \frac{2\pi ϵ}{\overline{d}}}\mathrm{R}e{\displaystyle \underset{j}{}}{\displaystyle \underset{p}{}}{\displaystyle \frac{m\sigma B_p}{\pi \mathrm{}^2}}\left({\displaystyle \frac{det\stackrel{~}{𝐀}^{21}}{det𝐀}}\right)^{1/2}`$ (79) $`\times {\displaystyle \frac{f_j(𝐪_f,𝐪_i)}{\|p_\alpha \|}}\left({\displaystyle \frac{S_p}{\lambda }}\right)e^{i(S_jS_p)/\mathrm{}ϵ(T_j+T_p)/\mathrm{}}_E`$ Because of the rapidly oscillating phases, the energy averaging will result in zero unless $`S_jS_p`$. As stated earlier, for every homoclinic orbit there is a periodic orbit that comes infinitesimally close to it. The same periodic orbit’s action can be nearly equal to the actions of different segments of the same homoclinic orbit. Thus, a diagonal approximation is used for the homoclinic segments $`\stackrel{~}{𝒞}_\alpha (\lambda )`$ $``$ $`{\displaystyle \frac{2\pi ϵg}{\overline{d}}}\mathrm{R}e{\displaystyle \underset{j}{}}{\displaystyle \frac{m\sigma B_j}{\pi \mathrm{}^2}}\left({\displaystyle \frac{det\stackrel{~}{𝐀}^{21}}{det𝐀}}\right)^{1/2}e^{2ϵT_j/\mathrm{}}`$ (81) $`\times \left({\displaystyle \frac{S_j}{\lambda }}\right){\displaystyle \frac{f_j(𝐪_f,𝐪_i)}{\|p_\alpha \|}}_E`$ Upon applying the sum rule for two-dimensional systems and other special cases, Eq. (B9), we have $`\stackrel{~}{𝒞}_\alpha (\lambda )`$ $``$ $`{\displaystyle \frac{ϵgm\sigma }{\pi \overline{d}\|p_\alpha \|}}{\displaystyle \frac{1}{\mathrm{}^2}}{\displaystyle _0^{\mathrm{}}}e^{2ϵT/\mathrm{}}`$ (83) $`\times \left({\displaystyle \frac{S_j}{\lambda }}\right)f_j(𝐪_f,𝐪_i)_jdT`$ The changes in action of the homoclinic excursions are now weighted by the $`f_j(𝐪_f,𝐪_i)`$’s. Without the additional weighting the average in the changes in the action would be zero for all positions of the Gaussian wave packet. In , a heuristic argument for the direction of the weighted level velocities was given. The argument basically states that the energy surface changes with the parameter such that the action changes are minimized. Eq. (83) differs from in that the proper weightings, $`f_j(𝐪_f,𝐪_i)`$, of the homoclinic orbits are derived here, and the action changes are not correlated with the inverse periods. Also, in the homoclinic orbits were strictly cut-off at the Heisenberg time whereas here there is an exponential decay on the order of the Heisenberg time with the energy smoothing term, $`ϵ`$, equal to $`\mathrm{}/\tau _H`$ . One reason that reported such good results is that since the number of homoclinic segments proliferate exponentially, most of the included segments occurred near the Heisenberg time and the expression in Eq. (83) is divided by the Heisenberg time (i. e. multiplied by $`ϵ`$). As $`\mathrm{}0`$ $`(\tau _H\mathrm{})`$, the integral in Eq. (83) would be dominated by the subset of $`\{S_j/\lambda \}`$ associated with very long orbits and would decouple from the weightings. For small enough $`\mathrm{}`$, as previously stated, the $`S_j/\lambda `$ for these orbits would approach a zero-centered Gaussian density, and the integral would vanish. In other words, we could use arguments analogous to those underlying Eq. (26) to write $`\stackrel{~}{𝒞}_\alpha (\lambda )`$ $``$ $`{\displaystyle \frac{gm\sigma }{\pi \overline{d}\|p_\alpha \|}}{\displaystyle _0^{\mathrm{}}}𝑑Te^{2T/\tau _H}{\displaystyle \frac{T}{\tau _H}}{\displaystyle \frac{F}{V}}`$ (85) $`\times {\displaystyle }d𝐩d𝐪{\displaystyle \frac{H(𝐩,𝐪;\lambda )}{\lambda }}\delta (EH(𝐩,𝐪;\lambda ))`$ where $`F`$ is the phase space average of the weightings. Note that the phase space average of $`H(𝐩,𝐪;\lambda )/\lambda `$ vanishes excluding an irrelevant drift of levels, so the RMT prediction of $`\stackrel{~}{𝒞}_\alpha (\lambda )`$ is recovered for $`\mathrm{}`$ small enough. The leading order in $`\mathrm{}`$ correction to this is more difficult to ascertain. $`\mathrm{}`$ enters into the exponential in the integral for the energy smoothing, but not for the classical decay of the action changes. Upon taking the integral, this yields two competing terms for the $`\mathrm{}`$-scaling which may depend upon the region of phase space the correlation is taken in. The numerics also show a large fluctuation of the scaling in the stadium (see the next section). ## V Stadium Billiard In this section the semiclassical theories just presented and the numerical results from the stadium billiard are compared. The stadium billiard, which was proven by Bunimovich to be classically chaotic, has become a paradigm for studies of quantum chaos. It is defined as a two-dimensional infinite well with the shape pictured in Fig. (4). We continuously vary the side length, $`2\lambda `$, while altering the radii of the endcaps, $`R`$, to keep the area of the stadium a constant. Throughout this section, the level velocities and intensities are evaluated for a stadium with $`\lambda =R=1`$. For billiards the average number of states below a given energy, $`E`$, is approximately $`\overline{N}(E)mAE/2\pi \mathrm{}^2`$ where $`A`$ is the area of the billiard. This is the first term in an asymptotic series in powers of $`\mathrm{}`$. The density of the states, $`d\overline{N}/dE`$, is then a constant not depending on $`\lambda `$ to the lowest power of $`\mathrm{}`$ if the area remains the same. We will examine three different energy regimes for the stadium. Since billiards are scaling systems, this will correspond to three different values of $`\mathrm{}`$. The energy regimes are separated by a factor of four in energy or conversely a factor of one half in $`\mathrm{}`$. Twice as many states are taken in each successive energy regime so that the averages will incorporate the same relative size interval in energy as $`\mathrm{}`$ is decreased. This corresponds to the increase in the density of states for varying $`\mathrm{}`$. The distributions of the level velocities for all three energy regimes are shown in Fig. (5) along with the random matrix theory prediction. The skewness occurs because of a class of marginally stable orbits in the stadium. These orbits are the bouncing ball orbits which only strike the straight edges. Their contribution do not seem to decrease as the semiclassical limit is approached though they should once $`\mathrm{}`$ is sufficiently small. There is no clear trend for the level velocity distribution to approach Gaussian behavior. The root mean square of the level velocities also deviates from our calculations of the $`\mathrm{}`$-scaling in Section IV (Fig. (6)). This is again explained by the bouncing ball orbits whose effects are missing from the trace formula. Quantizing only these orbits using WKB yields a dimensionless level velocity scaling of $`\mathrm{}^2`$, while the trace formula gives a scaling of $`\mathrm{}^{3/2}`$. The numerical results give a scaling of approximately $`\mathrm{}^{1.8}`$ which lies in between the two suggesting that the marginally stable orbits significantly effect the level velocities. To study the intensities, the eigenstates must first be constructed. Bogomolny’s transfer operator method was used to find the eigenstates. This method uses a $`(d1)`$ dimensional surface of section. A convenient choice is the boundary of the stadium (Fig. (4)). The generation of a full phase space picture of the stadium would otherwise require four dimensions, two positions and two momenta. The position coordinate is measured along the perimeter and the momentum coordinate is defined by $`\mathrm{cos}\theta `$. The classical dynamics have a quantum analog that uses source points on the boundary. Thus, all of the eigenfunction’s localization behavior can be explored using wave packets defined in these coordinates. A coherent state on the boundary is a one-dimensional Gaussian wave packet; see the lower figure in Fig. (7). The corresponding wave packet in the interior of the stadium can be generated by a Green’s function and is shown in the upper figure of Fig. (7). For billiards the Green’s function is proportional to a zero<sup>th</sup> order Hankel function of the first kind, $`H_0^{(1)}(kr)/2i\mathrm{}^2`$. The centroid of the Gaussian wave packet is moved along the boundary and its momentum is changed according to the Birkhoff coordinate system. Thus, the entire phase space of the stadium is explored. The results for the average and the standard deviation of the intensities using Birkhoff coordinates are shown in Figs. (8) and (9), respectively. The average is flat except for peaks associated with the two symmetry lines of the stadium. The eigenstates used here were even-even states, so there is twice the intensity along the two symmetry lines that bisect the end caps and the straight edges. The standard deviation has two large peaks centered around the bouncing ball orbits. The rest of the figure is relatively flat with a few small bumps. Random matrix theory would predict this to be a flat figure with small oscillations. The marginal stability of the bouncing ball orbits can be seen but no other feature of the stadium, except for the horizontal bounce, is picked out by looking at the intensities. Fig. (10) shows the $`\mathrm{}`$-scaling of the root mean square for the intensities where the wave packet is placed on various periodic orbits. The theory from Section IV predicts a smaller scaling than the numerical results of the stadium. The heights of the bouncing ball peaks can be approximated by quantizing the rectangular region of the stadium. The intensities obtained from this calculation are weighted by the ratio of the density of the bouncing ball states to the total density of states. The Gaussian wave packet is placed in the center of the straight edge and the middle energy regime is used. The results of this approximation are 80.6 and 320.5 for the average intensity and rms intensity, respectively, compared to 80.7 and 386.3 for the numerical calculations of the stadium. Random matrix theory suggests that the correlation coefficient for a generic chaotic system should result in zero. On the other hand, using the correlation coefficient for the stadium in Birkhoff coordinates, we found that some of the states gave nonzero correlations, Fig. (11). In fact, large correlations are found for nearly all the states in the stadium billiard which means that there exists phase space localization for most of the states. The large positive values of the correlation coefficient in the center of the figure again correspond to the bouncing ball states. Classically, this area of phase space is difficult to enter and leave. Hence, the localization is expected to be stronger for this area of phase space. The area beneath the peaks is several standard deviations $`(N^{1/2}=(114)^{1/2}0.09)`$ away from zero as predicted by random matrix theory. Thus, phase space localization is also occurring in this region. The point exactly in between the peaks is the point in phase space associated with the horizontal bounce. The series of smaller peaks leading up to the large peaks are the gateways into the vertical bouncing ball area. Fig. (12) is a plot of the orbits corresponding to these peaks. They are periodic orbits which only strike the endcaps twice and become almost vertical. Orbits must pass through these regions in order to enter or exit the vertical bounce states. As the energy of the system is increased (i. e. $`\mathrm{}`$ is decreased), the results of the correlation function remain qualitatively the same, Fig. (13). All the peaks and valleys stay in the same place. The numerical results of the overlap correlation coefficient fluctuate depending upon the area of phase space being considered. This is consistent with the the semiclassical theory in Section IV. More details of the system are explored as $`\mathrm{}`$ is decreased, since the phase space is divided into finer areas. Thus, more detailed information about the phase space localization of the system is observed in the overlap correlation coefficient at smaller values of $`\mathrm{}`$. ## VI Conclusions We have shown that intensity weighted level velocities are a good measure of the localization properties for chaotic systems. They are far more sensitive to localization than similarly weighted level curvatures (which are closely related to level statistics). Thus, a system can be RMT-like, yet the eigenstates are not behaving ergodically (as RMT predicts). The stadium eigenstates show a great deal of localization. Not only are the vertical bouncing ball orbits predicted by the measure, but also other orbits. The overlap correlation coefficient is very parameter dependent. Choosing a different parameter to vary would highlight other sets of orbits depending on how strong the perturbation effects those orbits. The degree of localization can be predicted by the return dynamics. In a chaotic system, all the return dynamics can be organized by the homoclinic orbits. The manner in which a chaotic system’s eigenstates approach ergodicity as $`\mathrm{}0`$ will depend on a new time scale, i.e. that required for the homoclinic excursions to explore the available phase space fully. Parametrically varied data exist that can be analyzed in this way. In the Coulomb-blockade conductance data to the extent that the resonance energy variations are related to a single particle level velocity (minus a constant charging energy and absent residual interaction effects) should show correlations. We mention also that the microwave cavity data can be studied with even more flexibility since they have measured the eigenstates and can therefore meticulously study a wide range of $`\|\alpha `$ to get a complete picture of the eigenstate localization properties. Finally, this analysis could be applied in a very fruitful way to near-integrable and mixed phase space systems. In these cases, standard random matrix theory would not give the zero<sup>th</sup> order statistical expectation, but the localization would still be determined by the return dynamics in the semiclassical approximation. ###### Acknowledgements. We gratefully acknowledge important discussions with B. Watkins and T. Nagano and support from the National Science Foundation under Grant No. NSF-PHY-9800106 and the Office of Naval Research under Grant No. N00014-98-1-0079. ## A Gaussian Integration Inserting Eq. (51) into Eq. (49), the strength function involves two $`N`$-dimensional integrals where $`N`$ is the system’s number of degrees of freedom, $`S_{\alpha ,osc}(E)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{I}m{\displaystyle \frac{1}{i\mathrm{}(2\pi i\mathrm{})^{(d1)/2}}}\left({\displaystyle \frac{1}{\pi \sigma ^2}}\right)^{d/2}`$ (A4) $`\times {\displaystyle }{\displaystyle \underset{j}{}}\|D_s\|^{1/2}\mathrm{exp}\{i𝐩_\alpha (𝐪𝐪^{})/\mathrm{}`$ $`(𝐪𝐪_\alpha )^2/2\sigma ^2(𝐪^{}𝐪_\alpha )^2/2\sigma ^2`$ $`+iS_j(𝐪,𝐪^{};E)/\mathrm{}i\nu _j^{}\pi /2\}d𝐪d𝐪^{}`$ To evaluate the integrals over $`𝐪`$ and $`𝐪^{}`$, the action is quadratically expanded about the points $`𝐪_f`$ and $`𝐪_i`$, $`S_j(𝐪,𝐪^{};E)`$ $`=`$ $`S_j(𝐪_f,𝐪_i;E)+𝐩_f(𝐪𝐪_f)𝐩_i(𝐪^{}𝐪_i)`$ (A5) $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,k}{\overset{N}{}}}[\left({\displaystyle \frac{p_f^{(i)}}{q^{(k)}}}\right)_{𝐪_f}(q^{(i)}q_f^{(i)})(q^{(k)}q_f^{(k)})`$ (A6) $``$ $`\left({\displaystyle \frac{p_i^{(i)}}{q^{(k)}}}\right)_{𝐪_i}(q^{(i)}q_i^{(i)})(q^{(k)}q_i^{(k)})`$ (A7) $`+`$ $`2\left({\displaystyle \frac{p_f^{(i)}}{q^{(k)}}}\right)_{𝐪_i}(q^{(i)}q_f^{(i)})(q^{(k)}q_i^{(k)})]`$ (A8) It is useful to define the vector $$𝐳=(z_1,\mathrm{},z_N,z_1^{},\mathrm{},z_N^{})$$ (A9) where $`z_i`$ $`=`$ $`(q^{(i)}q_f^{(i)})/\sigma `$ (A10) $`z_i^{}`$ $`=`$ $`(q^{(i)}q_i^{(i)})/\sigma `$ (A11) Thus, the integrals become $`S_{\alpha ,osc}(E)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{I}m{\displaystyle \frac{1}{i\mathrm{}(2\pi i\mathrm{})^{(d1)/2}}}\left({\displaystyle \frac{1}{\pi \sigma ^2}}\right)^{d/2}`$ (A14) $`\times {\displaystyle }{\displaystyle \underset{j}{}}\|D_s\|^{1/2}\sigma ^{2d}`$ $`\times \mathrm{exp}\{𝐳𝐀𝐳𝐛𝐳+c\}d𝐳`$ where $`𝐀`$ is composed of four $`N`$-dimensional matrices $$𝐀=\left(\begin{array}{cc}𝐀^{11}& 𝐀^{12}\\ 𝐀^{21}& 𝐀^{22}\end{array}\right)$$ (A15) and $`𝐛`$ $`=`$ $`(i\delta p_f^{(1)}\delta q_f^{(1)},\mathrm{},i\delta p_f^{(N)}\delta q_f^{(N)},`$ (A17) $`i\delta p_i^{(1)}\delta q_i^{(1)},\mathrm{},i\delta p_i^{(N)}\delta q_i^{(N)})`$ with $`\delta p_f^{(i)}`$ $`=`$ $`(p_\alpha ^{(i)}p_f^{(i)})\sigma /\mathrm{}`$ (A18) $`\delta p_i^{(i)}`$ $`=`$ $`(p_\alpha ^{(i)}p_i^{(i)})\sigma /\mathrm{}`$ (A19) $`\delta q_f^{(i)}`$ $`=`$ $`(q_\alpha ^{(i)}q_f^{(i)})/\sigma `$ (A20) $`\delta q_i^{(i)}`$ $`=`$ $`(q_\alpha ^{(i)}q_i^{(i)})/\sigma `$ (A21) and $`c`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}S_j(𝐪_f,𝐪_i;E){\displaystyle \frac{i}{\mathrm{}}}𝐩_\alpha (𝐪_f𝐪_i)`$ (A23) $`{\displaystyle \frac{(𝐪_f𝐪_\alpha )^2}{2\sigma ^2}}{\displaystyle \frac{(𝐪_i𝐪_\alpha )^2}{2\sigma ^2}}{\displaystyle \frac{i\nu _j^{}\pi }{2}}`$ The matrix $`𝐀`$ can be expressed in terms of the stability matrix, $`𝐌`$, where $`𝐌`$ has the same form as Eq. (A15) $$\left(\begin{array}{c}𝐩\\ 𝐪\end{array}\right)=𝐌\left(\begin{array}{c}𝐩^{}\\ 𝐪^{}\end{array}\right)$$ (A24) Thus, $`𝐀_{ab}^{11}`$ $`=`$ $`{\displaystyle \frac{\delta _{a,b}}{2}}{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}\left({\displaystyle \frac{p_f^{(a)}}{q^{(b)}}}\right)_{𝐪_f}`$ (A25) $`=`$ $`{\displaystyle \frac{\delta _{a,b}}{2}}{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}{\displaystyle \frac{\underset{i}{\overset{N}{}}m_{a,i}\mathrm{c}of\left(𝐌_{b,i}^{21}\right)}{det𝐌^{21}}}`$ (A26) $`=`$ $`{\displaystyle \frac{𝐈}{2}}{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}𝐌^{11}(𝐌^{21})^1`$ (A27) $`𝐀_{ab}^{12}`$ $`=`$ $`{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}\left({\displaystyle \frac{p_f^{(a)}}{q^{(b)}}}\right)_{𝐪_i}`$ (A28) $`=`$ $`{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}{\displaystyle \frac{\mathrm{c}of\left(𝐌_{b,a}^{21}\right)}{det𝐌^{21}}}`$ (A29) $`=`$ $`{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}((𝐌^{21})^1)^T`$ (A30) $`𝐀_{ab}^{21}`$ $`=`$ $`{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}\left({\displaystyle \frac{p_i^{(a)}}{q^{(b)}}}\right)_{𝐪_f}`$ (A31) $`=`$ $`{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}{\displaystyle \frac{\mathrm{c}of\left(𝐌_{a,b}^{21}\right)}{det𝐌^{21}}}`$ (A32) $`=`$ $`{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}(𝐌^{21})^1`$ (A33) $`𝐀_{ab}^{22}`$ $`=`$ $`{\displaystyle \frac{\delta _{a,b}}{2}}+{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}\left({\displaystyle \frac{p_i^{(a)}}{q^{(b)}}}\right)_{𝐪_i}`$ (A34) $`=`$ $`{\displaystyle \frac{\delta _{a,b}}{2}}{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}{\displaystyle \frac{\underset{i}{\overset{N}{}}m_{i+N,b+N}\mathrm{c}of\left(𝐌_{i,a}^{21}\right)}{det𝐌^{21}}}`$ (A35) $`=`$ $`{\displaystyle \frac{𝐈}{2}}{\displaystyle \frac{i\sigma ^2}{2\mathrm{}}}(𝐌^{21})^1𝐌^{22}`$ (A36) where $`m_{i,k}`$ are elements of the stability matrix and $`\mathrm{c}of\left(𝐌_{ik}^{21}\right)`$ is the signed minor of $`𝐌_{ik}^{21}`$. $`D_s`$ is a determinate involving second derivatives of the actions, $`D_s=\left\|\begin{array}{cc}\frac{^2S}{𝐪𝐪^{}}& \frac{^2S}{𝐪E}\\ \frac{^2S}{E𝐪^{}}& \frac{^2S}{E^2}\end{array}\right\|`$ $`=`$ $`{\displaystyle \frac{1}{\|\dot{q}^{(N)}\|\|\dot{q}^{(N)}\|}}\left\|{\displaystyle \frac{^2S}{\stackrel{~}{𝐪}\stackrel{~}{𝐪}^{}}}\right\|`$ (A39) $`=`$ $`{\displaystyle \frac{1}{\|\dot{q}^{(N)}\|\|\dot{q}^{(N)}\|}}\left({\displaystyle \frac{2\mathrm{}}{i\sigma ^2}}\right)^{(d1)}`$ (A41) $`\times \left\|\stackrel{~}{𝐀}^{21}\right\|`$ The tildes in the determinants in the above equation are used to exclude the $`N`$th coordinate. To obtain the above result $`q^{(N)}`$ and $`q^{(N)}`$ are chosen to be locally oriented along the trajectory where the dots indicate time derivatives. Since $`𝐀`$ is a symmetric matrix, the result for a general Gaussian integral is used and, hence, the strength function becomes $`S_{\alpha ,osc}(E)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{I}m{\displaystyle \frac{1}{i\mathrm{}(2\pi i\mathrm{})^{(d1)/2}}}\left({\displaystyle \frac{\sigma ^2}{\pi }}\right)^{d/2}\left({\displaystyle \frac{2\mathrm{}}{i\sigma ^2}}\right)^{(d1)/2}`$ (A44) $`{\displaystyle \underset{j}{}}\left({\displaystyle \frac{\pi ^{2d}}{det𝐀}}\right)^{1/2}\left({\displaystyle \frac{1}{\|\dot{q}^{(N)}\|\|\dot{q}^{(N)}\|}}\right)^{1/2}\left\|\stackrel{~}{𝐀}^{21}\right\|^{1/2}`$ $`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{4}}𝐛𝐀^1𝐛+c\right\}`$ $`=`$ $`{\displaystyle \frac{\sigma }{\pi ^{1/2}\mathrm{}}}\mathrm{R}e{\displaystyle \underset{j}{}}\left({\displaystyle \frac{det\stackrel{~}{𝐀}^{21}}{det𝐀}}\right)^{1/2}\left({\displaystyle \frac{1}{\|\dot{q}^{(N)}\|\|\dot{q}^{(N)}\|}}\right)^{1/2}`$ (A46) $`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{4}}𝐛𝐀^1𝐛+c\right\}`$ where the time derivatives are evaluated at the saddle points. ## B Sum Rule for the Strength Function The determinant of the $`(2N\times 2N)`$ matrix $`𝐀`$, $$det𝐀=det\left(\begin{array}{cc}\frac{𝐈}{2}\frac{i\sigma ^2}{2\mathrm{}}𝐌^{11}(𝐌^{21})^1& \frac{i\sigma ^2}{2\mathrm{}}((𝐌^{21})^1)^T\\ \frac{i\sigma ^2}{2\mathrm{}}(𝐌^{21})^1& \frac{𝐈}{2}\frac{i\sigma ^2}{2\mathrm{}}(𝐌^{21})^1𝐌^{22}\end{array}\right)$$ (B1) can be reduced to determinants of $`(N\times N)`$ matrices by using the relation $`((𝐌^{21})^1)^T=𝐌^{12}+𝐌^{11}(𝐌^{21})^1𝐌^{22}`$ and some row and column manipulations , so that $`det𝐀`$ $`=`$ $`det({\displaystyle \frac{i\sigma ^2}{4\mathrm{}}}[𝐌^{11}+𝐌^{22}`$ (B3) $`+i({\displaystyle \frac{\mathrm{}}{\sigma ^2}}𝐌^{21}{\displaystyle \frac{\sigma ^2}{\mathrm{}}}𝐌^{12})])/det(𝐌^{21})`$ The coordinates parallel to the trajectory do not mix with the transverse coordinates, since a point on an orbit will remain on that particular orbit. Thus, the $`N`$th rows and columns of the individual matrices in the above expression are zero except for the $`(N,N)`$ elements. It is convenient to re-express the submatrices of the stability matrix in terms of the Lyapunov exponents. Let $`\{\lambda _i\}`$ be the set of Lyapunov exponents whose real part is positive ordered such that $`\lambda _1>\lambda _2>\mathrm{}>\lambda _{N1}`$. The Lyapunov exponents along the parallel coordinate are zero and we will only work with the reduced $`(2(N1)\times 2(N1))`$ stability matrix in what follows. Let $`𝚲`$ be the diagonal matrix of the eigenvalues of the reduced stability matrix, $$𝚲=\left(\begin{array}{cccccc}e^{\lambda _1t}& \mathrm{}& 0& 0& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& e^{\lambda _{N1}t}& 0& \mathrm{}& 0\\ 0& \mathrm{}& 0& e^{\lambda _1t}& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& 0& 0& \mathrm{}& e^{\lambda _{N1}t}\end{array}\right)$$ (B4) Thus, by a similarity transform the reduced stability matrix be can written in terms of the Lyapunov exponents, i.e. $$𝐌=𝐋𝚲𝐋^1$$ (B5) Hence, each of the elements of the stability matrix can be written as $$m_{ij}=\underset{k}{\overset{N}{}}a_{ij}^{(k)}e^{\lambda _kt}+b_{ij}^{(k)}e^{\lambda _kt}$$ (B6) where $`a_{ij}^{(k)}`$ and $`b_{ij}^{(k)}`$ are linear combinations of the elements of the $`𝐋`$ and $`𝐋^1`$ matrices. Because in general chaotic systems $`\lambda 0`$, the $`b_{ij}^{(k)}`$’s may be omitted without seriously effecting the above sum. All the determinants including the numerator and denominator of $`det𝐀`$ as well as $`det\stackrel{~}{𝐀}^{21}`$, thus, will involve products of Eq. (B6). The homoclinic orbits in the sum begin and end at the intersections of the stable and unstable manifolds near the Gaussian centroid. Since neither manifolds may cross themselves, then in the vicinity of the Gaussian centroid the branches of each manifold are nearly parallel to themselves. Thus, to an excellent approximation, the same similarity transformation will diagonalize the stability matrix for each individual orbit, regardless of the period. Consequently, the elements of $`𝐋`$ and $`𝐋^1`$ are period independent. Connections can be made between the determinants and the Kolmogorov-Sinai entropy. The Kolmogorov-Sinai entropy, $`h_{KS}`$, of a system can be expressed using Pesin’s Theorem as the sum of the Lyapunov exponents with positive real part, $$h_{KS}=\underset{i}{\overset{N1}{}}\lambda _i$$ (B7) If there is no mixing between the different coordinates, then the individual matrices $`𝐌^{11}`$, $`𝐌^{12}`$, $`𝐌^{21}`$ and $`𝐌^{22}`$ are diagonal. Thus, each matrix element depends only upon one Lyapunov exponent and the determinants are proportional to $`\mathrm{exp}(h_{KS}t)`$. This is the case for two dimensional systems where the parallel and perpendicular coordinates in the stability matrix separate as mentioned above. Hence, we have $$\left\|\frac{det\stackrel{~}{𝐀}^{21}}{det𝐀}\right\|\mathrm{exp}(h_{KS}t)$$ (B8) Unlike the periodic orbits, the homoclinic sum is over segments of the orbits. The number of homoclinic points will proliferate exponentially at the same rate as the fixed points in the neighborhood which is proportional to $`\mathrm{exp}(h_TT)`$ where $`h_T`$ is the topological entropy. This is demonstrated by examining the partitioning of the phase space mentioned in Sect. IV B which has exponential growth. The partitioning reflects the symbolic dynamics of the system. The symbolic code uniquely describes each orbit so that amount of code (partitions) cannot grow faster than the number of periodic points, since each code (partition) cannot represent more than one periodic point. Finally, the sum rule is obtained by setting the topological entropy and Kolmogorov-Sinai entropy equal to each other. Then, for the special case of no mixing in the stability matrix as mentioned above the combination of the amplitudes and the number of orbits yields $$\underset{j}{}\left\|\frac{det\stackrel{~}{𝐀}^{21}}{det𝐀}\right\|\mathrm{}𝑑T\mathrm{}$$ (B9)
warning/0006/cond-mat0006198.html	ar5iv	text	# 1 Introduction ## 1 Introduction Fluorescence resonance energy transfer (FRET) is a powerful technique to study many aspects of structure and dynamics of polymers and biopolymers.<sup>1-8</sup> In this technique, the polymer is doped with a donor and an acceptor at suitable positions along the chain. The donor is initially excited optically by a laser light which subsequently transfers its energy to the acceptor which is located at a distance R from the donor. In many applications, the distance between the donor and the acceptor is fixed, as in the case of rigid biopolymers.<sup>4,5,7,8</sup> However, in many cases of interest, this distance is a fluctuating quantity.<sup>1-3,6</sup> For example, the distance between the two ends of a polymer in solution executes a Brownian motion. If the polymer is assumed to be an ideal Gaussian chain, then the mean square distance between the two ends is $`Nb^2`$, where $`N`$ is the number of the monomer units and $`b`$ is monomer (or the Kuhn) length. Thus, as the number of monomers in the chain increases, the average distance between the two ends also increases. The probability distribution of $`R`$, $`P\left(R\right)`$, however, remains peaked at $`R=0`$ for ideal chain, although its height decreases as $`1/\sqrt{(}N)`$. All these properties play important parts in determining the observed dynamics of excitation energy transfer. The usually assumed mechanism for excitation energy transfer between a donor and an acceptor is the Forster mechanism<sup>6</sup> which gives the following expression for the singlet-singlet resonance energy transfer rate, $`k\left(R\right)`$, $$k\left(R\right)=k_F\frac{1}{1+\left(R/R_F\right)^6}$$ (1) where $`R_F`$ is the Forster radius and $`k_F`$ is the rate of excitation transfer when the separation between the donor and the acceptor is very small, that is $`R/R_F0`$. The Forster radius is usually obtained from the overlap of the donor fluorescence with the acceptor absorption and several other available parameters<sup>5</sup>. The dynamics of Forster energy migration has been traditionally investigated by performing time domain measurements of the decay of the fluorescence (due to excitation transfer) from the donor.<sup>1,4-6</sup> More recently, this technique has been used in single molecule spectroscopy of biopolymers<sup>7,8</sup>. In the latter, distance dependence of FRET provides relevant information about the conformation and dynamics of single biopolymers. Recently, FRET from single protein molecules has also been used to study protein folding<sup>9</sup>. At any given time after the initial excitation, the fluorescence intensity is a measure of the ”unreacted” donor concentration. As both $`k_F`$ and $`R_F`$ are determined by the donor- acceptor pair, the rate of decay of the fluorescence intensity provides a direct probe of the conformational dynamics of the polymer. When the polymer is in solution, each monomer (or polymer bead) executes Brownian motion. Because of the connectivity of the polymer chain, this Brownian motion of the monomers are strongly correlated. This many body nature of polymer dynamics can be described by a joint, time dependent probability distribution $`P(𝐫^N,t)`$ where $`𝐫^N`$ denotes the positions of all the N polymer beads. The time dependence of the probability distribution $`P(𝐫^N,t)`$ can be described by the following reaction-diffusion equation<sup>10,11</sup> $$\frac{}{t}P(𝐫^N,t)=_B(𝐫^N,t)P(𝐫^N,t)k\left(R\right)P(𝐫^N,t)$$ (2) where $`_B`$ is the full $`3N`$ dimensional diffusion operator given by, $$_B(𝐫,t)=D\underset{j=1}{\overset{N}{}}\frac{}{r_j}P_{eq}(𝐫,t)\frac{}{r_j}\frac{P(𝐫,t)}{P_{eq}(𝐫,t)}$$ (3) $`R`$ is the scalar distance between the two ends of the polymer chain and $`D`$ is the center of mass diffusion coefficient. The solution of Eq.2 with the sink term given by the Forster expression is highly non-trivial. In two seminal papers, Wilemski and Fixman presented a nearly analytic solution of the problem for any arbitrary sink.<sup>10,11</sup> The WF theory has been tested ($`\mathrm{𝑜𝑛𝑙𝑦}`$ for the average rate) by computer simulations when the sink is a Heaviside function.<sup>11-13</sup> We are not aware of any such simulation study with the Forster rate. Such a study is clearly important. The objective of this paper is mainly two fold. First, to present results of detailed Brownian dynamics (BD) simulations of Eq.2, with k(R) given by the Forster rate (Eq.1). Detailed investigation into the time dependence of the survival probability at a time t after the initial excitation ($`S_P\left(t\right)`$) is presented. We believe that this is the first such calculation of $`S_P\left(t\right)`$ for this kind of a problem. We find that the time dependence of $`S_P\left(t\right)`$ can be non-exponential for a large range of the relevant parameter space (N,$`k_F,R_F`$). Second, we present a detailed comparison of the simulated rate with the WF theory. This comparison has been carried out at the level of time dependent survival probability, again, we believe, for the first time. The organization of the rest of the paper is as follows. In next section, we present the details of the BD simulation. In section III, we discuss the WF theory. In section IV we present the simulation results and the comparison with the WF theory. Section V concluded with a discussion. ## 2 Simulation Details Brownian Dynamics (BD) simulations of polymer motion have been carried out with an idealized Rouse chain, in which the set of beads are connected by the harmonic potential, $$\beta U=\frac{3}{2b^2}\underset{j=1}{\overset{N1}{}}\left(r_jr_{j+1}\right)^2$$ (4) where the position vector of a bead $`j`$ is denoted by $`r_j`$, and $`N`$ is the number of beads constituting the polymer chain. The mean squared bond length is $`b^2`$. Equilibrium end-to-end distance of the polymer chain is, $$\left(r_Nr_0\right)^2=L^2=Nb^2$$ (5) In Rouse model, there exists no excluded volume forces and the hydrodynamic interactions between the monomer beads are ignored<sup>13</sup>. More details on this model can be found elsewhere<sup>13</sup>. In our present study the polymer chain is additionally characterized by the presence of two reactive end groups. This essentially implies that within the time interval $`\mathrm{\Delta }t`$, the two end groups react with a probability $`k\left(R\right)\mathrm{\Delta }t`$. The initial configuration for each trajectory has been selected from a Monte Carlo generated random distribution. The following equation of motion has been used in the simulations, $$r_j\left(t+\mathrm{\Delta }t\right)=r_j\left(t\right)+F_j\left(t\right)\mathrm{\Delta }t+\mathrm{\Delta }X^G\left(t\right)$$ (6) where the positions of the j-th bead at time $`t`$ and $`t+\mathrm{\Delta }t`$ are denoted by $`r_j\left(t\right)`$ and $`r_j\left(t+\mathrm{\Delta }t\right)`$, respectively. $`F_j\left(t\right)`$ is the total force acting on bead $`j`$ and $`\mathrm{\Delta }X^G\left(t\right)`$ is a random Brownian displacement, taken from a Gaussian distribution with zero mean and the variance $`\left(X^G\right)^2=2\mathrm{\Delta }t`$. In writing Eq.1, we have set both $`D_0`$ and $`K_BT`$ equal to unity; the latter is Boltzmann constant times the temperature. Here we adopted a time step, $`\mathrm{\Delta }t=0.001`$. However, smaller time steps as low as 0.0001 has been employed in the limit of both large $`k_F`$ and large $`R_F`$ values, to account for the faster dynamics. During the simulation both the mass $`m`$ and the root mean-square bond length $`b^2`$ of the bead have been set to unity for computational convenience. The trajectory generated by using the above procedure needs to be terminated when the two end groups react. This has been done in simulations in the following way. Each time the trajectory is updated, the existing end to end distance $`R`$ is used in Eq.1 to calculate the distance dependent rate constant, $`k\left(R\right)`$. We then call a random number generator to get a value between zero and unity. If this value is less than $`k\left(R\right)\mathrm{\Delta }t`$, then the trajectory is terminated. Otherwise, the trajectory is continued. One forms a histogram over many such trajectories. This procedure generates an irreversible FRET. For each polymer chain constructed randomly, we have equilibrated it for 10,000 time steps before switching on the reaction. Subsequently, 50,000 to 1 Lakh trajectories with different initial configurations were generated and the survival probability $`S_p\left(t\right)`$ was obtained by averaging over all the trajectories. This procedure was systematically applied for the polymer chains containing the beads, $`N=20,50`$ and $`100`$. Before proceeding with the simulations of the Forster transfer, we reproduced the results of Pastor,Zwanzig and Szabo (PZS)<sup>12</sup> on the mean first passage time with Heaviside sink function of infinite strength. Our simulation results agreed with those of PZS within the uncertainity given by PZS. ## 3 Wilemski-Fixman Theory <sup>10,11</sup> Several decades ago Wilemski and Fixman (WF) developed a non-trivial theory for the diffusion limited intrachain reaction of a flexible polymer.<sup>10,11</sup> To account for the chemical reaction they have added a sink term $`𝒮`$ to the manybody diffusion equation, just as in Eq.2. The WF equation of motion is well-known and we present it below for the sake of completion $$\frac{}{t}P(𝐫^N,t)+_BP(𝐫^N,t)=k_0𝒮\left(R\right)P(𝐫^N,t).$$ (7) In the notation of the present work $$k_0=k_F;S\left(R\right)=\frac{1}{1+\left(R/R_F\right)^6}.$$ (8) The operator $`_B(𝐫^N,t)`$ is given by Eq.3. As already mentioned, the treatment of WF is general and can be applied to a reaction with arbitrary sink. Let us define a survival probability $`S_p\left(t\right)`$ as the probability that the chain has not reacted after time $`t`$. $`S_p\left(t\right)`$ is then given by, $$S_p\left(t\right)=P(𝐫^N,t)𝑑𝐫_1𝑑𝐫_2\mathrm{}𝑑𝐫_N$$ (9) In order to obtain the survival probability WF made a closer approximation, according to which the Laplace transform of $`S_P\left(t\right)`$ can be written as, $$\widehat{S}_p\left(s\right)=\frac{1}{s}\frac{k\upsilon _{eq}}{s^2\left(1+k\widehat{D}\left(s\right)/\upsilon _{eq}\right)}$$ (10) where $`\widehat{D}\left(s\right)`$ is defined as, $$\widehat{D}\left(s\right)=_0^{\mathrm{}}e^{st}D\left(t\right)𝑑t$$ (11) which is the Laplace transform of sink-sink time correlation function $`D\left(t\right)`$ defined as, $$D\left(t\right)=d^3𝐑_1d^3𝐑_2𝒮\left(𝐑_1\right)𝒮\left(𝐑_2\right)G(𝐑_1,𝐑_2,t)$$ (12) The Green function appearing in the above equation is given by, $`G(𝐑_1,𝐑_2,t)=\left({\displaystyle \frac{3}{2\pi L^2}}\right)^{3/2}({\displaystyle \frac{1}{\left(1\rho ^2\right)^{3/2}}})exp({\displaystyle \frac{R_1^22\rho \left(t\right)𝐑_1.𝐑_2+R_{}^{2}{}_{2}{}^{}}{2L^2\left(1\rho ^2\right)}})`$ (13) where $`\rho \left(t\right)`$ is the normalized time correlation function of end-to-end vector $`𝐑\left(0\right).𝐑\left(t\right)/R^2`$ which can be obtained analytically and is given by the following equation, $$\rho \left(t\right)=\frac{8}{\pi ^2}\underset{l;odd}{}\frac{4}{l^2}exp\left(\lambda _lt\right)$$ (14) If we neglect excluded volume and hydrodynamic interactions, $`\lambda _l`$ is given by the following expression<sup>10,11</sup>. $$\lambda _l=3D_0\left(l\pi /Nb\right)^2,$$ (15) Finally $`\upsilon _{eq}`$ is given as, $$\underset{t\mathrm{}}{lim}D\left(t\right)=\left(\upsilon _{eq}\right)^2$$ (16) Note that $`\upsilon _{eq}`$ is the rate when the distribution of the polymer ends is at equilibrium. Thus, $`\upsilon _{eq}`$ gives the initial rate of decay of $`S_P\left(t\right)`$ which will show up as the transient behavior. In most cases, the rate of decay should become progressively slower, as the population from the sink region decreases as the reaction proceeds. Once the choice of sink function specified, it is straight forward to calculate the survival probability by utilizing the above set of equations. WF choice was the heaviside sink function. Later, Doi showed that WF method is easy to apply if the heaviside sink function is replaced with a Gaussian sink function<sup>14</sup>. Bettizzeti and Perico studied the dependence of the rate on the choice of sink function with in the frame work WF theory and supported the WF closure approximation<sup>15</sup>. Surprisingly, no analysis of the time dependence of the survival probability, $`S_P\left(t\right)`$ has ever been reported. In this study we follow the orignal scheme proposed by WF to obtain $`S_p\left(t\right)`$ analytically. In doing so we use the Stefest algorithm to obtain $`S_p\left(t\right)`$ through the Laplace inversion of Eq. 10. ## 4 Results and Discussion Before discussing the results let us describe the scaling that has been used to compare the results obtained by simulation with the theory. In the reduced unit notation adopted in simulation, the rate constant has been scaled as $`\stackrel{~}{k}_F=k_Fb^2/D_0`$ and the real time has been scaled by $`b^2/D_0`$. However, in the original WF theory, time is scaled by $`6D/L^2`$ where $`D`$ is the center of mass diffusion. In the free draining limit, so $`D=D_0/N`$ and $`L^2=Nb^2`$. Thus, the numbers obtained from WF theory is to be converted to the simulation scaling for a comparison of results. The Forster radius is scaled by the bead diameter, $`b`$. Another important parameter in this problem is the root mean square radius of the polymer as this determines the end to end distribution. Although we have carried out simulations for N=20, 50 and 100, in this report we shall concentrate mostly on N=50. Figures 1 and 2 depict the time dependence of the survival probability $`S_p\left(t\right)`$ for two different values of $`R_F`$, $`R_F=1`$ and $`R_F=5`$, respectively, for a fixed $`N=50`$. In figure 1, $`k_F`$ has been varied from 50 to 0.1 , that is, over two orders of magnitude. The decay remains non-exponential over the whole range. In figure 2, $`k_F`$ has been varied from 10 to 0.1. Here the decay is exponential-like. These two figures demonstrate the strong dependence of the decay profile of $`S_P\left(t\right)`$ on $`R_F`$. Note that the earlier experiments which fitted the quantum yield to the Forster expression obtained values which are rather large, comparable to the ones shown in figure 2. This could have been due to the use of an equilibrium end-to-end distribution in the fitting, instead of a time dependent probability distribution. In model calculations, one usually assumes a small value of $`R_F`$ ( often in the form of a Heaviside sink function). This strong dependence of decay profile on $`R_F`$ could be potentially useful in unravelling mechanism and dynamics of energy transfer. It is not difficult to understand the above results qualitatively. For an ideal Gaussian chain, the maximum in the probability distribution that the two ends are separated by a distance R is located at $`\sqrt{(}2N/3)b`$. For N=50, this value is 5.773b. Therefore, when $`R_F`$ is equal to 5, the decay is facilitated by the presence of a large fraction of the distribution at a distance of separation where the transfer rate is large. This can explain the exponential-like decay for $`R_F=5`$ (figure 2). However, the situation is completely different for $`R_F=1`$. Here the probability of finding a polymer with end to end distance so small is negligible and the transfer rate where the bulk of the population is located is very small because of the strong $`R`$ dependence of the Forster transfer rate. Therefore, the decay of the survival probability starts slowly (Fig.1) and is determined by the interplay between the diffusion and the rate. This explains the shape of figure 1. The above discussion also suggests that the shape of the survival probability can depend rather strongly on the length of the polymer chain. This is because the Forster distance for a given donor-acceptor pair is likely to be independent of the length of the polymer chain. But the distribution and also the diffusion rate will be determined by N. However, this dependence is not trivial and will be discussed elsewhere. In figures 3 and 4, we have compared predictions of the WF theory with the simulations. In figure 3, $`S_P\left(t\right)`$ is plotted for two very different values of $`k_F`$ ($`K_F=`$ 1 and 10) at $`R_F=5`$ for $`N=50`$. It is seen that while the agreement is satisfactory at short times for both the cases, the same is not true at long times, particularly for the smaller $`k_F`$. In the latter case the simulation also finds a larger non-exponentiality, as discussed later. In figure 4, $`S_P\left(t\right)`$ is plotted for for $`R_F=1`$ , $`k_F=`$1.0 (and $`N=50`$). It is seen that the WF theory breaks down in this limit. This is one of the main results of the present study. The agreement improves for smaller $`K_F`$ but becomes worse in the opposite limit. We have also simulated both larger and smaller chains. Since the size of the polymer scales with N, it is not possible to compare results for different sizes. In figure 5, we show the comparison between the simulation results and the WF theory for $`N=100`$ at $`R_F=8`$ and $`k_F=0.1`$. This is the most favorable parameter space for the WF theory. Although the simulated $`S_P\left(t\right)`$ decays somewhat faster, the WF theory prediction is not totally off. We have not yet searched for any scaling laws (in N, $`R_F`$ and L) dependence – work is under progress in this direction. In figures 6 and 7, we present logarithmic plots of $`S_P\left(t\right)`$ to show the extent of non-exponentiality, for different values of $`R_F`$ and $`k_F`$. It can be seen that while the decay is nearly exponential for small $`K_F`$ (which is expected), it is strongly non-exponential when the value of $`k_F`$ becomes comparable to or larger than the bead diffusion rate, $`D_0/b^2`$. The inability of the Wilemski-Fixman theory to explain the time dependence of the survival probability is surprising. We note that earlier theoretical studies have considered only the mean first passage time. In figure 8 we have compared the simulated end-to-end vector time correlation function ($`\rho \left(t\right)`$) with the slightly approximate expression used by WF. The agreement is good, as expected. This agreement improves further for larger N. Thus, the failure of WF (as shown in Fig.4) must be due to the closure approximation. ## 5 Conclusions Use of FRET in single molecule spectroscopy of polymers and biopolymers requires accurate knowledge of the mechanism of energy transfer, more importantly, the distance dependence of the transfer rate. The fluorescence quantum yield can provide only an average estimate of the distance between the donor and the acceptor if the mechanism is well-understood. This could be sufficient for rigid systems. For many systems of interest, for example for understanding the dynamics of protein folding or in the collapse of polymers, one requires the time dependence of the excitation migration. This will be measured in terms of the time dependent survival probability. In this work we have presented results of detailed Brownian dynamics simulations of Forster energy transfer between the two ends of an ideal Gaussian chain. As noted by previous workers<sup>12</sup>, this apparently simple problem is actually highly non-trivial because this is a manybody problem. We have calculated survival probability for a large number of values of the transfer rate $`k_F`$ and the Forster radius, $`R_F`$. It is found that while the survival probability is exponential-like for small values of $`k_F`$ (compared to the monomer diffusion rate, $`D_0/b^2`$) and intermediate $`R_F`$, it is strongly non-exponential for small (compared to L) $`R_F`$. We have compared the results of the simulation with the well-known theory of Wilemski and Fixman. It is found that the theory is reliable when the Forster radius $`R_F`$ is comparable to the root mean square radius L of the polymer chain and the transfer rate $`k_F`$ is comparable to or smaller than the monomer diffusion rate $`D_0/b^2`$. However, the agrement is not at all satisfactory in the limit when $`R_F`$ is much smaller or larger than L. What is the reason for the failure of the WF theory when $`R_F`$ is substantially different from the root mean square radius of the polymer? While it is obvious that the WF closure approximation is inadequate in many situations, the exact reason for the failure is not clear. In fact, for the nature of the decay curve for large or small $`R_F`$ can possibly be understood even from a one dimensional theory, provided the end to end distance correlation function $`\rho \left(t\right)`$ is given. This problem will be discussed elsewhere. Note that the distance dependent rate appears in several other chemical processes, like in electron transfer reactions where the rate of transfer is known to show an exponential distance dependence. It will be interesting to study this problem with the method employed here. Another important, long standing problem is the study of reactions in realistic polymer chains with excluded volume and hydrodynamic interactions. Work in this direction is under progress. Acknowledgement This work is supported in parts by the Council of Scientific and Industrial Research (CSIR) and the Department of Science and Technology, India. The visit of A. Yethiraj was supported by the Jawaharlal Nehru Center for Advanced Scientific Research, Bangalore, India. G. Srinivas thanks CSIR, New Delhi, India for a research fellowship. Figure captions: Figure 1. The survival probability obtained from Brownian dynamics (BD) simulations of Eq.2 is plotted against the scaled time for several values of $`k_F`$ at $`R_F=`$ 1. The curves from top to bottom represent the cases with $`k_F=`$ 0.1, 1, 10 and 50, respectively. Figure 2. The survival probability $`S_p\left(t\right)`$, obtained from BD simulations is plotted for $`k_F=`$ 0.1, 1 and 10 at $`R_F=`$ 5. Curves from top to bottom show $`S_p\left(t\right)`$ at $`k_F=`$ 0.1, 1 and 10, respectively. Figure 3. BD simulation results have been compared with WF theory at a Forster radius, $`R_F=`$ 5. The upper set shows the case with $`k_F=`$ 1 and the lower set is for $`k_F=10`$. In both the cases, symbols shows the simulation results while the WF theory predictions are represented by the full lines. Figure 4. WF theory has been compared with the simulation results at a lower value of $`R_F`$, namely $`R_F=1`$ and for $`k_F`$=1. WF theory prediction has been shown by the full line while the symbols represent the simulation results. As seen from the figure, WF theory seems to break down in this limit. Figure 5. The comparison between WF theory and simulation results has been shown for a larger polymer chain, $`N=100`$, for $`R_F=8`$ and $`k_F`$ $`=0.1`$. Symbols and the full line represents the results of simulation and WF theory, respectively. Figure 6. The semilog plot of the survival probability $`S_p\left(t\right)`$ which has been obtained from simulations, is plotted against the scaled time at $`k_F=`$ 0.1 and $`R_F=`$ 5. This figure shows that the decay is nearly exponential for $`R_F=`$ 5 and small $`k_F`$. Figure 7. The semilog plot of the survival probability, obtained from simulations is plotted against the scaled time at $`k_F=`$ 1.0 and $`R_F=`$ 1. Highly non-exponential behavior of $`S_p\left(t\right)`$, is very clear in this limit. Figure 8. The end-to-end vector time correlation function $`\rho \left(t\right)`$ is plotted against the scaled time for a polymer of mean square length $`L^2=50b^2`$. Symbols shows the simulated $`\rho \left(t\right)`$ while the $`\rho \left(t\right)`$ obtained from WF expression is shown by full line.
warning/0006/astro-ph0006062.html	ar5iv	text	# Comparison of the ENEAR Peculiar Velocities with the PSCz Gravity Field ## 1 Introduction In the standard picture for the formation of cosmic structures via gravitational instability the peculiar velocity of a galaxy is generated by fluctuations in the mass distribution. For galaxies outside virialized systems, linear perturbation theory predicts $$𝐯(𝐫)\frac{\mathrm{\Omega }^{0.6}H_o}{4\pi }d^3r^{}\delta _m\frac{(𝐫^{}𝐫)}{\|𝐫^{}𝐫\|^3},$$ (1) where $`\mathrm{\Omega }`$ is the mass density parameter, $`H_o`$ is the Hubble constant and $`\delta _m`$ is the mass density fluctuation field. If the relationship between the galaxy distribution, $`\delta _g`$, and $`\delta _m`$ is approximately linear, $`\delta _g=b\delta _m`$, then the parameter $`\beta =\mathrm{\Omega }^{0.6}/b`$ can be derived from the comparison between the observed peculiar velocity field and that predicted from the galaxy distribution. A particularly useful method for performing a velocity-velocity comparison is the modal expansion method developed by Nusser & Davis (1995, hereafter ND95). This method expands the velocity fields by means of smooth functions defined in redshift space, thus alleviating the Malmquist biases inherent in real space analysis. Furthermore, the modal expansion filters the observed and predicted velocities in the same way, so that the smoothed fields can be compared directly. Because the number of modes is substantially smaller than the number of data points, the method also provides the means of estimating $`\beta `$ from a likelihood analysis carried out on a mode-by-mode basis, instead of galaxy-by-galaxy. The similar smoothing and the mode-by-mode comparison substantially simplify the error analysis. The modal expansion method has previously been used in comparisons between the 1.2 Jy IRAS predicted velocities and observed velocities inferred from Tully-Fisher (TF) measurements (Davis, Nusser & Willick 1996, hereafter DNW, da Costa et. al. 1998). In this paper, we perform a similar analysis using the recently completed redshift-distance survey of early-type galaxies (hereafter ENEAR, da Costa et al. 2000) and the IRAS PSCz redshift survey (Saunders et. al. 2000). Because of differences in the nature of the data sets considered some slight changes in the method are required and are described below. Our goal is to investigate how well the velocity field mapped by early-type galaxies matches the velocity field inferred from the PSCz survey, and to obtain the parameter $`\beta `$ yielding the best match. In section 2, we briefly describe the ENEAR redshift-distance catalog. In section 3, we describe the modal expansion method as used here, present maps of the ENEAR and PSCz radial peculiar velocity field and perform a likelihood analysis to derive $`\beta `$. A brief summary of our conclusions is presented in section 5. ## 2 Data We use a sub-sample extracted from the all-sky ENEAR redshift-distance survey (da Costa et al. 2000) comprising 578 objects within $`cz6000`$ km s<sup>-1</sup>, 355 field galaxies and 223 groups/clusters. Galaxies have been objectively assigned to groups and clusters using redshifts taken from complete redshift surveys sampling the same volume. Individual galaxy distances were estimated from an inverse $`D_n\sigma `$ template relation derived by combining cluster data (e.g., Bernardi et al. 2000). The cluster sample consists of 569 galaxies in 28 clusters. Over 80% of the galaxies in the magnitude-limited sample and roughly 60% of the cluster galaxies have new spectroscopic and photometric data obtained by the ENEAR survey. Multiple observations using different telescope/instrument configurations ensure the homogeneity of the data. Furthermore, the sample completeness is uniform across the sky. ## 3 The modal expansion An unbiased estimate of $`\beta =\mathrm{\Omega }^{0.6}/b`$ can be obtained from the comparison between smooth velocity fields with similar spatial resolution, derived from the ENEAR and PSCz data. To generate smooth fields we expand the peculiar velocities of both data in terms of smooth base functions. The expansion carried out here shares the general properties of that used by ND95, but differs in details. In their application to TF catalogs, ND95 defined $`P_i5\mathrm{log}(1u_i/s_i)`$, where $`s_i=cz_i`$ is the galaxy redshift in km s<sup>-1</sup> and $`u_i`$ its radial peculiar velocity. The function $`P`$ was then expressed by an expansion involving smooth functions. The final estimate of the smoothed velocity field was that obtained by minimizing the scatter of the rotational speeds given the magnitudes in the inverse TF relation. The scatter was also simultaneously minimized with respect to the the parameter of the TF relation. This led to an unbiased calibration of the inverse TF relation because the sample was mainly magnitude selected. The galaxy angular size and velocity dispersion in the $`D_n\sigma `$ relation do not uniquely fix the magnitude according to which the ENEAR sample is selected. So simultaneous minimization might lead to a biased estimate of the parameters of the $`D_n\sigma `$ relation. Although the bias is mild we use the calibration of the inverse $`D_n\sigma `$ given by Bernardi et al. (2000) by a regression of $`\sigma `$ on $`D_n`$ in clusters. We also express the peculiar velocity, $`u`$, rather than the function $`P`$ in terms of smooth functions. Another difference is that ND95 used TF catalogs with all galaxies having the same relative distance error which allowed an additional simplification in the application of the modal expansion method, namely, the expansion in terms of orthogonal smooth functions. This made the TF velocity error covariance matrices diagonal. In the ENEAR sample, the relative distance error is not the same for all objects (galaxies and groups/clusters), so using orthogonal functions does not offer any further simplification since the ENEAR error matrix remains non-diagonal. The lack of orthogonality slightly complicates the error analysis but does not affect the efficiency of the expansion. Choosing the spherical harmonics and Bessel functions to be our base smooth functions we write the radial peculiar velocity model as $$\stackrel{~}{u}(s,\theta ,\varphi )=\underset{l,m,n}{}\alpha _{nlm}\left[j_l^{}\left(k_ny(s)\right)c_{l1}\right]Y_{lm}(\theta ,\varphi ).$$ (2) where the sum is over $`m=l`$ to $`+l`$, $`l=0`$ to $`l_{max}`$, and $`n=0`$ to $`n_{max}`$. For the reasons given in DNW, we formulate our model to describe the velocity field with respect to the motion of the Local Group. The constant $`c_{l1}`$ is non-zero only for the dipole term ensuring that $`u=0`$ at the origin. The function $`y(s)`$ in the argument of the Bessel functions makes their oscillations match the radial distribution of the ENEAR data. Here we take $`y^2=\mathrm{ln}[1+(\frac{s}{1000})^2]`$, but other similar forms can be used as well. Also the expansion does not include a Hubble-like flow ($`us`$) so we assume that any such flow has been consistently removed from the ENEAR and PSCz velocities. The coefficients $`\alpha _{nlm}`$ are found by minimizing $$\chi ^2=\sigma _i^2[\stackrel{~}{u}_iu_i^\mathrm{o}]^2$$ (3) where $`u_i^o`$ are the raw observed velocities and $`\sigma _i`$ is the error of the velocity estimate resulting from observational uncertainties and intrinsic scatter in the $`D_n\sigma `$. For field galaxies $`\sigma _i=0.23s_i`$ and for groups of galaxies it is reduced by $`1/\sqrt{N_g}`$, where $`N_g`$ is the number of galaxies in the group. ## 4 Smooth Velocity Maps and the determination of $`\beta `$ We apply the modal expansion method to smooth the raw measured velocities of the $`578`$ ENEAR objects within a redshift of 6000 km s<sup>-1</sup> (Bernardi et. al. 2000). We use 51 modes corresponding to $`l_{max}=4,n_{max}=3`$ in (2). The smoothing scale of these functions is linear with redshift and matches the low resolution filter used in da Costa et al. (1998) (see their Figure 1). The smoothed velocities were then derived by minimizing (3) with respect to $`\alpha _{lmn}`$ assuming an error of $`\sigma _i=0.23s_i/\sqrt{N_g}`$ in the raw velocities of the ENEAR objects. The reduced $`\chi ^2`$ per d.o.f of the fit was 1.017, a satisfactory value in this type of analysis (see DNW, da Costa et al. 1998). Given an assumed value for $`\beta `$ we interpolate the PSCz predicted velocity field, computed by Branchini et. al. (2000), to the positions of the ENEAR galaxies. Branchini et. al. obtained the PSCz velocities from the PSCz galaxy distribution with a Top-Hat window of width equal to half the mean particle separation at a given redshift. The PSCz fields are then expanded in the same orthogonal set of basis functions as employed for the ENEAR velocities. The PSCz and ENEAR velocities are guaranteed to have the same resolution because the original smoothing of the PSCz density field is small compared to the resolution of the modal expansion. The smoothed ENEAR velocities are shown in Figure 1, in redshift shells 2000km s<sup>-1</sup> thick. Comparison of this figure and Figure 3 of da Costa et al. (1998) shows that the general flow pattern is remarkably similar. In the case of ENEAR, in the innermost shell very few prominent structures are probed by bright ellipticals. However, in the next two shells a strong dipole pattern can be easily recognized, being of comparable amplitude to that of observed with the SFI galaxies. This dipole corresponds to the reflex motion of the Local Group, with infalling galaxies in the Hydra-Centaurus direction and an outflow towards the Perseus-Pisces complex. The quality of the match can be evaluated from Figure 2 which shows the residual velocity field obtained subtracting the smoothed PSCz field from that of the ENEAR, assuming $`\beta =0.5`$. As can be seen the overall agreement is good with only a few more distant galaxies giving large residuals. Note, however, that even though with a larger amplitude, the mismatch seen in the outermost redshift shell at $`l0^{}`$, $`60^{}<b<15^{}`$ between ENEAR and PSCz correspond to mismatches in the comparison between SFI and 1.2 Jy IRAS velocity fields. This may correspond to a real mismatch between measured and predicted velocities which deserves further investigation. The filtered ENEAR and PSCz velocity fields are fully described by the modal expansion coefficients, $`\alpha _{en}`$ and $`\alpha _{ps}`$, of the ENEAR and PSCz fields, respectively. Since the number of these coefficients is significantly smaller than the number of galaxies, it is more efficient to estimate $`\beta `$ by comparing the modes rather than the individual galaxy velocities. As in da Costa et al. (1998) we define our best estimate for $`\beta `$ as the value that corresponds to the minimum of the pseudo-$`\chi ^2`$ $$\stackrel{~}{\chi }^2(\beta )=\underset{j,j^{}}{}\left[\alpha _{en}^j\alpha _{ps}^j(\beta )\right]\left[𝐓+𝐌(\beta )\right]^1\left[\alpha _{en}^j^{}\alpha _{ps}^j^{}(\beta )\right],$$ (4) where $`𝐓<\delta \alpha _{en}^j\delta \alpha _{en}^j^{}>`$ and $`𝐌<\delta \alpha _{ps}^j\delta \alpha _{ps}^j^{}>`$ are the the error covariance matrices of the coefficients $`\alpha _{en}^j`$ and $`\alpha _{ps}^j`$, respectively. For brevity of notation we have replaced the triplet $`n,l,m`$ with one index $`j`$. The PSCz covariance matrix $`𝐌`$ incorporates errors due to (i) the uncertainty in the LG motion, which creates a dipole discrepancy between the ENEAR and the PSCz velocities, $`(ii)`$ the discreteness in distribution of galaxies which propagates into the velocity field. and $`(iii)`$ small scale coherent (as in triple valued zones) nonlinear velocities that are not included in the PSCz recovered velocities. Details of how these error contributions are computed are in da Costa et al. (1998). Since the expansion functions are not orthogonal the ENEAR covariance matrix $`𝐓`$ has nonzero off-diagonal elements. This matrix is simply the inverse of $`\frac{^2\chi ^2}{\alpha _{en}^j\alpha _{en}^j^{}}`$ where the derivatives are computed at the minimum of $`\chi ^2`$ given by (3). Given the covariance matrices, we compute the curve of the reduced $`\stackrel{~}{\chi }^2(\beta )`$ as a function of $`\beta `$, which is shown in the top panel of Figure 3. The curve was computed with an error of 150 km s<sup>-1</sup> in the estimation of the LG motion and 160 km s<sup>-1</sup> for the amplitude of nonlinear error in the PSCz field (see da Costa et. al. 1998). This amplitude of the nonlinear error was chosen to make the $`\stackrel{~}{\chi }^2`$ per d.o.f equal to unity at the minimum. In their analysis of the SFI and 1.2 Jy IRAS, da Costa et al. (1998) obtained a lower value of 90 km s<sup>-1</sup> for the amplitude of this error. The difference can be attributed, as expected, to a better match between the SFI and IRAS velocities and the increased nonlinearities in the PSCz velocity at the positions of the ENEAR galaxies which preferentially reside in high density regions. The minimum value of the $`\stackrel{~}{\chi }^2`$ is attained at $`\beta =0.5`$, with the 1-sigma error being less than $`0.1`$. We note that this result is not sensitive to the exact values adopted for the error estimates Another statistic indicating the goodness of the match between the fields for various $`\beta `$ is the correlation function of the residual $`u_{en}u_{ps}`$ between the smoothed ENEAR and PSCz radial velocities. This is shown in the bottom panel of Figure 3 for $`\beta =0.2`$, 0.5, and 0.9. The amplitude of the PSCz field is small for $`\beta =0.2`$, so the correlation function for this $`\beta `$ is close to the correlation function of $`u_{en}`$ alone, while the opposite is true for $`\beta =0.9`$. On the other hand, for $`\beta =0.5`$ the correlation of the residual velocity field is significantly smaller, indicating a good match between the measured and predicted velocity fields. ## 5 Summary and Discussion Using the modal expansion method of ND95 and the recently completed ENEAR redshift-distance survey and the PSCz redshift survey we have carried out a comparison between the observed peculiar velocity field and that predicted from the distribution of PSCz galaxies. We find that the corresponding smoothed fields agree well and the best match is obtained with $`\beta =0.5\pm 0.1`$. This value is intermediate to those derived using the Mark III and SFI catalogs both based primarily on spiral galaxies. It is also consistent with the results obtained by Borgani et al. (2000) using an independent method based on modeling the velocity correlation function. Note, however, that the discrepancy between the values determined from these methods and those obtained from the power spectrum analysis (e.g., Zaroubi et al. 2000) and density-density comparisons (e.g., Sigad et al. 1998) still persist. The good agreement between SFI and 1.2 Jy IRAS and between ENEAR and PSCz implies that the SFI and ENEAR velocity fields are also in good agreement. This suggests that the velocity maps obtained from the new distance-redshift surveys are a fair representation of the underlying velocity field, as the general characteristics of the observed flow fields are independent of the type of galaxies and distance indicators used. The good agreement among the values of $`\beta `$ obtained using Mark III, SFI, ENEAR, 1.2 Jy and PSCz catalogs gives further support to low values of $`\beta `$ and point toward low-density cosmologies. ## Acknowledgements The authors would like to thank M. Maia, C. Rité and O. Chaves for their contribution over the years. The results of this paper are based on observations conducted at the European Southern Observatory (ESO) and the MDM Observatory. CNAW acknowledges support from NSF AST95-29028.
warning/0006/hep-th0006071.html	ar5iv	text	# Untitled Document hep-th/0006071 Noncommutative Tachyons And String Field Theory Edward Witten Dept. of Physics, Cal Tech, Pasadena, CA and CIT-USC Center For Theoretical Physics, USC, Los Angeles CA It has been shown recently that by turning on a large noncommutativity parameter, the description of tachyon condensation in string theory can be drastically simplified. We reconsider these issues from the standpoint of string field theory, showing that, from this point of view, the key fact is that in the limit of a large $`B`$-field, the string field algebra factors as the product of an algebra that acts on the string center of mass only and an algebra that acts on all other degrees of freedom carried by the string. June, 2000 Recently, solitons in scalar field theories on noncommutative spacetimes with very large noncommutativity parameter $`\theta `$ have been constructed in a strikingly simple way . This insight has been applied \[2,,3\] to string theory with strong noncommutativity to test the predictions about tachyon condensation and brane annihilation . The purpose of the present paper is to examine these issues in the context of open string field theory, as formulated in the language of noncommutative geometry . String field theory has been used to explore tachyon condensation \[6--12\], generally at $`\theta =0`$. As one might expect \[13,,14\], there is considerable simplification for large $`\theta `$, and so we will in this paper consider in string field theory language the solutions studied in \[1--3\]. We will carry out this discussion in flat $`𝐑^{26}`$ or $`𝐑^{10}`$. Factorization Of The Algebra In string field theory, the starting point is an associative algebra $`𝒜`$ that is built by multiplying string fields. The naive idea for constructing such an algebra is to start with the association of open string states with vertex operators $`𝒱`$. An open string vertex operator is of course inserted at the boundary of the open string, at some proper time $`\tau `$ along the boundary. Naively speaking, to multiply string states we would like to just multiply the corresponding vertex operators. The trouble is that the product $`𝒱𝒱^{}`$ of open string vertex operators at the same point on the boundary is not well-defined, because of the familiar short distance singularities of products of local quantum field operators. We do have an operator product expansion (OPE) $$𝒱(\tau )𝒱^{}(\tau ^{})\underset{k}{}c_k\|\tau \tau ^{}\|^{a_k}𝒱_k(\tau ^{})\mathrm{for}\tau \tau ^{}.$$ The coefficients $`c_k`$ in this expansion depend on whether $`\tau >\tau ^{}`$ or $`\tau <\tau ^{}`$; this is the origin of noncommutativity. There does not seem to be any elegant way to eliminate the $`\tau `$ dependence and extract an associative algebra $`𝒜`$ from the operator product expansion. In open string field theory, this is done by making rather special choices of local parameters for insertions of vertex operators; the construction is perhaps most naturally described in terms of gluing of open string states . At any rate, in this paper, we will only need properties of the algebra $`𝒜`$ that follow in a very general way from the properties of the operator product algebra. Technical details in the definition of $`𝒜`$ will not be important. The operator product expansion conserves the “ghost number” of the vertex operators, and hence $`𝒜`$ is graded by ghost number. The classical string field $`A`$ is a ghost number one element of $`𝒜`$. The worldsheet BRST operator $`Q`$ is a ghost number one derivation of the algebra (that is, $`Q(AA^{})=QAA^{}+(1)^AAQA^{}`$), and the equation of motion of bosonic open string field theory is $$QA+AA=0.$$ A similar construction can be made for open superstrings, but has been argued to have technical difficulties associated with meeting of picture-changing operators on the worldsheet . A modification has been proposed to circumvent this difficulty . The effect of this is to replace (1) by a nonpolynomial equation which, for our purposes in the present paper, can be treated in precisely the same way. Now, the OPE algebra of open string vertex operators has a subalgebra $`𝒜_0`$ in which one does not use the string center of mass coordinate. Thus, $`𝒜_0`$ contains vertex operators that depend in an arbitrary fashion on the ghosts $`b`$ and $`c`$ and the derivatives of the spacetime coordinates $`X^i`$, $`i=1,\mathrm{},26`$, all taken at zero spacetime momentum. $`𝒜_0`$ contains, for example, $`bc(X^1)^{22}^3X^2`$, but not $`bc(X^1)^{22}^3X^2e^{ipX}`$ with $`p0`$. Operators of zero momentum are closed under OPE’s, and so $`𝒜_0`$ is a subalgebra of $`𝒜`$. One may ask whether there is a complementary subalgebra $`𝒜_1`$ generated only by the $`e^{ipX}`$ and without $`X`$, $`^2X`$, etc. Normally, the answer to this question is “no,” since even if one merely makes a classical Taylor series expansion, the operator products of exponentials involve also the derivatives of $`X`$: $$e^{ipX}(\tau )e^{iqX}(\tau ^{})e^{i(p+q)X}(\tau ^{})+i(\tau \tau ^{})pXe^{i(p+q)X}(\tau ^{})+\mathrm{}.$$ However, there is a limit in which one actually can factorize $`𝒜`$ in terms of commuting subalgebras as $`𝒜=𝒜_0𝒜_1`$, where as suggested above $`𝒜_0`$ consists of vertex operators of zero momentum and $`𝒜_1`$ is generated by operators $`e^{ipX}`$. This is the limit in which the NS two-form field $`B`$ is constant and large \[13,,14\]. In fact, we assume that $`B`$ is of maximal rank in 10- or 26-dimensional Euclidean space. To see this factorization, we first recall the form of the worldsheet propagator in the presence of a $`B`$-field. We take the string world-sheet to consist of the upper half plane. The propagator between boundary points $`\tau ,\tau ^{}`$ on the real axis is then \[17,,18\] $$X^i(\tau )X^j(\tau ^{})=\alpha ^{}\left(\frac{1}{g+2\pi \alpha ^{}B}\right)_S^{ij}\mathrm{ln}(\tau \tau ^{})^2+i\pi \alpha ^{}\left(\frac{1}{g+2\pi \alpha ^{}B}\right)_A^{ij}ϵ(\tau \tau ^{}).$$ Here $`g`$ is the closed string metric and $`()_S`$, $`()_A`$ denote the symmetric and antisymmetric part of a matrix. Now we take the limit $`B\mathrm{}`$ with $`g,\alpha ^{}`$ fixed. To be definite, take $`B=tB_0`$ with $`t\mathrm{}`$. (This is the same limit as in , but parametrized differently.) If we set $$X^i=Y^i/\sqrt{t},$$ then the propagator becomes $$Y^i(\tau )Y^j(\tau ^{})=\frac{1}{t(2\pi )^2\alpha ^{}}\left(\theta ^2\right)^{ij}\mathrm{ln}(\tau \tau ^{})^2+\frac{i}{2}\theta ^{ij}ϵ(\tau \tau ^{}),$$ where $`\theta =1/B_0`$. Now for $`t\mathrm{}`$, the $`e^{iqY}`$ do generate a closed algebra as the $`\mathrm{ln}(\tau \tau ^{})^2`$ term does not contribute. (We will make this more explicit momentarily.) This is the center of mass algebra $`𝒜_1`$. What about $`𝒜_0`$, the algebra that doesn’t contain the center of mass momentum? As $`Y(\tau )Y(\tau ^{})t^1/(\tau \tau ^{})^2`$, we see that if we use $`\sqrt{t}^nY`$ as the algebra generators of $`𝒜_0`$, then the structure constants are independent of $`t`$. A typical element of $`𝒜_0`$ is then $$bc\sqrt{t}^{n_1}Y^{a_1}\sqrt{t}^{n_2}Y^{a_2}\sqrt{t}^{n_3}Y^{a_3},$$ with a factor of $`\sqrt{t}`$ for each $`^nY`$. So in this limit, we have two algebras $`𝒜_0`$ and $`𝒜_1`$. They commute and the full algebra is $`𝒜=𝒜_0𝒜_1`$. To verify that $`𝒜_0`$ and $`𝒜_1`$ commute, we simply have to observe that $$\sqrt{t}^nY(\tau )e^{iqY}(\tau ^{})\frac{1}{\sqrt{t}}(\tau \tau ^{})^ne^{iqY},$$ since the only term in the propagator that contributes is the logarithmic term, proportional to $`t^1`$. The right hand side vanishes for $`t\mathrm{}`$. Finally, to verify that $`𝒜_1`$ is closed under OPE’s, we note that when we expand $$e^{ipY}(\tau )e^{iqY}(\tau ^{})e^{\frac{1}{2}i\theta _{ij}p^iq^j}e^{i(p+q)Y}(\tau ^{})(1+i(\tau \tau ^{})pY+\mathrm{}),$$ the corrections proportional to $`Y`$ can be dropped since the factor of $`Y`$ is not accompanied by a factor of $`\sqrt{t}`$. The OPE algebra $`𝒜_1`$ thus reduces to $$e^{ipY}(\tau )e^{iqY}(\tau ^{})e^{\frac{1}{2}i\theta _{ij}p^iq^j}e^{i(p+q)Y}(\tau ^{})$$ which is the algebra of functions on noncommutative $`𝐑^{10}`$ or $`𝐑^{26}`$; as explained in , the usual complications of the OPE disappear, because the dimensions vanish and the right hand side has no dependence on $`\tau \tau ^{}`$. So far we have assumed that the $`B`$-field has maximal rank. If $`B`$ has less than maximal rank, we modify the factorization so that $`𝒜_1`$ is generated by $`e^{ipX}`$ where $`p`$ is tangent to the noncommutative directions and $`𝒜_0`$ is generated by all other operators. We still get a factorization $`𝒜=𝒜_0𝒜_1`$ in terms of commuting subalgebras. In this more general case, an operator in $`𝒜_0`$ may carry momentum, but in commutative directions only, while $`𝒜_1`$, on the other hand, will be the algebra of functions in the noncommutative directions. $`𝒜_1`$ is a down-to-earth algebra that can be described concretely in finite-dimensional terms, while $`𝒜_0`$ contains all of the mysterious stringy complications. By a similar scaling, the BRST operator Q acts on $`𝒜_0`$ and commutes with $`𝒜_1`$. So the string field equation $`0=QA+AA`$ makes sense for $`A𝒜_0`$. Tachyon Condensation For bosonic strings, there are at least two important solutions known with $`A𝒜_0`$. One of them is $`A=0`$ and describes the ordinary open string vacuum. The other, which will here be called $`A_0`$, was first explored numerically in and is now understood to describe tachyon condensation to “nothing,” that is, to a state with only closed strings. This means, in particular, that the corresponding $`A_0`$-$`A_0`$ open strings, with boundary conditions at both ends determined by the classical open string solution $`A_0`$, have no physical excitations. Now more generally, we could introduce $`2\times 2`$ Chan-Paton factors and start with two $`D25`$-branes. The solution $$A=\left(\begin{array}{cc}0& 0\\ 0& A_0\end{array}\right)$$ describes annihilation to a state with just one $`D25`$-brane. There are now several kinds of open strings: (1) The $`0`$-$`0`$ open strings describe physical excitations of the surviving $`D25`$-brane. (2) The $`0`$-$`A_0`$ and $`A_0`$-$`A_0`$ open strings are expected to have no physical excitations. Since the solution $`A_0`$ lies in $`𝒜_0`$, which (after suitably rescaling the coordinates) is completely independent of $`B`$, the solution $`A_0`$ is completely insensitive to the $`B`$-field. Now, let us specialize to the limit of large $`B`$ and invoke the idea of \[1\]. Let $`\rho 𝒜_1`$ be any projection operator, that is any element with $$\rho ^2=\rho .$$ Then as $`[Q,\rho ]=0`$, we see that $`A=A_0\rho `$ obeys $`0=QA+AA`$ if $`A_0`$ does. Equivalently, since (1) implies that $`(1\rho )^2=(1\rho )`$, we can solve the equation with $$A=A_0(1\rho ).$$ Now, following , represent $`𝒜_1`$ as the algebra of operators on a Hilbert space $``$. The endpoint of a string has a Chan-Paton label that takes values in $``$. Take $`\rho `$ to be the projector onto a finite-dimensional subspace of $``$, say a subspace $`V`$ of dimension $`n`$. Write $`=VW`$ where $`W`$ is the orthocomplement of $`V`$; so $`\rho \|_V=1`$ and $`\rho \|_W=0`$. In expanding around the solution $`A=A_0(1\rho )`$, we get several kinds of open strings: (1) The $`V`$-$`V`$ open strings, that is strings each of whose endpoints are labeled by vectors in $`V`$, are governed by the usual equations of open string theory except that the effective open string algebra for these strings is $`𝒜_0M_n`$, where $`M_n`$ is the algebra of $`n\times n`$ complex matrices acting on $`V`$. Hence the momentum of these strings is always tangent to the commutative directions in spacetime, and there are effective $`n\times n`$ Chan-Paton factors. These modes describe the physical excitations of $`n`$ $`D(252p)`$-branes, where $`2p`$ is the number of noncommutative directions. (2) The $`V`$-$`W`$ and $`W`$-$`W`$ open strings are governed by the same equations that describe the $`0`$-$`A_0`$ and $`A_0`$-$`A_0`$ open strings discussed above; so if the usual conjectures about tachyon condensation are true, then these open strings have no physical excitations. Thus, this solution describes annihilation of a $`D25`$-brane to a collection of $`n`$ parallel $`D(252p)`$-branes, for arbitrary $`p`$ and $`n`$. Type IIA At the very formal level of our discussion, we can consider tachyon condensation for unstable $`D9`$-branes of Type IIA in much the same way. In factorizing $`𝒜=𝒜_0𝒜_1`$, we include the superconformal ghosts and worldsheet fermions in $`𝒜_0`$; $`𝒜_1`$ is as in the case of the bosonic string the algebra of functions in the noncommuting directions of spacetime. There is, conjecturally, still a solution $`A_0`$ that describes tachyon condensation, and a more general solution $`A=A_0(1\rho )`$ which, if $`\rho `$ is the projection operator to an $`n`$-dimensional subspace, describes annihilation to a system of $`n`$ $`D(92p)`$-branes. As noted in \[2,,3\], a further subtlety arises because for open string excitations of a Type IIA $`D9`$-brane, there is a $`𝐙_2`$ symmetry that changes the sign of the tachyon field. Let $`A_0^{}`$ be the conjugate solution with opposite tachyon field. The solutions $`A_0`$ and $`A_0^{}`$ describe tachyon condensation to two different closed string vacua that differ by the sign of the tachyon field; in fact, they differ by one unit of the Ramond-Ramond zero-form $`G_0/2\pi `$ . If $`\rho _1,\rho _2`$ obey $`\rho _1^2=\rho _1`$, $`\rho _2^2=\rho _2`$, $`\rho _1\rho _2=\rho _2\rho _1=0`$, we can make a more general solution with $`A=A_0\rho _1+A_0^{}\rho _2`$. A special case of this with $`\rho _1=\rho `$, $`\rho _2=1\rho _1`$ is $$A=A_0(1\rho )+A_1\rho .$$ Let us consider $`\rho `$ to be the projector onto all quantum states that are supported within a large region $`\mathrm{\Omega }`$ in the noncommutative phase space. This is only a rough, semiclassical description of $`\rho `$, but it should be good if $`\mathrm{\Omega }`$ is large enough. (1) describes tachyon condensation to a state in which the tachyon field has one sign outside of $`\mathrm{\Omega }`$ and another sign inside $`\mathrm{\Omega }`$; because the two states differ by one unit of $`G_0/2\pi `$, there is a supersymmetric $`D8`$-brane wrapped on the boundary of $`\mathrm{\Omega }`$. As noted in \[2,,3\], it is perplexing that in this approximation the tension of the $`D8`$-brane appears to be zero. To explore this puzzle in a little more detail (but without claiming to resolve it), consider the case of two noncommutative directions with coordinates $`x,y`$ obeying $`[x,y]=i\theta `$. Suppose that $`\mathrm{\Omega }`$ is a disc and that we want $`\rho `$ to be a projector onto an $`n`$-dimensional subspace. Then the area of $`\mathrm{\Omega }`$ should be $`A=2\pi \theta n`$, so its radius is $`r=\sqrt{2\theta n}`$. $`\rho `$ is approximately 1 deep inside $`\mathrm{\Omega }`$ and approximately 0 far from $`\mathrm{\Omega }`$; the scale of variation of $`\rho `$ is approximately the same as the width in space of the outermost quantum state onto which $`\rho `$ projects. (If we try to make $`\rho `$ vary more slowly than that, there will be states on which it is not equal approximately to either 0 or 1.) That outermost state fills a cylindrical shell near the boundary of $`\mathrm{\Omega }`$ of area $`2\pi \theta `$; the radial thickness of the shell is thus $`\mathrm{\Delta }r=\theta /r=\sqrt{\theta /2n}`$. The validity of our description rests on neglecting the logarithmic term in (1), which is proportional to $`\theta ^2/t\alpha ^{}`$; this term can be considered small if the scale of variation of the solution is large compared to $`\theta /\sqrt{t\alpha ^{}}`$. The condition we need is thus $`\mathrm{\Delta }r>>\theta /\sqrt{t\alpha ^{}}`$ or $$\frac{A}{t}<<\alpha ^{}.$$ Since $`A/t`$ is the area of $`\mathrm{\Omega }`$ in the original coordinates $`X`$, before the rescaling (1), this means that the solution (1) is actually only valid if the area of $`\mathrm{\Omega }`$ in string units is much less than one. Type IIB Now what can we say about tachyon condensation for Type IIB superstrings? For Type IIB, we could first of all consider a $`D9`$\- or $`D\overline{9}`$-brane with $`A=0`$. The corresponding boundary conditions for open strings we will call $`0`$ and $`\overline{0}`$, respectively. The combined $`D9`$-$`D\overline{9}`$ system is believed to admit a somewhat more interesting solution. First of all, to describe a $`D9`$-$`D\overline{9}`$ system, we use $`2\times 2`$ Chan-Paton matrices, but with the opposite GSO projection for the off-diagonal terms. Thus the string field takes the form $$A=\left(\begin{array}{cc}B& T\\ \overline{T}& B^{}\end{array}\right),$$ where $`B`$ and $`B^{}`$ have the usual GSO projections and $`T`$ and $`\overline{T}`$ have the opposite ones; thus $`B`$ and $`B^{}`$ describe gauge fields as well as stringy excitations, while $`T`$ and $`\overline{T}`$ have a tachyon at the lowest level. There is a symmetry $$Te^{i\theta }T,\overline{T}e^{i\theta }\overline{T}.$$ It is believed that there exists a solution that describes tachyon condensation to the closed string vacuum. It has been explored numerically \[9,,11,,12\] but is not known in closed form; we merely write it as $$A_0=\left(\begin{array}{cc}\alpha & \beta \\ \overline{\beta }& \gamma \end{array}\right)$$ Because of the symmetry (1), it can be generalized to $$A_\theta =\left(\begin{array}{cc}\alpha & e^{i\theta }\beta \\ e^{i\theta }\overline{\beta }& \gamma \end{array}\right).$$ As in the discussion of the bosonic string, we can add extra uncondensed $`D9`$’s and $`D\overline{9}`$’s and mix the two solutions. Thus, in a larger space, we can consider the string field $$A=\left(\begin{array}{cc}0& 0\\ 0& A_\theta \end{array}\right),$$ where the upper left block describes excitations of a $`D9`$ or $`D\overline{9}`$. This field certainly obeys the equations of string field theory, and describes partial annihilation of a system of $`D9`$’s and $`D\overline{9}`$’s, leaving a single brane. The open strings in expanding around this solution can be classified as $`0`$-$`0`$ open strings, which describe ordinary open string excitations, and $`0`$-$`\theta `$, $`\theta `$-$`0`$, and $`\theta `$-$`\theta `$ open strings, all of which describe no physical modes at all if the usual hypotheses about tachyon condensation are correct. In the large $`B`$ limit, the string field algebra factors as $`𝒜=𝒜_0𝒜_1`$ just as before. We want to generalize the solution $`A_\theta 𝒜_0`$ to include the $`𝒜_1`$ factor, by a suitable generalization of the previous ansatz $`A=A_0\rho `$. For this, we let $`\sigma `$ be an element of $`𝒜_1`$ and $`\overline{\sigma }`$ its complex conjugate (or hermitian adjoint in a Hilbert space representation). We take $$A=\left(\begin{array}{cc}\alpha \overline{\sigma }\sigma & \beta \overline{\sigma }\\ \overline{\beta }\sigma & \gamma \sigma \overline{\sigma }\end{array}\right).$$ To obey $`QA+AA=0`$ given that $`A_0`$ obeys this equation, <sup>1</sup> And similarly for any other equation, like the one in , that is constructed by multiplication of string fields as well as operations like $`Q`$ that commute with $`𝒜_1`$. the properties we need are $$\sigma \overline{\sigma }\sigma =\sigma ,\overline{\sigma }\sigma \overline{\sigma }=\overline{\sigma }.$$ If $`\sigma `$ is invertible, these equations say that $`\sigma `$ is unitary. On an eigenspace with $`\sigma =e^{i\theta }`$, we get $$A=A_\theta =\left(\begin{array}{cc}\alpha & e^{i\theta }\beta \\ e^{i\theta }\beta & \gamma \end{array}\right).$$ This solution describes tachyon condensation to a state with no physical excitations. The fun comes when $`\sigma `$ is not invertible. Let $`V`$ be the kernel of $`\sigma `$, and $`W`$ the kernel of $`\overline{\sigma }`$ (or equivalently the cokernel of $`\sigma `$). Let $`n`$ and $`m`$ be the dimensions of $`V`$ and $`W`$, and $`M_n`$, $`M_m`$ the algebras of matrices acting on $`V`$ and $`W`$ respectively; we suppose that $`n`$ and $`m`$ are finite. The equations (1) mean in general that $`\sigma `$ is a unitary isomorphism between the orthocomplement of $`V`$ and the orthocomplement of $`W`$. If $`2p`$ is the number of noncommutative directions, then string states whose endpoints are labeled by vectors in $`V`$ describe the excitations of $`n`$ $`D(92p)`$-branes, and those whose endpoints are labeled by vectors in $`W`$ describe the excitations of $`m`$ $`D(\overline{92p})`$\- branes. The $`V`$-$`W`$ open strings are equivalent to conventional $`D(92p)`$-$`D(\overline{92p})`$ open strings. Other open string excitations of this system are governed by the same equations as the $`0`$-$`\theta `$ and $`\theta `$-$`\theta `$ excitations of the the solution (1), and describe no physical excitations at all, if the conventional hypotheses are correct. This solution thus describes tachyon condensation down to a system with $`n`$ $`D(92p)`$-branes and $`m`$ $`D(\overline{92p})`$-branes. The net $`D(92p)`$-brane charge is $`nm`$, which is the same as the index of the operator $`\sigma `$.<sup>2</sup> This solution has been described in a more detailed setting in section 4 of . Note in eqn. (4.8) of that paper an operator of index 1. To describe explicitly a solution of (1) with nonzero index, suppose that there are two noncommutative directions with coordinates $`x,y`$, with $$[x,y]=i\theta ,\theta >0.$$ We introduce the creation and annihilation operators $$a=\frac{xiy}{\sqrt{2\theta }},\overline{a}=\frac{x+iy}{\sqrt{2\theta }},[a,\overline{a}]=1.$$ $`a`$ and $`\overline{a}`$ are represented on a Hilbert space $``$ that contains a vector $`\|0`$ with $`a\|0=0`$. The kernel of $`a`$ is generated by $`\|0`$, and $`\overline{a}`$ has no kernel. We let $$\begin{array}{cc}\hfill \sigma & =\frac{1}{\sqrt{\overline{a}a+1}}a\hfill \\ \hfill \overline{\sigma }& =\overline{a}\frac{1}{\sqrt{\overline{a}a+1}}.\hfill \end{array}$$ Clearly, the kernel of $`\sigma `$ is generated by $`\|0`$, and the kernel of $`\overline{\sigma }`$ is empty; so the index of $`\sigma `$ is 1. A short computation shows that $`\sigma \overline{\sigma }=1`$, which implies (1). From (1), we have $$\sigma =\frac{1}{\sqrt{x^2+y^2+\theta }}(xiy).$$ If $`\sigma `$ is treated as a classical function of $`x`$ and $`y`$, then for $`x,y\mathrm{}`$ we have $`\|\sigma \|=1`$. Ignoring the commutator $`[x,y]`$ is valid in describing the behavior near infinity. We can thus regard $`\sigma `$ near infinity as a $`U(1)`$-valued function on a circle; as such, its winding number is $`1`$. Thus, in this particular case, the index equals minus the winding number. This relation is a special case of the general Atiyah-Singer index theorem. Since the index and the winding number are both topological invariants, the relation between the index and the winding number can be proved, in this particular problem, by computing the index for one operator of every possible index, for example $`\sigma =(1/\sqrt{a^n\overline{a}^n})a^n`$ for positive index or $`\sigma =\overline{a}^n(1/\sqrt{a^n\overline{a}^n})`$ for negative index. We leave details to the reader. 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warning/0006/quant-ph0006018.html	ar5iv	text	# Can Quantum Computer Perform Better than Classical ? ## Abstract We present a theoretical model of a quantum device which can factorize any number $`N`$ in two steps, by preparing an input state and then performing a proper measurement process. However, the analysis reveals that the duration of state preparation and measurement is at least proportional to $`N`$ and hence the computation is not efficient. On the other hand, the energy consumption of this quantum computer grows like $`\mathrm{log}N`$ while for classical ones is exponential in the input bit size $`m=\mathrm{log}_2N`$. These results suggest the existence of a generalized Heisenberg relation which put limits on the efficiency of quantum computers in terms of the total computation time, the total energy consumption and the classical complexity of the problem. The action of a quantum computer is described on the abstract level in the following way. We have a quantum system with $`N`$ orthogonal states (computational basis) which can store $`n=\mathrm{log}_2N`$ bits of information. Firstly, we have to prepare our quantum system in an initial (input) state. Then a quantum algorithm is performed which is realized as a sequence of $`T`$ unitary transformation called quantum gates. Finally, the output state is measured. The number $`T`$ of involved quantum gates is assumed to be proportional to the physical time needed to achieve a given task. The algorithm is called efficient if $`T`$ is polynomial in $`m=\mathrm{log}_2N`$. The celebrated Shor’s quantum algorithm which factorizes numbers into primes is efficient while its classical counterparts are not. This remarkable result had an enormous impact on the development of the whole field of quantum information and quantum computing . There are, however, the following factors which can spoil the performance of a quantum computer: A) The unavoidable inaccuracies in manufacturing of a hardware which produce deviations of the real computer’s Hamiltonian from the designed one. B) The number of elementary transformations which form a quantum algorithm is by no means unique and does not determine the real physical time of computation. In particular, by increasing the energy level spacing of the corresponding computer’s Hamiltonian we can speed up its time evolution. On the other hand the decoherence effects due to an interaction with an environment typically grow in a nonlinear way with the energy level spacing. C) A finite duration of the input state preparation procedure and the output state measurement process should be taken into account. In particular, one can expect that the Heisenberg energy-time uncertainty relation may put here some universal limits. The factor A) poses a rather technological problem and will be not discussed here. The questions raised in B) will be discussed in details in the forthcomming publication . Our goal is to show by constructing an explicit but purely theoretical model that there exist fundamental obstacles related to C) which could destroy the efficiency of quantum algorithms. We describe the operation of a fictitious quantum machine which can factorize any number $`N`$ with $`T=0`$ quantum gates, it means by preparing an input state and a proper measurement of it only. Although formally, there are no quantum gates in this model the quantum algorithm exists in the form of a specially designed Hamiltonian of the system. The dynamics governed by the system’s Hamiltonian influences state preparation and measurements processes. In this case the problems raised in B) are neglected. Consider a resonant cavity which supports the states of a photon (radiation modes) with the frequencies being the logarithms of prime numbers times a fixed frequency unit $`\omega `$ $$\omega _q=\omega \mathrm{log}q,q=2,3,5,7,11,13,\mathrm{}$$ $`(1)`$ The second quantization Hamiltonian of the electromagnetic field in the cavity written in terms of annihilation and creation operators $`\widehat{a}_q,\widehat{a}_q^+`$ $$\widehat{H}=\mathrm{}\omega \underset{q}{}(\mathrm{log}q)\widehat{a}_q^+\widehat{a}_q$$ $`(2)`$ possesses nondegenerated eigenvalues being proportional to the logarithms of all natural numbers $$\widehat{H}\psi _N=E_N\psi _N,E_N=\mathrm{}\omega \mathrm{log}N,N=1,2,3,4,\mathrm{}$$ $`(3)`$ The structure of the corresponding eigenstates reveals the factorization of $`N`$ into prime numbers. Namely, $$\psi _N(\widehat{a}_{q_1}^+)^{m_1}(\widehat{a}_{q_2}^+)^{m_2}\mathrm{}(\widehat{a}_{q_r}^+)^{m_r}\psi _1$$ $`(4)`$ where $$N=(q_1)^{m_1}(q_2)^{m_2}\mathrm{}(q_r)^{m_r}$$ $`(5)`$ and $`\psi _1`$ is a vacuum state. In words, at the state $`\psi _N`$ we have $`m_1`$ photons of the frequency $`\omega \mathrm{log}q_1`$, $`m_2`$ photons of the frequency $`\omega \mathrm{log}q_2`$, …, and $`m_r`$ photons of the frequency $`\omega \mathrm{log}q_r`$. Therefore, in principle one can find the factorization of any number $`N`$ in two steps. First we prepare the system in the state $`\psi _N`$ of a given energy $`E_N=\mathrm{}\omega \mathrm{log}N`$ by transferring this portion of energy into the empty cavity. Then we open the cavity and perform a spectral analysis of the corresponding radiation field counting photons in different modes. Let us discuss the possible preparation process. We perturb our quantum system being initially in a vacuum state by a weak external interaction Hamiltonian $`\widehat{V}(t)`$ which is periodic in time with a tunable frequency $`\mathrm{\Omega }`$ $$\widehat{V}(t)=\widehat{W}\mathrm{cos}(\mathrm{\Omega }t).$$ $`(6)`$ Instead of a usual electromagnetic interaction linear in field we assume that $`\widehat{W}`$ is a sufficiently nonlinear function of electromagnetic field operators which allows multiphoton excitations such that the (virtual) transitions between the vacuum $`\psi _1`$ and any state $`\psi _N`$ are possible i.e. $$<\psi _1,\widehat{W}\psi _N>0,forallN=2,3,4,\mathrm{}$$ $`(7)`$ For a given number $`N`$ which we want to factorize we tune the frequency $`\mathrm{\Omega }`$ to the value $`\omega \mathrm{log}N`$. The time dependent first-order perturbation calculus gives us the probability of excitation of the state $`\psi _M`$ $$p_M(t)=\frac{2}{\mathrm{}^2}\|<\psi _1,\widehat{W}\psi _M>\|^2\frac{\mathrm{sin}^2\left\{\frac{1}{2}\omega \left(\mathrm{log}M\mathrm{log}N\right)t\right\}}{\omega ^2\left(\mathrm{log}M\mathrm{log}N\right)^2}.$$ $`(8)`$ As the energy level spacing around $`E_N=\mathrm{}\omega \mathrm{log}N`$ is $`\delta E_N\mathrm{}\omega /N`$ it follows from the formula (8) that we have to wait for a time at least of the order $$tN\omega ^1$$ $`(9)`$ to be sure that the selected state $`\psi _N`$ has been prepared with much larger probability than the other neighboring states $`\psi _M`$. The similar estimation can be easily obtained for the duration of the measurement process. Therefore the total computation time $`t_c`$ grows exponentially with $`\mathrm{log}N`$ like in the classical situation. It is obvious that the result obtained above can be treated as a special case of the Heisenberg time-energy uncertainty relation $$\mathrm{\Delta }t\mathrm{\Delta }E\frac{\mathrm{}}{2}.$$ $`(10)`$ Indeed in order to identify the energy of a quantum state with an accuracy $`\mathrm{}\omega /N`$ we need a time longer than $`N/\omega `$. One should notice, however, that our quantum computer is superior to classical ones in respect of energy consumption, at least for the case of existing irreversible computers (see for the theory of reversible computations). The energy used for the factorization of $`N`$ is equal to $`E_c=\mathrm{}\omega \mathrm{log}N`$ while for the classical irreversible computers any logical step consumes an energy portion and hence the total energy cost of factorization grows exponentially with the input bit size $`m=\mathrm{log}_2N`$. Taking into account the eq.(9) we obtain the following inequality independent of an arbitrary frequency scale $`\omega `$ $$t_cE_c>>\mathrm{}N\mathrm{log}N.$$ $`(11)`$ The form of the inequality (11) suggests the following general hypothesis. There exists an inequality which puts universal limits on the performance of a quantum computer in terms of the total computation time $`t_c`$, the total energy consumption $`E_c`$ and the complexity $`𝒞(\mathrm{log}_2N)`$ of the problem to be solved. This ”generalized Heisenberg relation” reads $$t_cE_c>>\mathrm{}𝒞(\mathrm{log}_2N).$$ $`(12)`$ The complexity $`𝒞(\mathrm{log}_2N)`$ is a function of the input bit size and is defined by a minimal number of logical steps needed to solve the problem. The inequality (12) means that for a non-efficient optimal classical algorithm the quantum computation is also not efficient either with respect to the computation time or the used energy. In order to prove this hypothesis we cannot restrict ourselves to counting quantum gates in the algorithm but we have to discuss physical implementation of all steps of quantum computing including state preparation and measurement processes. ###### Acknowledgements. The author thanks Michał , Paweł , and Ryszard Horodecki’s and S. Kryszewski for discussions.The work is supported by the Grant KBN PB/273/PO3/99/16.
warning/0006/hep-ph0006126.html	ar5iv	text	# 1 Introduction ## 1 Introduction This talk is given at Fradkin memorial conference on Friday, 9 June 2000. Friday was always a traditional day of seminars on quantum field theory in our department, so-called Fradkin seminars. Two talks on wavelets were given at this seminar last year just before Fradkin was put in a hospital. Therefore I decided to talk here about wavelet application in particle physics and gave to organizers of the conference the title containing the first part of the above one. However soon I learned that the topic on QCD (the second part of the present title) would fit the conference schedule better and asked organizers for replacement. It was shifted to QCD session but with an old title on wavelets. Thus I decided to give two talks at one, and it explains how these two topics appear here together. ## 2 Wavelets First I learned about wavelets from Pete Carruthers in 1993. He applied them for analysis of some scaling cascade models used, in particular, in multiparticle production modelling. It was briefly described in our review paper . Then I proposed to use wavelets for pattern recognition in high energy nucleus-nucleus collisions, and it was applied to experimental data . Wavelets became a powerful mathematical tool in many investigations. They are used in those cases when the result of the analysis of a particular signal<sup>3</sup><sup>3</sup>3The notion of a signal is used here for any recorded information about some processes, objects, functions etc. should contain not only the list of its typical frequencies (scales) but also the knowledge of the definite local coordinates where these properties are important. Wavelets form a complete orthonormalized system of functions with a finite support by using dilations and translations. That is why by changing a scale (dilations) they can describe the local characteristics of a signal, and by translations they cover the whole region in which it is studied. Due to the completeness of the system, they also allow for the inverse transformation to be done. The locality property of wavelets leads to their substantial superiority over Fourier transform which provides us only with the knowledge of global frequencies (scales) of the object under investigation because the system of functions used (sine, cosine) is defined on the infinite interval. High energy collisions of elementary particles result in production of many new particles in a single event. Each newly created particle is depicted kinematically by its momentum vector i.e. by a dot in the three-dimensional phase space. Different patterns formed by these dots in the phase space would correspond to different dynamics. To understand this dynamics is a main goal of all studies done at accelerators and in cosmic rays. Especially intriguing is a problem of the quark-gluon plasma, the state of matter with deconfined quarks and gluons which could exist during an extremely short intervals of time. One hopes to create it in collisions of high energy nuclei. Nowadays, the data about Pb-Pb collisions are available where in a single event more than 1000 particles are produced. We are waiting for RHIC accelerator in Brookhaven and LHC in CERN to provide events with up to 20000 new particles created. Therefore the problem of phase space pattern recognition in an event-by-event analysis becomes meaningful. It is believed that the detailed characterization of each collision event could reveal the rare new phenomena, and it will be statistically reliable due to a large number of particles produced in a single event. When individual events are imaged visually, the human eye has a tendency to observe different kinds of intricate patterns with dense clusters (spikes) and rarefied voids. However, the observed effects are often dominated by statistical fluctuations. The method of factorial moments was proposed to remove them but it is hard to use in event-by-event approach. The wavelet analysis avoids smooth polynomial trends typical for the statistical component. It was first applied to analyze the individual high multiplicity event of Pb-Pb interaction at energy 158 GeV per nucleon. With emulsion technique used in experiment the angles of particle emission are often measured only, and the two-dimensional phase space is considered therefore. The experimental statistics is rather low but acceptance is high and homogeneous that is important for proper pattern recognition. To simplify the analysis , the two-dimensional target diagram representing the polar and azimuthal angles of created charged particles was split into 24 one-dimensional functions representing the polar angle distribution of these particles in 24 azimuthal angle sectors of $`\pi /12`$ and in each of them particles were projected onto the polar angle $`\theta `$ axis. Thus one-dimensional functions of the rapidity distribution of these particles in 24 sectors were obtained. Then the wavelet coefficients were calculated in all of them and tied up together (continuous MHAT wavelet was used). The resulting pattern showed that many particles are concentrated close to some value of the polar angle i.e. reveal the ring-like structure in the target diagram. The interest to such patterns is related to the fact that they can result from the so-called gluon Cherenkov radiation or, more generally, from the gluon bremsstrahlung at a finite length within a quark-gluon medium (plasma, in particular). More elaborated two-dimensional analysis was done recently and confirmed these conclusions with jet regions tending to lie on some ring-like formations. The jet-like substructure of the event becomes more pronounced, and ring-like correlations of jettty regions are noticeable. With higher statistics, one can learn if the angular distribution of these rings corresponds to theoretical expectations. It is due to wavelet analysis that for the first time the fluctuation structure of an event is shown in a way similar to the target diagram on the two-dimensional plot. Previously, some attempts to consider such events with different methods of treating the traditional projection and correlation measures revealed just that such substructures lead to spikes in the angular (pseudorapidity) distribution and are somewhat jetty. Various Monte Carlo simulations of the process were compared to the data and failed to describe this jettiness in its full strength. More careful analysis of large statistics data on hadron-hadron interactions (unfortunately, however, for rather low multiplicity) with dense groups of particles separated showed some ”anomaly” in the angular distribution of these groups awaited from the theoretical side. Further analysis using the results of wavelet transform are needed to check this conclusion in high multiplicity nucleus-nucleus interactions when many events of this kind become available. ## 3 QCD Here I briefly describe recent advances in theoretical understanding of multiplicity distributions of quark and gluon jets. The extended survey with more detailed comparison to experimental data can be found in the recent review paper . The progress in experimental studies of properties of quark and gluon jets is very impressive. Therefore the study of the energy evolution of such parameters of multiplicity distributions of jets as their average multiplicities and widths becomes possible. It is well known that the average multiplicities of quark and gluon jets increase quite fast with energy but their ratio has a much slower dependence. The perturbative QCD provides quite definite predictions which can be confronted to experiment. In brief, the results can be summarized by saying that the energy dependence of the mean jet multiplicity can be perfectly fitted but the ratio of gluon to quark jet multiplicities can be described within the precision of 15-20$`\%`$ only. Moreover, one can understand why next-to-leading approximation is good enough for describing the energy dependence, but it is not quite satisfactory yet for the ratio value. I show this by presenting the analytical expressions. For the corresponding Figures, I refer the reader to the review paper . The theoretical asymptotical value of the ratio of average multiplicities equal 2.25 is much higher than its experimental values, which are in the range from 1.05 at comparatively low energies of $`\mathrm{{\rm Y}}`$ resonance to 1.5 at $`Z^0`$ resonance. The next-to leading order (NLO) corrections reduce this ratio from its asymptotical value by about 10$`\%`$ at $`Z^0`$ energy. The NNLO and 3NLO terms diminish it further and show the tendency to approximate the data with better accuracy. The computer solution of QCD equations for the generating functions has shown even better agreement with experiment not only on this ratio but on higher moments of multiplicity distributions as well. Being perfect at $`Z^0`$ energy, the agreement in the ratio is not as good at lower energies where the theoretical curve is still about 15-20$`\%`$ above the experimental one. In other words, the theoretically predicted slope of the ratio of multiplicities in gluon and quark jets is noticeably smaller than its experimental value. Nevertheless, one can speak about the steady convergence of theory and experiment with subsequent improvements being done. Moreover, it is even surprising that any agreement is achieved in view of the expansion parameter being extremely large (about 0.5) at present energies. The importance of studying the slopes stems from the fact that some of them are extremely sensitive (while others are not) to higher order perturbative corrections and to non-perturbative terms in the available energy region. Thus they provide us with a good chance to learn more about the structure of the perturbation series from experiment. In the perturbative QCD, the general approach to studying the multiplicity distributions is formulated in the framework of equations for generating functions. Therefrom, one can get equations for average multiplicities and, in general, for any moment of the multiplicity distributions . In particular, two equations for average multiplicities of gluon and quark jets are written as $`n_G(y)^{^{}}={\displaystyle }dx\gamma _0^2[K_G^G(x)(n_G(y+\mathrm{ln}x)+n_G(y+\mathrm{ln}(1x)n_G(y))`$ $`+n_fK_G^F(x)(n_F(y+\mathrm{ln}x)+n_F(y+\mathrm{ln}(1x)n_G(y))],`$ (1) $$n_F(y)^{^{}}=dx\gamma _0^2K_F^G(x)(n_G(y+\mathrm{ln}x)+n_F(y+\mathrm{ln}(1x)n_F(y)).$$ (2) Herefrom one can learn about the energy evolution of the ratio of multiplicities $`r`$ and of the QCD anomalous dimension $`\gamma `$ (the slope of the logarithm of average multiplicity in a gluon jet) defined as $$r=\frac{n_G}{n_F},\gamma =\frac{n_G^{^{}}}{n_G}=(\mathrm{ln}n_G)^{^{}}.$$ (3) Here, prime denotes the derivative over the evolution parameter $`y=\mathrm{ln}(p\mathrm{\Theta }/Q_0),p,\mathrm{\Theta }`$ are the momentum and the initial angular spread of the jet, related to the parton virtuality $`Q=p\mathrm{\Theta }/2`$, $`Q_0`$=const, $`K`$’s are the well known splitting functions, $`n_G`$ and $`n_F`$ are the average multiplicities in gluon and quark jets, $`n_G^{^{}}`$ is the slope of $`n_G`$, $`n_f`$ is the number of active flavours. The perturbative expansion of $`\gamma `$ and $`r`$ is written as $$\gamma =\gamma _0(1a_1\gamma _0a_2\gamma _0^2a_3\gamma _0^3)+O(\gamma _0^5),$$ (4) $$r=r_0(1r_1\gamma _0r_2\gamma _0^2r_3\gamma _0^3)+O(\gamma _0^4),$$ (5) where $`\gamma _0=\sqrt{2N_c\alpha _S/\pi },\alpha _S`$ is the strong coupling constant, $$\alpha _S=\frac{2\pi }{\beta _0y}\left[1\frac{\beta _1\mathrm{ln}(2y)}{\beta _0^2y}\right]+O(y^3),$$ (6) $`\beta _0=(11N_c2n_f)/3,\beta _1=(51N_c19n_f)/3,r_0=N_c/C_F,`$ and in QCD $`N_c=3`$ is the number of colours, $`C_F=4/3`$. The limits of integration in eqs. (1), (2) used to be chosen equal either to 0 and 1 or to $`e^y`$ and $`1e^y`$. This difference, being negligibly small at high energies $`y`$, is quite important at low energies. Moreover, it is of physics significance. With limits equal to $`e^y`$ and $`1e^y`$, the partonic cascade terminates at the perturbative level $`Q_0`$ as is seen from the arguments of multiplicities in the integrals. With limits equal to 0 and 1, one extends the cascade into the non-perturbative region with low virtualities $`Q_1xp\mathrm{\Theta }/2`$ and $`Q_2(1x)p\mathrm{\Theta }/2`$ less than $`Q_0/2`$. This region contributes terms of the order of $`e^y`$, power-suppressed in energy. It is not clear whether the equations and LPHD hypothesis are valid down to some $`Q_0`$ only or the non-perturbative region can be included as well. Nevertheless, the purely perturbative expansion (4), (5) with constant coefficients $`a_i,r_i`$ and energy-dependent $`\gamma _0`$ is at work just in the case of limits 0 and 1. The values of $`a_i,r_i`$ for different number of active flavors $`n_f`$ are tabulated in . At $`Z^0`$-energy the subsequent terms in (5) diminish the value of $`r`$ compared with its asymptotics $`r_0=2.25`$ approximately by $`10\%,\mathrm{\hspace{0.17em}13}\%,\mathrm{\hspace{0.17em}1}\%`$ for $`n_f=4`$ getting closer to experiment. However the theoretical value of $`r`$ still exceeds its experimental values by 15-20$`\%`$. The energy dependence of mean multiplicities can be obtained from the definition (3) by inserting there the value of $`\gamma `$ (4) and integrating over $`y`$. Keeping the terms as small as $`y^1`$ at large $`y`$ in the exponent, one gets the following expressions for energy dependence of multiplicities of gluon (G) and quark (F) jets $$n_G=Ky^{a_1C^2}\mathrm{exp}\left[2C\sqrt{y}+\delta _G(y)\right],$$ (7) with $`K`$ an overall normalization constant, $`C=\sqrt{4N_c/\beta _0}`$, and $$\delta _G(y)=\frac{C}{\sqrt{y}}\left[2a_2C^2+\frac{\beta _1}{\beta _0^2}[\mathrm{ln}(2y)+2]\right]+\frac{C^2}{y}\left[a_3C^2\frac{a_1\beta _1}{\beta _0^2}[\mathrm{ln}(2y)+1]\right];$$ (8) $$n_F=\frac{K}{r_0}y^{a_1C^2}\mathrm{exp}\left[2C\sqrt{y}+\delta _F(y)\right],$$ (9) with $$\delta _F(y)=\delta _G(y)+\frac{C}{\sqrt{y}}r_1+\frac{C^2}{y}(r_2+\frac{r_1^2}{2}).$$ (10) It happens that 2NLO and 3NLO terms (contributing to $`y^{1/2}`$ and $`y^1`$ terms in the exponent) are almost constant at present energies and do not change the energy dependence prescribed in NLO approximation. It explains why the energy dependence is well fitted by both NLO and 3NLO formulas while 2NLO correction to the ratio $`r`$ is large and important. The rather small difference in $`r`$ values results in quite noticeable disagreement of the slopes $`r^{}`$. Theoretical estimates can be shown to be quite predictive for the ratio of the slopes of multiplicities but it is much less reliable to use the perturbative estimates even at $`Z^0`$-energy for such quantities as the slope of $`r`$ or the ratio of slopes of logarithms of multiplicities (the logarithmic slopes). Much higher energies are needed to do that. Thus the values of $`r^{^{}}`$ and/or of the logarithmic slopes can be used to verify the structure of the perturbative expansion. I demonstrate it here on the example of the slope value. The slope $`r^{}`$ is extremely sensitive to higher order perturbative corrections. The role of higher order corrections is increased here compared with $`r`$ because each $`n`$th order term proportional to $`\gamma _0^n`$ gets an additional factor $`n`$ in front of it when differentiated, the main constant term disappears and the large ratio $`r_2/r_1`$ becomes crucial: $$r^{^{}}=Br_0r_1\gamma _0^3\left[1+\frac{2r_2\gamma _0}{r_1}+\left(\frac{3r_3}{r_1}+B_1\right)\gamma _0^2+O(\gamma _0^3)\right],$$ (11) where the relation $`\gamma _0^{^{}}B\gamma _0^3(1+B_1\gamma _0^2)`$ has been used with $`B=\beta _0/8N_c;B_1=\beta _1/4N_c\beta _0`$. The factor in front of the bracket is very small already at present energies: $`Br_0r_10.156`$ and $`\gamma _00.5`$. However, the numerical estimate of $`r^{^{}}`$ is still indefinite due to the expression inside the brackets. Let us note that each differentiation leads to a factor $`\alpha _S`$ or $`\gamma _0^2`$, i.e., to terms of higher order. For values of $`r_1`$, $`r_2`$, $`r_3`$ tabulated above ($`n_f=4`$) one estimates $`2r_2/r_14.9`$, $`(3r_3/r_1)+B_11.5`$). The simplest correction proportional to $`\gamma _0`$ is more than twice larger 1 at energies studied and the next one is about 0.4. Therefore the ever higher order terms should be calculated for the perturbative values of $`r^{}`$ to be trusted. The slope $`r^{}`$ is equal to 0 for a fixed coupling constant. The higher order terms are also important for the moments of the multiplicity distributions . The normalized second factorial moment $`F_2`$ defines the width of the multiplicity distribution. The asymptotical $`(\gamma _00)`$ values of $`F_2^G`$ and $`F_2^F`$ are different: $$F_{2,as}^G=\frac{4}{3},F_{2,as}^F=\frac{7}{4}.$$ (12) At $`Z^0`$ energy the widths of the distributions are smaller $$F_2^G1.12,F_2^F1.34.$$ (13) but still much larger than their experimental values 1.02 and 1.08, correspondingly. The rather large difference of the perturbative (13) and experimental values at $`Z^0`$ indicates that moments of the distributions should be sensitive to corrections. The conclusions about the third moments are similar. Nonetheless, the computer solution of QCD equation happened to be quite successful in fitting experimental data even for higher moments and their ratios $`H_q`$ introduced in . It shows that the role of conservation laws treated approximately in the analytical approach and accurately accounted in computer calculations becomes more important for higher moments. Thus it is shown that the analytical approach is quite successful in demonstrating that all features of QCD predictions about multiplicities of quark and gluon jets correspond to the general trends of experimental data. Some disagreement at the level of 15-20$`\%`$ is understandable due to incomplete account for the energy-momentum conservation in such an approach. Further accurate computer solutions are needed to check if these trends persist at the higher precision level.
warning/0006/hep-ph0006186.html	ar5iv	text	# Relativistic quark-antiquark potential and heavy quarkonium mass spectra ## Abstract: A general approach to accounting for retardation effects in the long-range (confining) part of the quark-antiquark potential is presented. The charmonium and bottomonium mass spectra are calculated with the systematic account of relativistic and retardation effects and the one-loop radiative corrections. A good fit to available experimental data on the mass spectra is obtained. conference: Heavy Quark Physics 5, Dubna, Russia, 6-8 April 2000 The relativistic properties of the quark-antiquark interaction potential play an important role in analysing different static and dynamical characteristics of heavy mesons. The Lorentz-structure of the confining quark-antiquark interaction is of particular interest. In the literature there is no consent on this item. For a long time the scalar confining kernel has been considered to be the most appropriate one . The main argument in favour of this choice is based on the nature of the heavy quark spin-orbit potential. The scalar potential gives a vanishing long-range magnetic interaction, which is in agreement with the flux tube picture of quark confinement of , and allows to get the fine structure for heavy quarkonia in accord with experimental data. However, the calculations of electroweak decay rates of heavy mesons with a scalar confining potential alone yield results which are in worse agreement with data than with a vector potential . The radiative $`M1`$-transitions in quarkonia such as e. g. $`J/\psi \eta _c\gamma `$ are the most sensitive to the Lorentz-structure of the confining potential. The relativistic corrections for these decays arising from vector and scalar potentials have different signs . In particular, as it has been shown in ref. , agreement with experiments for these decays can be achieved only for a specific mixture of vector and scalar potentials. In this context, it is worth noting, that the recent study of the $`q\overline{q}`$ interaction in the Wilson loop approach indicates that it cannot be considered as purely scalar. Moreover, the found structure of spin-independent relativistic corrections is not compatible with a scalar potential. A similar conclusion has been obtained in ref. on the basis of a Foldy-Wouthuysen reduction of the full Coulomb gauge Hamiltonian of QCD. There, the Lorentz-structure of the confinement has been found to be of vector nature. The scalar nature of spin splittings in heavy quarkonia in this approach is dynamically generated through the interaction with collective gluonic degrees of freedom. Thus we see that while the spin-dependent structure of $`q\overline{q}`$ interaction is well established now, the spin-independent part is still controversial in the literature. The uncertainty in the Lorentz-structure of the confining interaction complicates the account of retardation corrections since the relativistic reconstruction of the static confining potential is not unique. Here we present the generalized prescription of such reconstruction and discuss its implications for the heavy quarkonium mass spectra. In our preceding papers we have developed the relativistic quark model based on the quasipotential approach. A meson is described by the wave function of the bound quark-antiquark state, which satisfies the quasipotential equation of the Schrödinger type $`\left({\displaystyle \frac{b^2(M)}{2\mu _R}}{\displaystyle \frac{𝐩^2}{2\mu _R}}\right)\mathrm{\Psi }_M(𝐩)`$ (1) $`={\displaystyle \frac{d^3q}{(2\pi )^3}V(𝐩,𝐪;M)\mathrm{\Psi }_M(𝐪)},`$ (2) where the relativistic reduced mass is $$\mu _R=\frac{E_aE_b}{E_a+E_b}=\frac{M^4(m_a^2m_b^2)^2}{4M^3},$$ (3) and $`E_a`$, $`E_b`$ are given by $$E_a=\frac{M^2m_b^2+m_a^2}{2M},E_b=\frac{M^2m_a^2+m_b^2}{2M}.$$ (4) Here $`M=E_a+E_b`$ is the meson mass, $`m_{a,b}`$ are the masses of light and heavy quarks, and $`𝐩`$ is their relative momentum. In the centre of mass system the relative momentum squared on mass shell reads $$b^2(M)=\frac{[M^2(m_a+m_b)^2][M^2(m_am_b)^2]}{4M^2}.$$ (5) The kernel $`V(𝐩,𝐪;M)`$ in Eq. (1) is the quasipotential operator of the quark-antiquark interaction. It is constructed with the help of the off-mass-shell scattering amplitude, projected onto the positive energy states. Constructing the quasipotential of the quark-antiquark interaction we have assumed that the effective interaction is the sum of the usual one-gluon exchange term with the mixture of long-range vector and scalar linear confining potentials, where the vector confining potential contains the Pauli interaction. The quasipotential is then defined by $`V(𝐩,𝐪;M)=\overline{u}_a(p)\overline{u}_b(p)\{{\displaystyle \frac{4}{3}}\alpha _sD_{\mu \nu }(𝐤)\gamma _a^\mu \gamma _b^\nu `$ (6) $`+V_V(𝐤)\mathrm{\Gamma }_a^\mu \mathrm{\Gamma }_{b;\mu }+V_S(𝐤)\}u_a(q)u_b(q),`$ (7) where $`\alpha _S`$ is the QCD coupling constant, $`D_{\mu \nu }`$ is the gluon propagator in the Coulomb gauge and $`𝐤=𝐩𝐪`$; $`\gamma _\mu `$ and $`u(p)`$ are the Dirac matrices and spinors with $`ϵ(p)=\sqrt{p^2+m^2}`$. The effective long-range vector vertex is given by $$\mathrm{\Gamma }_\mu (𝐤)=\gamma _\mu +\frac{i\kappa }{2m}\sigma _{\mu \nu }k^\nu ,$$ (8) where $`\kappa `$ is the Pauli interaction constant characterizing the anomalous chromomagnetic moment of quarks. Vector and scalar confining potentials in the nonrelativistic limit reduce to $$V_V(r)=(1\epsilon )Ar+B,V_S(r)=\epsilon Ar,$$ (9) reproducing $$V_{\mathrm{conf}}(r)=V_S(r)+V_V(r)=Ar+B,$$ (10) where $`\epsilon `$ is the mixing coefficient. The retardation contribution to the one-gluon exchange part of the $`q\overline{q}`$ potential is well known. For the confining part of the $`q\overline{q}`$ potential the retardation contribution is much more indefinite. It is a consequence of our poor knowledge of the confining potential especially in what concerns its relativistic properties: the Lorentz structure (scalar, vector, etc.) and the dependence on the covariant variables such as $`k^2=k_0^2𝐤^2`$. Nevertheless we can perform some general considerations and then apply them to a particular case of the linearly rising potential. To this end we note that for any nonrelativistic potential $`V(𝐤^2)`$ the simplest relativistic generalization is to replace it by $`V(k_0^2𝐤^2)`$. In the case of the Lorentz-vector confining potential we can use the same approach as for the one-gluon exchange even with more general vertices containing the Pauli terms, since the mass-shell vector currents are conserved here as well. It is possible to introduce alongside with the “diagonal gauge” the so-called “instantaneous gauge” which is the generalization of the Coulomb gauge: $`V_V(k_0^2𝐤^2)\overline{u}_a(𝐩)\overline{u}_b(𝐩)\mathrm{\Gamma }_a^\mu \mathrm{\Gamma }_{b\mu }u_a(𝐪)u_b(𝐪)`$ (11) $`=\overline{u}_a(𝐩)\overline{u}_b(𝐩)\{V_V(𝐤^2)\mathrm{\Gamma }_a^0\mathrm{\Gamma }_b^0`$ (12) $`[V_V(𝐤^2)𝚪_a𝚪_b+V_V^{}(𝐤^2)(𝚪_a𝐤)`$ (13) $`\times (𝚪_b𝐤)]\}u_a(𝐪)u_b(𝐪),`$ (14) where $$V_V(k_0^2𝐤^2)V_V(𝐤^2)+k_0^2V_V^{}(𝐤^2)$$ and $$k_0^2=(ϵ_a(𝐩)ϵ_a(𝐪))(ϵ_b(𝐪)ϵ_b(𝐩))\frac{(𝐩^2𝐪^2)^2}{4m_am_b}$$ (15) with the correct Dirac limit in which the retardation contribution vanishes when one of the particles becomes infinitely heavy . For the case of the Lorentz-scalar potential we can make the same expansion in $`k_0^2`$, which yields $$V_S(k_0^2𝐤^2)V_S(𝐤^2)+k_0^2V_S^{}(𝐤^2).$$ (16) But in this case we have no reasons to fix $`k_0^2`$ in the only way (15). The other possibility is to take a half sum instead of a symmetrized product, namely to set (see e. g. ) $`k_0^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(ϵ_a(𝐩)ϵ_a(𝐪))^2+(ϵ_b(𝐪)ϵ_b(𝐩))^2\right]`$ (17) $``$ $`{\displaystyle \frac{1}{8}}(𝐩^2𝐪)^2\left({\displaystyle \frac{1}{m_a^2}}+{\displaystyle \frac{1}{m_b^2}}\right).`$ (18) The Dirac limit is not fulfilled by this choice, but this cannot serve as a decisive argument. Thus the most general expression for the energy transfer squared, which incorporates both possibilities (15) and (17) has the form $`k_0^2=\lambda (ϵ_a(𝐩)ϵ_a(𝐪))(ϵ_b(𝐪)ϵ_b(𝐩))+(1\lambda )`$ (19) $`\times {\displaystyle \frac{1}{2}}\left[(ϵ_a(𝐩)ϵ_a(𝐪))^2+(ϵ_b(𝐪)ϵ_b(𝐩))^2\right]`$ (20) $`\lambda {\displaystyle \frac{(𝐩^2𝐪^2)^2}{4m_am_b}}`$ (21) $`+(1\lambda ){\displaystyle \frac{1}{8}}(𝐩^2𝐪)^2\left({\displaystyle \frac{1}{m_a^2}}+{\displaystyle \frac{1}{m_b^2}}\right),`$ (22) where $`\lambda `$ is the mixing parameter. Thus the spin-independent part of $`q\overline{q}`$ potential with the account of retardation corrections takes the form: $`V_{\mathrm{SI}}(r)=V_C(r)+V_{\mathrm{conf}}(r)+V_{\mathrm{VD}}(r)`$ (23) $`+{\displaystyle \frac{1}{8}}\left({\displaystyle \frac{1}{m_a^2}}+{\displaystyle \frac{1}{m_b^2}}\right)\mathrm{\Delta }\left[V_C(r)+(1+2\kappa )V_V(r)\right],`$ (24) where the velocity-dependent part $`V_{\mathrm{VD}}(r)=V_{\mathrm{VD}}^C(r)+V_{\mathrm{VD}}^V(r)+V_{\mathrm{VD}}^S(r),`$ $`V_{\mathrm{VD}}^C(r)={\displaystyle \frac{1}{2m_am_b}}\left\{V_C(r)\left[𝐩^2+{\displaystyle \frac{(𝐩𝐫)^2}{r^2}}\right]\right\}_W`$ $`V_{\mathrm{VD}}^V(r)={\displaystyle \frac{1}{m_am_b}}\left\{V_V(r)𝐩^2\right\}_W`$ $`+{\displaystyle \frac{1}{4}}\left[(1\lambda _V)\left({\displaystyle \frac{1}{m_a^2}}+{\displaystyle \frac{1}{m_b^2}}\right){\displaystyle \frac{2\lambda _V}{m_am_b}}\right]`$ $`\times \left\{V_V(r)𝐩^2+V_V^{}(r){\displaystyle \frac{(𝐩𝐫)^2}{r}}\right\}_W,`$ $`V_{\mathrm{VD}}^S(r)={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{m_a^2}}+{\displaystyle \frac{1}{m_b^2}}\right)\left\{V_V(r)𝐩^2\right\}_W`$ $`+{\displaystyle \frac{1}{4}}\left[(1\lambda _S)\left({\displaystyle \frac{1}{m_a^2}}+{\displaystyle \frac{1}{m_b^2}}\right){\displaystyle \frac{2\lambda _S}{m_am_b}}\right]`$ $`\times \left\{V_V(r)𝐩^2+V_V^{}(r){\displaystyle \frac{(𝐩𝐫)^2}{r}}\right\}_W`$ and $`\{\mathrm{}\}_W`$ denotes the Weyl ordering of operators. Making the natural decomposition $`V_{\mathrm{VD}}(r)={\displaystyle \frac{1}{m_am_b}}\left\{𝐩^2V_{bc}(r)+{\displaystyle \frac{(𝐩𝐫)^2}{r^2}}V_c(r)\right\}_W`$ (32) $`+\left({\displaystyle \frac{1}{m_a^2}}+{\displaystyle \frac{1}{m_b^2}}\right)\left\{𝐩^2V_{de}(r){\displaystyle \frac{(𝐩𝐫)^2}{r^2}}V_e(r)\right\}_W`$ (33) we obtain for the corresponding structures with $`\lambda _V=1`$ and including one-loop radiative corrections in $`\overline{MS}`$ renormalization scheme: $`V_C(r)={\displaystyle \frac{4}{3}}{\displaystyle \frac{\overline{\alpha }_V(\mu ^2)}{r}}{\displaystyle \frac{4}{3}}{\displaystyle \frac{\beta _0\alpha _s^2(\mu ^2)}{2\pi }}{\displaystyle \frac{\mathrm{ln}(\mu r)}{r}},`$ (34) $`V_{bc}(r)={\displaystyle \frac{2}{3}}{\displaystyle \frac{\overline{\alpha }_V(\mu ^2)}{r}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\beta _0\alpha _s^2(\mu ^2)}{2\pi }}{\displaystyle \frac{\mathrm{ln}(\mu r)}{r}}`$ (35) $`+\left({\displaystyle \frac{1\epsilon }{2}}{\displaystyle \frac{\epsilon \lambda _S}{2}}\right)Ar+B,`$ (36) $`V_c(r)={\displaystyle \frac{2}{3}}{\displaystyle \frac{\overline{\alpha }_V(\mu ^2)}{r}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\beta _0\alpha _s^2(\mu ^2)}{2\pi }}`$ (37) $`\times \left[{\displaystyle \frac{\mathrm{ln}(\mu r)}{r}}{\displaystyle \frac{1}{r}}\right]\left({\displaystyle \frac{1\epsilon }{2}}+{\displaystyle \frac{\epsilon \lambda _S}{2}}\right)Ar,`$ (38) $`V_{de}(r)={\displaystyle \frac{\epsilon }{4}}(1+\lambda _S)Ar+B,`$ (39) $`V_e(r)={\displaystyle \frac{\epsilon }{4}}(1\lambda _S)Ar,`$ (40) where $`\overline{\alpha }_V(\mu ^2)`$ $`=`$ $`\alpha _s(\mu ^2)\left[1+\left({\displaystyle \frac{a_1}{4}}+{\displaystyle \frac{\gamma _E\beta _0}{2}}\right){\displaystyle \frac{\alpha _s(\mu ^2)}{\pi }}\right],`$ $`a_1`$ $`=`$ $`{\displaystyle \frac{31}{3}}{\displaystyle \frac{10}{9}}n_f,\beta _0=11{\displaystyle \frac{2}{3}}n_f.`$ Here $`n_f`$ is a number of flavours and $`\mu `$ is a renormalization scale. It is easy to check that the exact Barchielli, Brambilla, Prosperi relations following from the Lorentz invariance of the Wilson loop $`V_{de}{\displaystyle \frac{1}{2}}V_{bc}+{\displaystyle \frac{1}{4}}(V_C+V_0)=0,`$ (41) $`V_e+{\displaystyle \frac{1}{2}}V_c+{\displaystyle \frac{r}{4}}{\displaystyle \frac{\mathrm{d}(V_C+V_0)}{\mathrm{d}r}}=0`$ (42) are exactly satisfied. The expression for spin-dependent part of the quark-antiquark potential with the inclusion of radiative corrections can be found in ref. . Now we can calculate the mass spectra of heavy quarkonia with the account of all relativistic corrections (including retardation effects) of order $`v^2/c^2`$ and one-loop radiative corrections. For this purpose we substitute the quasipotential which is a sum of the spin-independent and spin-dependent parts into the quasipotential equation. Then we multiply the resulting expression from the left by the quasipotential wave function of a bound state and integrate with respect to the relative momentum. Taking into account the accuracy of the calculations, we can use for the resulting matrix elements the wave functions of Eq. (1) with the static potential $$V_{\mathrm{NR}}(r)=\frac{4}{3}\frac{\overline{\alpha }_V(\mu ^2)}{r}+Ar+B.$$ (43) As a result we obtain the mass formula ($`m_a=m_b=m`$) $`{\displaystyle \frac{b^2(M)}{2\mu _R}}=W+a𝐋𝐒+b[(𝐒_a𝐒_b)`$ (44) $`+{\displaystyle \frac{3}{r^2}}(𝐒_a𝐫)(𝐒_b𝐫)]+c𝐒_a𝐒_b,`$ (45) where the first term on the right-hand side of the mass formula contains all spin-independent contributions, the second term describes the spin-orbit interaction, the third term is responsible for the tensor interaction, while the last term gives the spin-spin interaction. To proceed further we need to discuss the parameters of our model. There is the following set of parameters: the quark masses ($`m_b`$ and $`m_c`$), the QCD constant $`\mathrm{\Lambda }`$ and renormalization point $`\mu `$ in the short-range part of the $`Q\overline{Q}`$ potential, the slope $`A`$ and intercept $`B`$ of the linear confining potential (10), the mixing coefficient $`\epsilon `$ (9), the long-range anomalous chromomagnetic moment $`\kappa `$ of the quark (8), and the mixing parameter $`\lambda _S`$ in the retardation correction for the scalar confining potential. We can fix the values of the parameters $`\epsilon =1`$ and $`\kappa =1`$ from the consideration of radiative decays and comparison of the heavy quark expansion in our model with the predictions of the heavy quark effective theory. We fix the slope of the linear confining potential $`A=0.18`$ GeV<sup>2</sup> which is a rather adopted value. In order to reduce the number of independent parameters we assume that the renormalization scale $`\mu `$ in the strong coupling constant $`\alpha _s(\mu ^2)`$ is equal to the quark mass. We also varied the quark masses in a reasonable range for the constituent quark masses. The numerical analysis and comparison with experimental data lead to the following values of our model parameters: $`m_c=1.55\mathrm{GeV},m_b=4.88\mathrm{GeV},A=0.18\mathrm{GeV}^2,B=0.16\mathrm{GeV},\mu =m_Q,\mathrm{\Lambda }=0.178\mathrm{GeV},\epsilon =1,\kappa =1,\lambda _S=0.`$ The quark masses $`m_{c,b}`$ have usual values for constituent quark models and coincide with those chosen in our previous analysis . The above value of the retardation parameter $`\lambda _S`$ for the scalar confining potential coincides with the minimal area low and flux tube models , with lattice results and Gromes suggestion . The found value for the QCD parameter $`\mathrm{\Lambda }`$ gives the following values for the strong coupling constants $`\alpha _s(m_c^2)0.32`$ and $`\alpha _s(m_b^2)0.22`$. The results of our numerical calculations of the mass spectra of charmonium and bottomonium (in GeV) are presented in Tables 1 and 2. We see that the calculated masses agree with experimental values within few MeV and this difference is compatible with the estimates of the higher order corrections in $`v^2/c^2`$ and $`\alpha _s`$. The model reproduces correctly both the positions of the centres of gravity of the levels and their fine and hyperfine splitting. Note that the good agreement of the calculated mass spectra with experimental data is achieved by systematic accounting for all relativistic corrections (including retardation corrections) of order $`v^2/c^2`$, both spin-dependent and spin-independent ones, while in most of potential models only the spin-dependent corrections are included. The calculated mass spectra of charmonium and bottomonium are close to the results of our previous calculation where retardation effects in the confining potential and radiative corrections to the one-gluon exchange potential were not taken into account. Both calculations give close values for the experimentally measured states as well as for the yet unobserved ones. The inclusion of radiative corrections allowed to get better results for the fine splittings of quarkonium states. Thus we can conclude from this comparison that the inclusion of retardation effects and spin-independent one-loop radiative corrections resulted only in the slight shift ($`10\%`$) in the value of the QCD parameter $`\mathrm{\Lambda }`$ and an approximately two-fold decrease of the constant $`B`$. Such changes of parameters almost do not influence the wave functions. As a result the decay matrix elements involving heavy quarkonium states remain mostly unchanged. We are grateful to the organizers for the nice meeting and stimulating discussions. Two of us (R.N.F and V.O.G.) were supported in part by Russian Foundation for Fundamental Research under Grant No. 00-02-17768.
warning/0006/physics0006057.html	ar5iv	text	# “Back-of-the-envelope” wind and altitude correction for 100 metre sprint times ## 1 Introduction Although not officially recognized by the International Amateur Athletic Federation (IAAF), correcting sprint race times for the effects of wind and altitude variation is a subject of increasing interest in the Track and Field community. With the number of men’s sub-9.90 s and women’s sub-10.80 s clockings on the rise, correcting these marks to their sea-level, 0-wind equivalents is useful in determining the overall quality of the performances (at the time of competition). A literature search reveals rather detailed experimental field studies of these effects , as well as several theoretical estimates based on mathematical and computational simulations . Physically, linear drag forces (scaled to units of mass) are expressed as $$F_d=\frac{1}{2}C_dA\rho (H)(v(t)w)^2,$$ (1) where $`C_d`$ is the drag coefficient, $`A`$ the frontal cross-sectional area, $`\rho (H)`$ the atmospheric density at altitude $`H`$, $`v(t)`$ the sprinter’s velocity, and $`w`$ the wind speed (co-linear to $`v(t)`$). It follows that $`F_d`$ will be smaller for tail-winds ($`w>0`$), and larger for head-winds ($`w<0`$) at a fixed altitude. Head-winds and tail-winds of equal magnitude will not provide equal-magnitude time corrections, due to the non-linear nature of the drag term. At 0 metres altitude with 0 wind, the base drag is $`F_0=1/2C_dA\rho _0v(t)^2`$, where $`\rho _0=1.184`$g cm<sup>-3</sup> is the mean sea-level density of air at 25 degrees Celsius. Since air density varies exponentially as a function of altitude, a convenient approximation can be written as $`\rho (H)=\rho _0\mathrm{exp}(0.000125H)`$ for the range of elevations considered herein (less than 2300 m for the majority of competition venues). The general consensus of most researchers in question is that for a 10.00 s race (average men’s world-class sprint), a tail-wind of +2 ms<sup>-1</sup> will provide an advantage of roughly 0.10 seconds (i.e. a faster time), whose value will vary slightly depending on the altitude of the competition venue. If the wind gauge reads in excess of +2 ms<sup>-1</sup>, the performance is termed wind-assisted, and is not eligible for any potential record status. Conversely with no wind, an altitude of 1000 m will produce an advantage of 0.03 s, above which performances are officially deemed altitude-assisted. Unlike wind-assisted marks, an altitude-assisted time can still count for a record. At 2000 m, the advantage will be about 0.06 s over a sea-level run. An 11.00 s time (average world-class women) will be boosted by about +0.12 s with a +2 ms<sup>-1</sup> tail-wind, and by 0.07 s (no wind) at 2000 m. As altitude increases, the magnitude of the wind effects will increase. Obviously, this is a reasonable explanation for the rash of World Records (WRs) experienced in the sprints and long jump at the 1968 Olympics in Mexico City, which resides at an altitude of approximately 2250 m. ## 2 “Back-of-the-Envelope” correction: Derivation A ”back-of-the-envelope” (BOTE) calculation is a simplified reduction of a complex (physical) model, from which one can make reasonable predictions with a minimal number of input parameters. An exact modeling of wind and altitude effects is a daunting task, since the mechanics involved are numerous and not easily representable by basic functions (see for such a model). A historically-based method of such simulations is via the velocity curve approach. This is a method of studying a sprinter’s performance, first introduced empirically by Hill in the early 1900s, and further investigated by Keller as an equation of motion of the form $$\dot{v}(t)=F_pv(t)\alpha ^1,$$ (2) The term $`F_p`$ is a propulsive term, while $`\alpha `$ is a decay term (representing internal physiological or biomechanical variables). Note that again Equations (2) is scaled in units of the sprinter’s mass $`M`$, so the interpretation of $`F_d`$ is force per unit mass (or, effectively, acceleration). Unless otherwise specified, this notation is used for the remainder of the article. This derivation roughly follows that of Reference , however the latter incorrectly estimates the numerical value of certain key variables, and omits the effects of altitude all together. In fact, the author of suggests that altitude effects on sprint times cannot be modeled by drag modification alone, which is not necessarily a correct assertion (as will be shown). Equation (2) may easily be altered to include drag effects by the addition of $`F_d`$, $$\dot{v}(t)=F_pv(t)\alpha ^1F_d,$$ (3) and a time dependence $`F_pF_p(t)`$ may also be added (see e.g. for such mechanisms). The BOTE expression presented herein, being simplistic by its namesake, does not include these substitutions, and furthermore imposes an additional simplification of time-independence, $`v(t)v`$, and $`\dot{v}(t)=0`$. To address the issue of wind and altitude correction, define as follows $`v(w,H)`$ (velocity of the sprinter at altitude $`H`$ with wind $`w`$); $`v(0,0)`$ (velocity with 0-wind, at sea-level); and $`F_d(w,H)`$ (effective drag for wind $`w`$ and altitude $`H`$). Subject to the constraint of constant velocity, Equation (3) may be rewritten as $`v(w,H)`$ $`=`$ $`\alpha [F_pF_d(w,H)],`$ $`v(0,0)`$ $`=`$ $`\alpha [F_pF_d(0,0)],`$ (4) for each case described above. Solving for $`\alpha `$ and equating the two expressions above yields $$\frac{v(0,0)}{v(w,H)}=\frac{(1F_d(0,0)/F_p)}{(1F_d(w,H)/F_p)},$$ (5) Define the ratio $`\delta =F_d(0,0)/F_p`$ as the effort required to overcome drag in 0-wind conditions at sea level. The numerical value of $`\delta `$ will be discussed shortly. Since velocity is constant (i.e. the average race velocity), one can write $`v(w,H)=100/t_{w,H}`$, where $`t_{w,H}`$ is the official time for the race under consideration, and rewrite Equation (5) as $$\frac{t_{0,0}}{t_{w,H}}=\frac{(1F_d(w,H)/F_p)}{(1\delta )},$$ (6) To simplify this expression further, note that the drag force for arbitrary $`w`$ and $`H`$ can be written as $`F_d(w,H)`$ $`=`$ $`{\displaystyle \frac{1}{2}}C_dA\rho (H)v(w,H)^2\left(1{\displaystyle \frac{w}{v(w,H)}}\right)^2`$ (7) $`=`$ $`F_d(0,0)\left({\displaystyle \frac{v(w,H)}{v_{0,0}}}\right)^2\mathrm{exp}(0.000125H)\left(1{\displaystyle \frac{wt_{w,H}}{100}}\right)^2`$ So, replacing $`(v(w,H)/v_0)=(t_{0,0}/t_{w,H})`$, $$\frac{F_d(w,H)}{F_p}=\delta \left(\frac{t_{0,0}}{t_{w,H}}\right)^2\mathrm{exp}(0.000125H)\left(1\frac{wt_{w,H}}{100}\right)^2,$$ (8) and thus $$\frac{t_{0,0}}{t_{w,H}}=\frac{1}{(1\delta )}\left[1\delta \left(\frac{t_{0,0}}{t_{w,H}}\right)^2exp(0.000125H)\left(1\frac{wt_{w,H}}{100)}\right)^2\right].$$ (9) Unfortunately, this is now a quadratic expression in $`(t_{0,0}/t_{w,H})`$, but this problem is quickly resolved by making the following substitution. Rewrite $`(t_{0,0}/t_{w,H})=1+\mathrm{\Delta }t/t_{w,H}`$, with $`\mathrm{\Delta }t=t_{0,0}t_{w,H}`$. Since $`\mathrm{\Delta }t`$ will seldom be larger than 0.3 s for a $``$10 s race, it is reasonable to make the substitution $$\left(\frac{t_{0,0}}{t_{w,H}}\right)^2=\left(1+\frac{\mathrm{\Delta }t}{t_{w,H}}\right)^21+2\frac{\mathrm{\Delta }t}{t_{w,H}}=2\frac{t_{0,0}}{t_{w,H}}1$$ (10) The numerical value of $`\delta `$ is determined as $$\delta =\frac{F_d(0,0)}{F_p}=\frac{1}{2}\frac{C_dA\rho _0v^2}{MF_p}$$ (11) (recall that the earlier definition of $`F`$ is scaled in units of inverse mass, hence the need for $`M`$ in the denominator). Pritchard initially found a value of $`\delta 0.032`$ (i.e. $`3.2\%`$ of a sprinter’s effort is required to overcome drag), however this assumed an overestimated value of the drag coefficient $`C_d=1.0`$, as well as the mean propulsive force $`F_p=12.1`$ ms<sup>-2</sup>. Current research suggests a drag coefficient of $`C_d(0.5,0.6)`$ , as well as an average $`F_p7`$ ms<sup>-2</sup> . For a $`9.9010.00`$ s race, the average velocity is between $`v=1010.1`$ m s<sup>-1</sup>. Taking the drag area to be $`C_dA=0.23`$m<sup>2</sup> (consistent with the quoted $`C_d`$ values, and cross-sectional area $`A(0.4,0.5)`$ m<sup>2</sup>, for a sprinter of mass 75 kg, one finds $`\delta 0.027`$. Since $`\delta `$ is small, $`1/(1\delta )(1+\delta )`$, and including the approximation of Equation (10), Equation (9) may be rearranged as $`{\displaystyle \frac{t_{0,0}}{t_{w,H}}}`$ $`=`$ $`\left({\displaystyle \frac{1}{1\delta }}\right)\left[1\delta \mathrm{exp}(0.000125H)\left(2{\displaystyle \frac{t_{0,0}}{t_{w,H}}}1\right)\left(1{\displaystyle \frac{wt_{w,H}}{100}}\right)^2\right]`$ (12) $``$ $`1+\delta \delta \mathrm{exp}(0.000125H)\left(1{\displaystyle \frac{wt_{w,H}}{100}}\right)^2+o(\delta ^2),`$ Thus, inserting the numerical value of $`\delta `$, one obtains the “back-of-the-envelope” calculation $$t_{0,0}t_{w,H}[1.0270.027\mathrm{exp}(0.000125H)(1wt_{w,H}/100)^2].$$ (13) For women, the input parameters $`F_p`$, $`A`$, $`v^2`$, and $`M`$ are smaller, so assuming values $`v=(100/11)=9.1`$ms<sup>-1</sup>, $`M=65`$ kg, $`F_p5`$ ms<sup>-2</sup>, and $`A0.35`$m<sup>2</sup>, $`\delta `$ remains essentially unchanged. Equation (13) provides an excellent match to the predictions of Reference , as well as those of Dapena and Linthorne . Thus, 100 metre sprint times may be corrected to their 0-wind, sea level equivalents by inputting only the official time, the wind gauge reading, and the altitude of the sporting venue. Furthermore, Equation (8) is easily programmable in most scientific calculators and personal computers, and hence may be used track-side by coaches, officials and the media immediately following a race to gauge its overall “quality”. ## 3 Applications To demonstrate the utility of Equation (13), Tables 123, and 4 present the corresponding corrections to the top five all-time men’s and women’s 100 m performances run with legal tail-winds, illegal winds ($`w>+2.0`$ ms<sup>-1</sup>), altitude effects ($`H>1000`$ m), and extreme head-winds ($`w<1`$ms<sup>-1</sup>). The current 100 m World Record (as of June 2000) of 9.79 s by Maurice Greene was run at low altitude with virtually no wind, and adjusts to 9.80 s. Note that (Table 1) the 9.86 s performances of Trinidadian Ato Boldon and Namibia’s Frank Fredericks were both run into head-winds of equal magnitude, but the altitude difference allows for a 0.02 s differential in corrected times. The former World Record of 9.84 s by Canada’s Donovan Bailey corrects to a 9.88 s. It is also interesting to note how exceptional performances can be hampered by strong head-winds (Table 4). The 10.49 s WR of the late Florence Griffith-Joyner is included in both Table 1 and Table 2, to demonstrate the common belief that this mark was strongly wind-aided (despite the fact that the official wind gauge reading was +0.0 ms<sup>-1</sup>, there is strong circumstantial evidence to suggest that the equipment malfunctioned). Griffith-Joyner’s legal personal records (PRs) correct to about 10.68 s, while her +2.1 ms<sup>-1</sup> wind-aided 10.54 s (Seoul, 1988) corrects to 10.66 s. Thus, the actual WR mark should probably be 10.60-10.65 s (or a 0-wind, 0-altitude equivalent of about 10.66-10.68 s). American Marion Jones’ current adjusted PR is 10.69 s, effectively on par with Griffith-Joyner’s best marks. From a historical perspective, it is interesting to note that Calvin Smith’s 10.04 s ($``$2.2 ms<sup>-1</sup>) in 1983 would have converted to a 0-wind 9.93 s at sea level. Smith’s actual WR of 9.93 s (+1.2 ms<sup>-1</sup>) was run at altitude (Colorado Springs, USA; 1850 m), correcting to only 10.03 s. The former WRs of Canada’s Ben Johnson, 9.79 s (+1.1 ms<sup>-1</sup>) in Seoul, SKR, and 9.83 s (+1.0 ms<sup>-1</sup>) in Rome, ITA, would correct to 9.85 s and 9.88 s, respectively. Unfortunately, these marks were stricken as a result of performance-enhancing drug infractions. ## 4 Conclusions The presented “back-of-the-envelope” calculation is simple to use, and is applicable to both men and women’s performances. An on-line JavaScript version is available at the author’s website, currently http://palmtree.physics.utoronto.ca/$``$newt/track/wind/ It is hoped that its use may be eventually adopted by the IAAF and/or other governing bodies of Athletics as a relative gauge of performance quality under differing competition conditions. Acknowledgements I thank Jesus Dapena for helpful discussions, and for providing the data of Reference prior to its publication. This work was supported in part by a Walter C. Sumner Memorial Fellowship, as well as a grant from the National Sciences and Engineering Research Council.
warning/0006/cs0006046.html	ar5iv	text	# 3-Coloring in Time 𝑂(1.3289^𝑛) ## 1 Introduction There are many known NP-complete problems including such important graph theoretic problems as coloring and independent sets. Unless P=NP, we know that no polynomial time algorithm for these problems can exist, but that does not obviate the need to solve them as efficiently as possible, indeed the fact that these problems are hard makes efficient algorithms for them especially important. We are interested in this paper in worst case analysis of algorithms for 3-coloring, a basic NP-complete problem. We will also discuss other related problems including 3-SAT, 3-edge-coloring and 3-list-coloring. Our algorithms for these problems are based on the following simple idea: to find a solution to a 3-coloring problem, it is not necessary to choose a color for each vertex (giving something like $`O(\mathrm{3}^n)`$ time). Instead, it suffices to only partially solve the problem by restricting each vertex to two of the three colors. We can then test whether the partial solution can be extended to a complete coloring in polynomial time (e.g. as a 2-SAT instance). This idea applied naively already gives a simple $`O(\mathrm{1}.\mathrm{5}^n)`$ time randomized algorithm; we improve this by taking advantage of local structure (if we choose a color for one vertex, this restricts the colors of several neighbors at once). It seems likely that our idea of only searching for a partial solution can be applied to many other combinatorial search problems. If we perform local reductions as above in a 3-coloring problem, we eventually reach a situation in which some uncolored vertices are surrounded by partially colored neighbors, and we run out of good local configurations to use. To avoid this problem, we translate our 3-coloring problem to one that also generalizes the other problems listed above: constraint satisfaction (CSP). In an $`(a,b)`$-CSP instance, we are given a collection of $`n`$ variables, each of which can be given one of $`a`$ different colors. However certain color combinations are disallowed: we also have input a collection of $`m`$ constraints, each of which forbids one coloring of some $`b`$-tuple of variables. Thus 3-satisfiability is exactly $`(\mathrm{2},\mathrm{3})`$-CSP, and 3-coloring is a special case of $`(\mathrm{3},\mathrm{2})`$-CSP in which the constraints disallow adjacent vertices from having the same color. As we show, $`(a,b)`$-CSP instances can be transformed in certain interesting and useful ways: in particular, one can transform $`(a,b)`$-CSP into $`(b,a)`$-CSP and vice versa, one can transform $`(a,b)`$-CSP into $`(max(a,b),\mathrm{2})`$-CSP, and in any $`(a,\mathrm{2})`$-CSP instance one can eliminate variables for which only two colors are allowed, reducing the problem to a smaller one of the same form. Because of this ability to eliminate partially colored variables immediately rather than saving them for a later 2-SAT instance, we can solve a $`(\mathrm{3},\mathrm{2})`$-CSP instance without running out of good local configurations. Our actual algorithm solves $`(\mathrm{3},\mathrm{2})`$-CSP by applying such reductions only until we reach instances with a certain simplified structure, which can then be solved in polynomial time as an instance of graph matching. We further improve our time bound for graph 3-vertex-coloring by using methods involving network flow to find a large set of good local reductions which we apply before treating the remaining problem as a $`(\mathrm{3},\mathrm{2})`$-CSP instance. And similarly, we solve 3-edge-coloring by using graph matching methods to find a large set of good local reductions which we apply before treating the remaining problem as a 3-vertex-coloring instance. ### 1.1 New Results We show the following: * A $`(\mathrm{3},\mathrm{2})`$-CSP instance with $`n`$ variables can be solved in worst case time $`O(\mathrm{1}.\mathrm{3}\mathrm{6}\mathrm{4}\mathrm{5}^n)`$, independent of the number of constraints. We also give a very simple randomized algorithm for solving this problem in expected time $`O(n^{O(\mathrm{1})}\mathrm{2}^{n/\mathrm{2}})O(\mathrm{1}.\mathrm{4}\mathrm{1}\mathrm{4}\mathrm{2}^n)`$. * A $`(d,\mathrm{2})`$-CSP instance with $`n`$ variables and $`d>\mathrm{3}`$ can be solved by a randomized algorithm in expected time $`O((\mathrm{0}.\mathrm{4}\mathrm{5}\mathrm{1}\mathrm{8}d)^n)`$. * 3-coloring in a graph of $`n`$ vertices can be solved in time $`O(\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^n)`$, independent of the number of edges in the graph. * 3-list-coloring (graph coloring given a list at each vertex of three possible colors chosen from some larger set) can be solved in time $`O(\mathrm{1}.\mathrm{3}\mathrm{6}\mathrm{4}\mathrm{5}^n)`$, independent of the number of edges. * 3-edge-coloring in an $`n`$-vertex graph can be solved in time $`O(\mathrm{2}^{n/\mathrm{2}})`$, again independent of the number of edges. * 3-satisfiability of a formula with $`t`$ 3-clauses can be solved in time $`O(n^{O(\mathrm{1})}+\mathrm{1}.\mathrm{3}\mathrm{6}\mathrm{4}\mathrm{5}^t)`$, independent of the number of variables or 2-clauses in the formula. Except where otherwise specified, $`n`$ denotes the number of vertices in a graph or variables in a SAT or CSP instance, while $`m`$ denotes the number of edges in a graph, constraints in an CSP instance, or clauses in a SAT problem. ### 1.2 Related Work There is a growing body of papers on worst case analysis of algorithms for NP-hard problems. Several authors have described algorithms for maximum independent sets ; the best of these is Robson’s , which takes time $`O(\mathrm{1}.\mathrm{2}\mathrm{1}\mathrm{0}\mathrm{8}^n)`$. Others have described algorithms for Boolean formula satisfiability ; the best of these satisfiability algorithms are Schöning’s, which solves 3-SAT in expected time $`O((\mathrm{4}/\mathrm{3})^n)`$ , and Hirsch’s, which solves SAT in time $`O(\mathrm{1}.\mathrm{2}\mathrm{3}\mathrm{9}^m)`$ . For three-coloring, we know of several relevant references. Lawler is primarily concerned with the general chromatic number, but he also gives the following very simple algorithm for 3-coloring: for each maximal independent set, test whether the complement is bipartite. The maximal independent sets can be listed with polynomial delay , and there are at most $`\mathrm{3}^{n/\mathrm{3}}`$ such sets , so this algorithm takes time $`O(\mathrm{1}.\mathrm{4}\mathrm{4}\mathrm{2}\mathrm{2}^n)`$. Schiermeyer gives a complicated algorithm for solving 3-colorability in time $`O(\mathrm{1}.\mathrm{4}\mathrm{1}\mathrm{5}^n)`$, based on the following idea: if there is one vertex $`v`$ of degree $`n\mathrm{1}`$ then the graph is 3-colorable iff $`Gv`$ is bipartite, and the problem is easily solved. Otherwise, Schiermeyer performs certain reductions involving maximal independent sets that attempt to increase the degree of $`G`$ while partitioning the problem into subproblems, at least one of which will remain solvable. Our $`O(\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^n)`$ bound significantly improves both of these results. There has also been some related work on approximate or heuristic 3-coloring algorithms. Blum and Karger show that any 3-chromatic graph can be colored with $`\stackrel{\textcolor[rgb]{1,0,0}{tilde}}{O}(n^{\mathrm{3}/\mathrm{1}\mathrm{4}})`$ colors in polynomial time. Alon and Kahale describe a technique for coloring random 3-chromatic graphs in expected polynomial time, and Petford and Welsh present a randomized algorithm for 3-coloring graphs which also works well empirically on random graphs although they prove no bounds on its running time. Finally, Vlasie has described a class of instances which are (unlike random 3-chromatic graphs) difficult to color. Very recently, Schöning has described a simple and powerful randomized algorithm for $`k`$-SAT and more general constraint satisfaction problems, including the CSP instances that we use in our solution of 3-coloring. However, for $`(d,\mathrm{2})`$-CSP, Schöning notes that his method is not as good as a randomized approach based on an idea from our previous conference paper : simply choose a random pair of values for each variable and solve the resulting 2-SAT instance in polynomial time. The table below compares the resulting $`(d/\mathrm{2})^n`$ bound with our new results; an entry with value $`x`$ in column $`d`$ indicates a time bound of $`O(x^n)`$ for $`(d,\mathrm{2})`$-CSP. We were unable to locate prior work on worst case edge coloring. Since any 3-edge-chromatic graph has at most $`\mathrm{3}n/\mathrm{2}`$ edges, one can transform the problem to 3-vertex-coloring at the expense of increasing $`n`$ by a factor of $`\mathrm{3}/\mathrm{2}`$. If we applied our vertex coloring algorithm we would then get time $`O(\mathrm{1}.\mathrm{5}\mathrm{3}\mathrm{1}\mathrm{9}^n)`$ which is significantly improved by the bound stated above. It is interesting that, historically, until the work of Schöning , the time bounds for 3-coloring have been smaller than those for 3-satisfiability (in terms of the number of vertices or variables respectively). Schöning’s $`O((\mathrm{4}/\mathrm{3})^n)`$ bound for 3-SAT reversed this pattern by being smaller than the previous $`O(\mathrm{1}.\mathrm{3}\mathrm{4}\mathrm{4}\mathrm{3}^n)`$ bound for 3-coloring from our 1995 conference paper . The present work restores 3-coloring to a smaller time bound than 3-SAT. ## 2 Constraint Satisfaction Problems We now describe a common generalization of satisfiability and graph coloring as a constraint satisfaction problem (CSP) . We are given a collection of $`n`$ variables, each of which has a list of possible colors allowed. We are also given a collection of $`m`$ constraints, consisting of a tuple of variables and a color for each variable. A constraint is satisfied by a coloring if not every variable in the tuple is colored in the way specified by the constraint. We would like to choose one color from the allowed list of each variable, in a way not conflicting with any constraints. For instance, 3-satisfiability can easily be expressed in this form. Each variable of the satisfiability problem may be colored (assigned the value) either true ($`T`$) or false ($`F`$). For each clause like $`(x_\mathrm{1}x_\mathrm{2}\neg x_\mathrm{3})`$, we make a constraint $`((v_\mathrm{1},F),(v_\mathrm{2},F),(v_\mathrm{3},T))`$. Such a constraint is satisfied if and only if at least one of the corresponding clause’s terms is true. In the $`(a,b)`$-CSP problem, we restrict our attention to instances in which each variable has at most $`a`$ possible colors and each constraint involves at most $`b`$ variables. The CSP instance constructed above from a 3-SAT instance is then a $`(\mathrm{2},\mathrm{3})`$-CSP instance, and in fact 3-SAT is easily seen to be equivalent to $`(\mathrm{2},\mathrm{3})`$-CSP. In this paper, we will concentrate our attention instead on $`(\mathrm{3},\mathrm{2})`$-CSP and $`(\mathrm{4},\mathrm{2})`$-CSP. We can represent a $`(d,\mathrm{2})`$-CSP instance graphically, by interpreting each variable as a vertex containing up to $`d`$ possible colors, and by drawing edges connecting incompatible pairs of vertex colors (Figure 1). Note that this graphical structure is not actually a graph, as the edges connect colors within a vertex rather than the vertices themselves. However, graph $`\mathrm{3}`$-colorability and graph $`\mathrm{3}`$-list-colorability can be translated directly to a form of $`(\mathrm{3},\mathrm{2})`$-CSP: we keep the original vertices of the graph and their possible colors, and add up to three constraints for each edge of the graph to enforce the condition that the edge’s endpoints have different colors (Figure 2). Of course, since these problems are all NP-complete, the theory of NP-completeness provides translations from one problem to the other, but the translations above are size-preserving and very simple. We will later describe more complicated translations from 3-coloring and 3-edge-coloring to $`(\mathrm{3},\mathrm{2})`$-CSP in which the input graph is partially colored before treating the remaining graph as an CSP instance, leading to improved time bounds over our pure CSP algorithm. As we now show, $`(a,b)`$-CSP instances can be transformed in certain interesting and useful ways. We first describe a form of duality that transforms $`(a,b)`$-CSP instances into $`(b,a)`$-CSP instances, exchanging constraints for variables and vice versa. ###### Lemma 1 If we are given an $`(a,b)`$-CSP instance, we can find an equivalent $`(b,a)`$-CSP instance in which each constraint of the $`(a,b)`$-CSP instance corresponds to a single variable of the transformed problem, and each constraint of the transformed problem corresponds to a single variable of the original problem. Proof: An assignment of colors to the original $`(a,b)`$-CSP instance’s variables solves the problem if and only if, for each constraint, there is at least one pair $`(V,C)`$ in the constraint that does not appear in the coloring. In our transformed problem, we choose one variable per original constraint, with the colors available to the new variable being these pairs $`(V,C)`$ in the corresponding constraint in the original problem. Choosing such a pair in a coloring of the transformed problem is interpreted as ruling out $`C`$ as a possible color for $`V`$ in the original problem. We then add constraints to our transformed problem to ensure that for each $`V`$ there remains at least one color that is not ruled out: we add one constraint for each $`a`$-tuple of colors of new variables—recall that each such color is a pair $`(V,C)`$—such that all colors in the $`a`$-tuple involve the same original variable $`V`$ and exhaust all the choices of colors for $`V`$. This duality may be easier to understand with a small example. As discussed above, 3-SAT is essentially the same as $`(\mathrm{2},\mathrm{3})`$-CSP, so Lemma 1 can be used to translate 3-SAT to $`(\mathrm{3},\mathrm{2})`$-CSP. Suppose we start with the 3-SAT instance $`(x_\mathrm{1}x_\mathrm{2}\neg x_\mathrm{3})(\neg x_\mathrm{1}x_\mathrm{3}x_\mathrm{4})(x_\mathrm{1}\neg x_\mathrm{2}\neg x_\mathrm{4})`$. Then we make a $`(\mathrm{3},\mathrm{2})`$-CSP instance (Figure 3) with three variables $`v_i`$, one for each 3-SAT clause. Each variable has three possible colors: $`(\mathrm{1},\mathrm{2},\mathrm{3})`$ for $`v_i`$, $`(\mathrm{1},\mathrm{3},\mathrm{4})`$ for $`v_\mathrm{2}`$, and $`(\mathrm{1},\mathrm{2},\mathrm{4})`$ for $`v_\mathrm{3}`$. The requirement that value $`T`$ or $`F`$ be available to $`x_\mathrm{1}`$ corresponds to the constraints $`((v_\mathrm{1},\mathrm{1}),(v_\mathrm{2},\mathrm{1}))`$ and $`((v_\mathrm{2},\mathrm{1}),(v_\mathrm{3},\mathrm{1}))`$; we similarly get constraints $`((v_\mathrm{1},\mathrm{2}),(v_\mathrm{3},\mathrm{2}))`$, $`((v_\mathrm{1},\mathrm{3}),(v_\mathrm{2},\mathrm{3}))`$, and $`((v_\mathrm{2},\mathrm{4}),(v_\mathrm{3},\mathrm{4}))`$. One possible coloring of this $`(\mathrm{3},\mathrm{2})`$-CSP instance would be to color $`v_\mathrm{1}`$ $`\mathrm{1}`$, $`v_\mathrm{2}`$ $`\mathrm{3}`$, and $`v_\mathrm{3}`$ $`\mathrm{4}`$; this would give satisfying assignments in which $`x_\mathrm{1}`$ and $`x_\mathrm{3}`$ are $`T`$, $`x_\mathrm{4}`$ is $`F`$, and $`x_\mathrm{2}`$ can be either $`T`$ or $`F`$. We can similarly translate an $`(a,a)`$-CSP instance into an $`(a,\mathrm{2})`$-CSP instance in which each variable corresponds to either a constraint or a variable, and each constraint forces the variable colorings to match up with the dual constraint colorings; we omit the details as we do not use this construction in our algorithms. ## 3 Simplification of CSP Instances Before we describe our CSP algorithms, we describe some situations in which the number of variables in an CSP instance may be reduced with little computational effort. ###### Lemma 2 Let $`v`$ be a variable in an $`(a,\mathrm{2})`$-CSP instance, such that only two of the $`a`$ colors are allowed at $`v`$. Then we can find an equivalent $`(a,\mathrm{2})`$-CSP instance with one fewer variable. Proof: Let the two colors allowed at $`v`$ be $`R`$ and $`G`$. Define $`\mathrm{conflict}(C)`$ to be the set of pairs $`\{(u,A):((u,A),(v,C))`$ is a constraint. We then include $`\mathrm{conflict}(R)\times \mathrm{conflict}(G)`$ to our set of constraints. Any pair $`((u,A),(w,B))\mathrm{conflict}(R)\times \mathrm{conflict}(G)`$ does not reduce the space of solutions to the original problem since if both $`(u,A)`$ and $`(w,B)`$ were present in a coloring there would be no possible color left for $`v`$. Conversely if all such constraints are satisfied, one of the two colors for $`v`$ must be available. Therefore we can now find a smaller equivalent problem by removing $`v`$, as shown in Figure 4. When we apply this variable elimination scheme, the number of constraints can increase, but there can exist only $`(an)^\mathrm{2}`$ distinct constraints, which in our applications will be a small polynomial. ###### Lemma 3 Let $`(v,X)`$ and $`(w,Y)`$ be (variable,color) pairs in an $`(a,\mathrm{2})`$-CSP instance, such that $`vw`$ the only constraints involving these pairs are either of the form $`((v,X),(w,Z))`$ with $`YZ`$, or $`((v,Z),(w,Y))`$ with $`XZ`$. Then we can find an equivalent $`(a,\mathrm{2})`$-CSP instance with two fewer variables. Proof: It is safe to choose the colors $`(v,X)`$ and $`(w,Y)`$, since these two choices do not conflict with each other nor with anything else in the CSP instance. ###### Lemma 4 Let $`(v,R)`$ and $`(v,B)`$ be (variable,color) pairs in an $`(a,\mathrm{2})`$-CSP instance, such that whenever the instance contains a constraint $`((v,R),(w,X))`$ it also contains a constraint $`((v,B),(w,X))`$. Then we can find an equivalent $`(a,\mathrm{2})`$-CSP instance with one fewer variable. Proof: Any solution involving $`(v,B)`$ can be changed to one involving $`(v,R)`$ without violating any additional constraints, so it is safe to remove the option of coloring $`v`$ with color $`B`$. Once we remove this option, $`v`$ is restricted to two colors, and we can apply Lemma 2. ###### Lemma 5 Let $`(v,R)`$ be a (variable,color) pair in an $`(a,b)`$-CSP instance that is not involved in any constraints. Then we can find an equivalent $`(a,b)`$-CSP instance with one fewer variable. Proof: We may safely assign color $`R`$ to $`v`$ and remove it from the instance. ###### Lemma 6 Let $`(v,R)`$ be a (variable,color) pair in an $`(a,\mathrm{2})`$-CSP instance that is involved in constraints with all three color options of another variable $`w`$. Then we can find an equivalent $`(a,b)`$-CSP instance with one fewer variable. Proof: No coloring of the instance can use $`(v,R)`$, so we can restrict $`v`$ to the remaining two colors and apply Lemma 2. We say that a CSP instance in which none of Lemmas 26 applies is reduced. ## 4 Simple Randomized CSP Algorithm We first demonstrate the usefulness of Lemma 2 by describing a very simple randomized algorithm for solving $`(\mathrm{3},\mathrm{2})`$-CSP instances in expected time $`O(\mathrm{2}^{n/\mathrm{2}}n^{O(\mathrm{1})})`$. ###### Lemma 7 If we are given a $`(\mathrm{3},\mathrm{2})`$-CSP instance $`I`$, then in random polynomial time we can find an instance $`I^{}`$ with two fewer variables, such that if $`I^{}`$ is solvable then so is $`I`$, and if $`I`$ is solvable then with probability at least $`\frac{\mathrm{1}}{\mathrm{2}}`$ so is $`I^{}`$. Proof: If no constraint exists, we can solve the problem immediately. Otherwise choose some constraint $`((v,X),(w,Y))`$. Rename the colors if necessary so that both $`v`$ and $`w`$ have available the same three colors $`R`$, $`G`$, and $`B`$, and so that $`X=Y=R`$. Restrict the colorings of $`v`$ and $`w`$ to two colors each in one of four ways, chosen uniformly at random from the four possible such restrictions in which exactly one of $`v`$ and $`w`$ is restricted to colors $`G`$ and $`B`$ (Figure 5). Then it can be verified by examination of cases that any valid coloring of the problem remains valid for exactly two of these four restrictions, so with probability $`\frac{\mathrm{1}}{\mathrm{2}}`$ it continues to be a solution to the restricted problem. Now apply Lemma 2 and eliminate both $`v`$ and $`w`$ from the problem. ###### Corollary 1 In expected time $`O(\mathrm{2}^{n/\mathrm{2}}n^{O(\mathrm{1})})`$ we can find a solution to a $`(\mathrm{3},\mathrm{2})`$-CSP instance if one exists. Proof: We perform the reduction above $`n/\mathrm{2}`$ times, taking polynomial time and giving probability at least $`\mathrm{2}^{n/\mathrm{2}}`$ of finding a correct solution. If we repeat this method until a solution is found, the expected number of repetitions is $`\mathrm{2}^{n/\mathrm{2}}`$. ## 5 Faster CSP Algorithm We now describe a more complicated method of solving $`(\mathrm{3},\mathrm{2})`$-CSP instances deterministically with the somewhat better time bound of $`O(\mathrm{1}.\mathrm{3}\mathrm{6}\mathrm{4}\mathrm{4}\mathrm{3}^n)`$. More generally, our algorithm can actually handle $`(\mathrm{4},\mathrm{2})`$-CSP instances. Any $`(\mathrm{4},\mathrm{2})`$-CSP instance can be transformed into a $`(\mathrm{3},\mathrm{2})`$-CSP instance by expanding each of its four-color variables to two three-color variables, each having two of the original four colors, with a constraint connecting the third color of each new variable (Figure 6). Therefore, the natural definition of the “size” of a $`(\mathrm{4},\mathrm{2})`$-CSP instance is $`n=n_\mathrm{3}+\mathrm{2}n_\mathrm{4}`$, where $`n_i`$ denotes the number of variables with $`i`$ colors. However, we instead define the size to be $`n=n_\mathrm{3}+(\mathrm{2}ϵ)n_\mathrm{4}`$, where $`ϵ\mathrm{0}.\mathrm{0}\mathrm{9}\mathrm{5}\mathrm{5}\mathrm{4}\mathrm{3}`$ is a constant to be determined more precisely later. In any case, the size of a $`(\mathrm{3},\mathrm{2})`$-CSP instance remains equal to its number of variables, so any bound on the running time of our algorithm in terms of $`n`$ applies directly to $`(\mathrm{3},\mathrm{2})`$-CSP. The basic idea of our algorithm is to find a set of local configurations that must occur within any $`(\mathrm{4},\mathrm{2})`$-CSP instance $`I`$, such that any instance containing such a configuration can be replaced by a small number of smaller instances. In more detail, for each configuration we describe a set of smaller instances $`I_i`$ of size $`\|I\|r_i`$ such that $`I`$ is solvable if and only if at least one of the instances $`I_i`$ is solvable. If one particular configuration occurred at each step of the algorithm, this would lead to a recurrence of the form $$T(n)=T(nr_i)+\text{poly}(n)=O(\lambda (r_\mathrm{1},r_\mathrm{2},\mathrm{})^n)$$ for the worst-case running time of our algorithm, where the base $`\lambda (r_\mathrm{1},r_\mathrm{2},\mathrm{})`$ of the exponent in the running time is the largest zero of the function $`f(x)=\mathrm{1}x^{r_i}`$ (such a function is not necessarily a polynomial because the $`r_i`$ will not necessarily be integers). We call this value $`\lambda (r_\mathrm{1},r_\mathrm{2},\mathrm{})`$ the work factor of the given local configuration. The overall time bound will be $`\lambda ^n`$ where $`\lambda `$ is the largest work factor among the configurations we have identified. This value $`\lambda `$ will depend on our previous choice of $`ϵ`$; we will choose $`ϵ`$ in such a way as to minimize $`\lambda `$. ### 5.1 Single Constraints and Multiple Adjacencies We first consider local configurations in which some (variable,color) pair is incident on only one constraint, or has multiple constraints to the same variable. First, suppose that (variable,color) pair $`(v,R)`$ is involved in only a single constraint $`((v,R),(w,R))`$. If this is also the only constraint involving $`(w,R)`$, we call it an isolated constraint. Otherwise, we call it a dangling constraint. ###### Lemma 8 Let $`((v,R),(w,R))`$ be an isolated constraint in a $`(\mathrm{4},\mathrm{2})`$-CSP instance, and let $`ϵ\mathrm{0}.\mathrm{5}\mathrm{4}\mathrm{5}`$. Then the instance can be replaced by smaller instances with work factor at most $`\lambda (\mathrm{2}ϵ,\mathrm{3}ϵ)`$. Proof: If $`v`$ and $`w`$ are both three-color variables, then the instance can be colored if and only if we can color the instance formed by replacing them with a single four-color variable, in which the four colors are the remaining choices for $`v`$ and $`w`$ other than $`R`$ (Figure 6). Thus in this case we can reduce the problem size by $`ϵ`$, with no additional work. Otherwise, if there exists a coloring of the given instance, there exists one in which exactly one of $`v`$ and $`w`$ is given color $`R`$. Suppose first that $`v`$ has four colors while $`w`$ has only three. Thus we can reduce the problem to two instances, in one of which $`(v,R)`$ is used (so $`v`$ is removed from the problem, and $`(w,R)`$ is removed as a choice for variable $`w`$, allowing us to remove the variable by Lemma 2) and in the other of which $`(w,R)`$ is used (Figure 7). The first subproblem has its size reduced by $`\mathrm{3}ϵ`$ since both variables are removed, while the second’s size is reduced by $`\mathrm{2}ϵ`$ since $`w`$ is removed while $`v`$ loses one of its colors but is not removed. Thus the work factor is $`\lambda (\mathrm{2}ϵ,\mathrm{3}ϵ)`$. Similarly, if both are four-color variables, the work factor is $`\lambda (\mathrm{3}\mathrm{2}ϵ,\mathrm{3}\mathrm{2}ϵ)`$. For the given range of $`ϵ`$, this second work factor is smaller than the first. ###### Lemma 9 Let $`((v,R),(w,R))`$ be a dangling constraint in a reduced $`(\mathrm{4},\mathrm{2})`$-CSP instance. Then the instance can be replaced by smaller instances with work factor at most $`\lambda (\mathrm{2}ϵ,\mathrm{3}ϵ)`$. Proof: The second constraint for $`(w,R)`$ can not involve $`v`$, or we would be able to apply Lemma 4. We choose either to use color $`(w,R)`$ or to restrict $`w`$ to avoid that color (Figure 8). If we use color $`(w,R)`$, we eliminate choice $`(v,R)`$ and another choice on the other neighbor of $`w`$. If we avoid color $`(w,R)`$, we may safely use color $`(v,R)`$. In the worst case, the other neighbor of $`(w,R)`$ has four colors, so removing one only reduces the problem size by $`\mathrm{1}ϵ`$. There are four cases depending on the number of colors of $`v`$ and $`w`$: If both have three colors, the work factor is $`\lambda (\mathrm{2},\mathrm{3}ϵ)`$. If only $`v`$ has four colors, the work factor is $`\lambda (\mathrm{3}ϵ,\mathrm{3}\mathrm{2}ϵ)`$. If only $`w`$ has four colors, the work factor is $`\lambda (\mathrm{2}ϵ,\mathrm{4}\mathrm{2}ϵ)`$. If both have four colors, the work factor is $`\lambda (\mathrm{3}\mathrm{2}ϵ,\mathrm{4}\mathrm{3}ϵ)`$. These factors are all dominated by the one in the statement of the lemma. ###### Lemma 10 Suppose a reduced $`(\mathrm{4},\mathrm{2})`$-CSP instance includes two constraints such as $`((v,R),(w,B))`$ and $`((v,R),(w,G))`$ that connect one color of variable $`v`$ with two colors of variable $`w`$, and let $`ϵ\mathrm{0}.\mathrm{4}`$. Then the instance can be replaced by smaller instances with work factor at most $`\lambda (\mathrm{2}ϵ,\mathrm{3}\mathrm{2}ϵ)`$. Proof: We assume that the instance has no color choice with only a single constraint, or we could apply one of Lemmas 8 and 9 to achieve the given work factor. We say that $`(v,R)`$ implies $`(w,R)`$ if there are constraints from $`(v,R)`$ to every other color choice of $`w`$. If the target $`(w,R)`$ of an implication is not the source of another implication, then using $`(w,R)`$ eliminates $`w`$ and at least two other colors, while avoiding $`(w,R)`$ forces us to also avoid $`(v,R)`$ (Figure 9). Thus, in this case we achieve work factor either $`\lambda (\mathrm{2}ϵ,\mathrm{3}\mathrm{2}ϵ)`$ if $`w`$ has three color choices, or $`\lambda (\mathrm{2}\mathrm{2}ϵ,\mathrm{4}\mathrm{3}ϵ)`$ if it has four. If the target of every implication is the source of another, then we can find a cycle of colors each of which implies the next in the cycle (Figure 10). If no other constraints involve colors in the cycle (as is true in the figure), we can use them all, reducing the problem by the length of the cycle for free. Otherwise, let $`(v,R)`$ be a color in the cycle that has an outside constraint. If we use $`(v,R)`$, we must use the colors in the rest of the cycle, and eliminate the (variable,color) pair outside the cycle constrained by $`(v,R)`$. If we avoid $`(v,R)`$, we must also avoid the colors in the rest of the cycle. The maximum work factor for this case is $`\lambda (\mathrm{2},\mathrm{3}ϵ)`$, and arises when the cycle consists of only two variables, both of which have only three allowed colors. Finally, if the situation described in the lemma exists without forming any implication, then $`w`$ must have four color choices, exactly two of which are constrained by $`(v,R)`$. In this case restricting $`w`$ to those two choices reduces the size by at least $`\mathrm{3}\mathrm{2}ϵ`$, while restricting it to the remaining two choices reduces the size by $`\mathrm{2}ϵ`$, again giving work factor $`\lambda (\mathrm{2}ϵ,\mathrm{3}\mathrm{2}ϵ)`$. ### 5.2 Highly Constrained Colors We next consider cases in which choosing one color for a variable eliminates many other choices, or in which adjacent (variable,color) pairs have different numbers of constraints. ###### Lemma 11 Suppose a reduced $`(\mathrm{4},\mathrm{2})`$-CSP instance includes a color pair $`(v,R)`$ involved in three or more constraints, where $`v`$ has four color choices, or a pair $`(v,R)`$ involved in four or more constraints, where $`v`$ has three color choices. Then the instance can be replaced by smaller instances with work factor at most $`\lambda (\mathrm{1}ϵ,\mathrm{5}\mathrm{4}ϵ)`$. Proof: We can assume from Lemma 10 that each constraint connects $`(v,R)`$ to a different variable. Then if we choose to use color $`(v,R)`$, we eliminate $`v`$ and remove a choice from each of its neighbors, either eliminating them or reducing their number of choices from four to three. If we don’t use $`(v,R)`$, we eliminate that color only. So if $`v`$ has four choices, the work factor is at most $`\lambda (\mathrm{1}ϵ,\mathrm{5}\mathrm{4}ϵ)`$, and if it has three choices and four or more constraints, the work factor is at most $`\lambda (\mathrm{1},\mathrm{5}\mathrm{4}ϵ)`$. ###### Lemma 12 Suppose a reduced $`(\mathrm{4},\mathrm{2})`$-CSP instance includes a (variable,color) pair $`(v,R)`$ with three constraints, one of which connects it to a variable with four color choices, and let $`ϵ\mathrm{0}.\mathrm{3}\mathrm{5}\mathrm{7}\mathrm{6}`$. Suppose also that none of the previous lemmas applies. Then the instance can be replaced by smaller instances with work factor at most $`\lambda (\mathrm{3}ϵ,\mathrm{4}ϵ,\mathrm{4}ϵ)`$. Proof: For convenience suppose that the four-color neighbor is $`(w,R)`$. We can assume $`(w,R)`$ has only two constraints, else it would be covered by a previous lemma. Then, if $`(v,R)`$ and $`(w,R)`$ do not form a triangle with a third (variable,color) pair (Figure 11, left), we choose either to use or avoid color $`(v,R)`$. If we use $`(v,R)`$, we eliminate $`v`$ and the three adjacent color choices. If we avoid $`(v,R)`$, we create a dangling constraint at $`(w,R)`$, which we have seen in Lemma 9 allows us to further subdivide the instance with work factor $`\lambda (\mathrm{3}ϵ,\mathrm{3}\mathrm{2}ϵ)`$ in addition to the elimination of $`v`$. Thus, the overall work factor in this case is $`\lambda (\mathrm{4}ϵ,\mathrm{4}\mathrm{2}ϵ,\mathrm{4}\mathrm{3}ϵ)`$. On the other hand, suppose we have a triangle of constraints formed by $`(v,R)`$, $`(w,R)`$, and a third (variable,color) pair $`(x,R)`$, as shown in Figure 11, right. Then $`(v,R)`$ and $`(x,R)`$ are the only choices constraining $`(w,R)`$, so if $`(v,R)`$ and $`(x,R)`$ are both not chosen, we can safely choose to use color $`(w,R)`$. Therefore, we make three smaller instances, in each of which we choose to use one of the three choices in the triangle. We can assume from the previous cases that $`(v,R)`$ has only three choices, and further its third neighbor (other than $`(w,R)`$ and $`(x,R)`$) must also have only three choices or we could apply the previous case of the lemma. In the worst case, $`(x,R)`$ has only two constraints and $`x`$ has only three color choices. Therefore, the size of the subproblems formed by choosing $`(v,R)`$, $`(w,R)`$, and $`(x,R)`$ is reduced by at least $`\mathrm{4}ϵ`$, $`\mathrm{4}ϵ`$, and $`\mathrm{3}ϵ`$ respectively, leading to a work factor of $`\lambda (\mathrm{3}ϵ,\mathrm{4}ϵ,\mathrm{4}ϵ)`$. If instead $`x`$ has four color choices, we get the better work factor $`\lambda (\mathrm{4}\mathrm{2}ϵ,\mathrm{4}\mathrm{2}ϵ,\mathrm{4}\mathrm{2}ϵ)`$. For the given range of $`ϵ`$, the largest of these work factors is $`\lambda (\mathrm{3}ϵ,\mathrm{4}ϵ,\mathrm{4}ϵ)`$. ###### Lemma 13 Suppose a reduced $`(\mathrm{4},\mathrm{2})`$-CSP instance includes a (variable,color) pair $`(v,R)`$ with three constraints, one of which connects it to a variable with two constraints. Suppose also that none of the previous lemmas applies. Then the instance can be replaced by smaller instances with work factor at most $`max\{\lambda (\mathrm{1}+ϵ,\mathrm{4}),\lambda (\mathrm{3},\mathrm{4}ϵ,\mathrm{4})\}`$. Proof: Let $`(w,R)`$ be the neighbor with two constraints. Note that (since the previous lemma is assumed not to apply) all neighbors of $`(v,R)`$ have only three color choices. First, suppose $`(v,R)`$ and $`(w,R)`$ are not part of a triangle of constraints (Figure 12, top). Then, if we choose to use color $`(v,R)`$ we eliminate four variables, while if we avoid using it we create a dangling constraint on $`(w,R)`$ which we further subdivide into two more instances according to Lemma 9. Thus, the work factor in this case is $`\lambda (\mathrm{3},\mathrm{4}ϵ,\mathrm{4})`$. Second, suppose that $`(v,R)`$ and $`(w,R)`$ are part of a triangle with a third (variable,color) pair $`(x,R)`$, and that $`(x,R)`$ has three constraints (Figure 12, bottom left). Then (as in the previous lemma) we may choose to use one of the three choices in the triangle, resulting in work factor $`\lambda (\mathrm{3},\mathrm{4},\mathrm{4})`$. Finally, suppose that $`(v,R)`$, $`(w,R)`$, and $`(x,R)`$ form a triangle as above, but that $`(x,R)`$ has only two constraints (Figure 12, bottom right). Then if we choose to use $`(v,R)`$ we eliminate four variables, while if we avoid using it we create an isolated constraint between $`(w,R)`$ and $`(x,R)`$. Thus in this case the work factor is $`\lambda (\mathrm{1}+ϵ,\mathrm{4})`$. If none of the above lemmas applies to an instance, then each color choice in the instance must have either two or three constraints, and each neighbor of that choice must have the same number of constraints. ### 5.3 Triply-Constrained Colors Within this section we assume that we have a $`(\mathrm{4},\mathrm{2})`$-CSP instance in which none of the previous reduction lemmas applies, so any (variable,color) pair must be involved in exactly as many constraints as each of its neighbors. We now consider the remaining (variable,color) pairs that have three constraints each. Define a three-component to be a subset of such pairs such that any pair in the subset is connected to any other by a path of constraints. We distinguish two such types of components: a small three-component is one that involves only four distinct variables, while a large three-component involves five or more variables. Note that we can assume by the previous lemmas that each variable in a component has only three color choices. ###### Lemma 14 Let $`C`$ be a small three-component involving $`k`$ (variable,color) pairs. Then $`k`$ must be a multiple of four, and each variable involved in the component has exactly $`k/\mathrm{4}`$ pairs in $`C`$. Proof: Let $`v`$ and $`w`$ be variables in a small component $`C`$. Then each (variable,color) pair in $`C`$ from variable $`v`$ has exactly one constraint to a distinct (variable,color) pair from variable $`w`$, so the numbers of pairs from $`v`$ equals the number of pairs from $`w`$. The assertions that each variable has the same number of pairs, and that the total number of pairs is a multiple of four, then follow. We say that a small three-component is good if $`k=\mathrm{4}`$ in the lemma above. ###### Lemma 15 Let $`C`$ be a small three-component that is not good. Then the instance can be replaced by smaller instances with work factor at most $`\lambda (\mathrm{4},\mathrm{4},\mathrm{4})`$. Proof: A component with $`k=\mathrm{1}\mathrm{2}`$ uses up all color choices for all four variables. Thus we may consider these variables in isolation from the rest of the instance, and either color them all (if possible) or determine that the instance is unsolvable. The remaining small components have $`k=\mathrm{8}`$. Such a component may be drawn with the four variables at the corners of a square, and the top, left, and right pairs of edges uncrossed (Figure 13). If only the center two pairs were crossed, we would actually have two $`k=\mathrm{4}`$ components, and if any other two or three of the remaining pairs were crossed, we could reduce the number of crossings in the drawing by swapping the colors at one of the variables. Thus, the only possible small components with $`k=\mathrm{8}`$ are the one with all six pairs uncrossed, and the one with only one pair crossed. The first of these allows all four variables to be colored and removed, while in the other case there exist only three maximal subsets of variables that can be colored. (In the figure, these three sets are formed by the bottom two vertices, and the two sets formed by removing one bottom vertex). We split into instances by choosing to color each of these maximal subsets, eliminating all four variables in the component and giving work factor $`\lambda (\mathrm{4},\mathrm{4},\mathrm{4})`$. Define a witness to a large three-component to be a set of five (variable,color) pairs with five distinct variables, such that there exist constraints from one pair to three others, and from at least one of those three to the fifth. By convention we use $`(v,R)`$ to denote the first pair, $`(w,R)`$, $`(x,R)`$, and $`(y,R)`$ to denote the pairs connected by constraints to $`(v,R)`$, and $`(z,R)`$ to be the fifth pair in the witness. ###### Lemma 16 Every large three-component has a witness. Proof: Choose some arbitrary pair $`(u,R)`$ as a starting point, and perform a breadth first search in the graph formed by the pairs and constraints in the component. Let $`(z,R)`$ be the first pair reached by this search where $`z`$ is not one of the variables adjacent to $`(u,R)`$, let $`(v,R)`$ be the grandparent of $`(z,R)`$ in the breadth first search tree, and let the other three pairs be the neighbors of $`(v,R)`$. Then it is easy to see that $`(v,R)`$ and its neighbors must use the same four variables as $`(u,R)`$ and its neighbors, while $`z`$ by definition uses a different variable. ###### Lemma 17 Suppose that a $`(\mathrm{4},\mathrm{2})`$-CSP instance contains a large three-component. Then the instance can be replaced by smaller instances with work factor at most $`\lambda (\mathrm{4},\mathrm{4},\mathrm{5},\mathrm{5})`$. Proof: Let $`(v,R)`$, $`(w,R)`$, $`(x,R)`$, $`(y,R)`$, and $`(z,R)`$ be a witness for the component. Then we distinguish subcases according to how many of the neighbors of $`(z,R)`$ are pairs in the witness. 1. If $`(z,R)`$ has a constraint with only one pair in the witness, say $`(w,R)`$, then we choose either to use color $`(z,R)`$ or to avoid it. If we use it, we eliminate some four variables. If we avoid it, then we cause $`(w,R)`$ to have only two constraints. If $`(w,R)`$ is also constrained by one of $`(x,R)`$ or $`(y,R)`$, we then have a triangle of constraints (Figure 14, top left). We can assume without loss of generality that the remaining constraint from this triangle does not connect to a different color of variable $`z`$, for if it did we could instead use the same five variables in a different order to get a witness of this form. We then further subdivide into three more instances, in each of which we choose to use one of the pairs in the triangle, as in the second case of Lemma 13. This gives overall work factor $`\lambda (\mathrm{4},\mathrm{4},\mathrm{5},\mathrm{5})`$. On the other hand, if $`(v,R)`$ and $`(w,R)`$ are not part of a triangle (Figure 14, top right), then (after avoiding $`(z,R)`$) we can apply the first case of Lemma 13 again achieving the same work factor. 2. If $`(z,R)`$ has constraints with two pairs in the witness (Figure 14, bottom left), then choosing to use $`(z,R)`$ eliminates four variables and causes $`(v,R)`$ to dangle, while avoiding $`(z,R)`$ eliminates a single variable. The work factor is thus $`\lambda (\mathrm{1},\mathrm{6},\mathrm{7})`$. 3. If $`(z,R)`$ has constraints with all three of $`(w,R)`$, $`(y,R)`$, and $`(z,R)`$ (Figure 14, bottom right), then choosing to use $`(z,R)`$ also allows us to use $`(v,R)`$, eliminating five variables. The work factor is $`\lambda (\mathrm{1},\mathrm{5})`$. The largest of the three work factors arising in these cases is the first one, $`\lambda (\mathrm{4},\mathrm{4},\mathrm{5},\mathrm{5})`$. ### 5.4 Doubly-Constrained Colors As in the previous section, we define a two-component to be a subset of (variable,color) pairs such that each has two constraints, and any pair in the subset is connected to any other by a path of constraints. A two-component must have the form of a cycle of pairs, but it is possible for more than one pair in the cycle to involve the same variable. We distinguish two such types of components: a small two-component is one that involves only three pairs, while a large two-component involves four or more pairs. ###### Lemma 18 Suppose a reduced $`(\mathrm{4},\mathrm{2})`$-CSP instance includes a large two-component, and let $`ϵ\mathrm{0}.\mathrm{2}\mathrm{8}\mathrm{7}`$. Then the instance can be replaced by smaller instances with work factor at most $`\lambda (\mathrm{3},\mathrm{3},\mathrm{5})`$. Proof: We split into subcases: 1. Suppose the cycle passes through five consecutive distinct variables, say $`(v,R)`$, $`(w,R)`$, $`(x,R)`$, $`(y,R)`$, and $`(z,R)`$. We can assume that, if any of these five variables has four color choices, then this is true of one of the first four variables. Any coloring that does not use both $`(v,R)`$ and $`(y,R)`$ can be made to use at least one of the two colors $`(w,R)`$ or $`(x,R)`$ without violating any of the constraints. Therefore, we can divide into three subproblems: one in which we use $`(w,R)`$, eliminating three variables, one in which we use $`(x,R)`$, again eliminating three variables, and one in which we use both $`(v,R)`$ and $`(y,R)`$, eliminating all five variables. If all five variables have only three color choices, The work factor resulting from this subdivision is $`\lambda (\mathrm{3},\mathrm{3},\mathrm{5})`$. If some of the variables have four color choices, the work factor is at most $`\lambda (\mathrm{3}ϵ,\mathrm{4}ϵ,\mathrm{5}\mathrm{2}ϵ)`$, which is smaller for the given range of $`ϵ`$. 2. Suppose two colors three constraints apart on a cycle belong to the same variable; for instance, the sequence of colors may be $`(v,R)`$, $`(w,R)`$, $`(x,R)`$, $`(v,G)`$. Then any coloring can be made to use one of $`(w,R)`$ or $`(x,R)`$ without violating any constraints. If we form one subproblem in which we use $`(w,R)`$ and one in which we use $`(x,R)`$, we get work factor at most $`\lambda (\mathrm{3}ϵ,\mathrm{3}ϵ)`$ (the worst case occurring when only $`v`$ has four color choices). 3. Any long cycle which does not contain one of the previous two subcases must pass through the same four variables in the same order one, two, or three times. If it passes through two or three times, all four variables may be safely colored using colors from the cycle, reducing the problem with work factor one. And if the cycle has length exactly four, we may choose one of two ways to use two diagonally opposite colors from the cycle, giving work factor at most $`\lambda (\mathrm{4},\mathrm{4})`$. For the given range of $`ϵ`$, the largest of these work factors is $`\lambda (\mathrm{3},\mathrm{3},\mathrm{5})`$. ### 5.5 Matching Suppose we have a $`(\mathrm{4},\mathrm{2})`$-CSP instance to which none of the preceding reduction lemmas applies. Then, every constraint must be part of a good three-component or a small two-component. As we now show, this simple structure enables us to solve the remaining problem quickly. ###### Lemma 19 If we are given a $`(\mathrm{4},\mathrm{2})`$-CSP instance in which every constraint must be part of a good three-component or a small two-component, then we can solve it or determine that it is not solvable in polynomial time. Proof: We form a bipartite graph, in which the vertices correspond to the variables and components of the instance. We connect a variable to a component by an edge if there is a (variable,color) pair using that variable and belonging to that component. Since each pair in a good three-component or small two-component is connected by a constraint to every other pair in the component, any solution to the instance can use at most one (variable,color) pair per component. Thus, a solution consists of a set of (variable,color) pairs, covering each variable once, and covering each component at most once. In terms of the bipartite graph constructed above, this is simply a matching. So, we can solve the problem by using a graph maximum matching algorithm to determine the existence of a matching that covers all the variables. ### 5.6 Overall CSP Algorithm This completes the case analysis needed for our result. ###### Theorem 1 We can solve any $`(\mathrm{3},\mathrm{2})`$-CSP instance in time $`O(\lambda (\mathrm{4},\mathrm{4},\mathrm{5},\mathrm{5})^n)O(\mathrm{1}.\mathrm{3}\mathrm{6}\mathrm{4}\mathrm{4}\mathrm{3}^n)`$. Proof: We employ a backtracking (depth first) search in a state space consisting of $`(\mathrm{3},\mathrm{2})`$-CSP instances. At each point in the search, we examine the current state, and attempt to find a set of smaller instances to replace it with, using one of the reduction lemmas above. Such a replacement can always be found in polynomial time by searching for various simple local configurations in the instance. We then recursively search each smaller instance in succession. If we ever reach an instance in which Lemma 19 applies, we perform a matching algorithm to test whether it is solvable. If so, we find a solution and terminate the search. If not, we backtrack to the most recent branching point of the search and continue with the next alternative at that point. A bound of $`\lambda ^n`$ on the number of recursive calls in this search algorithm, where $`\lambda `$ is the maximum work factor occurring in our reduction lemmas, can be proven by induction on the size of an instance. The work within each call is polynomial and does not add appreciably to the overall time bound. To determine the maximum work factor, we need to set a value for the parameter $`ϵ`$. We used Mathematica to find a numerical value of $`ϵ`$ minimizing the maximum of the work factors involving $`ϵ`$, and found that for $`ϵ\mathrm{0}.\mathrm{0}\mathrm{9}\mathrm{5}\mathrm{5}\mathrm{4}\mathrm{3}`$ the work factor is $`\mathrm{1}.\mathrm{3}\mathrm{6}\mathrm{4}\mathrm{4}\mathrm{3}\lambda (\mathrm{4},\mathrm{4},\mathrm{5},\mathrm{5})`$. For $`ϵ`$ near this value, the two largest work factors are $`\lambda (\mathrm{3}ϵ,\mathrm{4}ϵ,\mathrm{4}ϵ)`$ (from Lemma 12) and $`\lambda (\mathrm{1}+ϵ,\mathrm{4})`$ (from Lemma 13); the remaining work factors are below 1.36. The true optimum value of $`ϵ`$ is thus the one for which $`\lambda (\mathrm{3}ϵ,\mathrm{4}ϵ,\mathrm{4}ϵ)=\lambda (\mathrm{1}+ϵ,\mathrm{4})`$. As we now show, for this optimum $`ϵ`$, $`\lambda (\mathrm{3}ϵ,\mathrm{4}ϵ,\mathrm{4}ϵ)=\lambda (\mathrm{1}+ϵ,\mathrm{4})=\lambda (\mathrm{4},\mathrm{4},\mathrm{5},\mathrm{5})`$, which also arises as a work factor in Lemma 17. Consider subdividing an instance of size $`n`$ into one of size $`n(\mathrm{1}+ϵ)`$ and another of size $`n\mathrm{4}`$, and then further subdividing the first instance into subinstances of size $`n(\mathrm{1}+ϵ)(\mathrm{3}ϵ)`$, $`n(\mathrm{1}+ϵ)(\mathrm{4}ϵ)`$, and $`n(\mathrm{1}+ϵ)(\mathrm{4}ϵ)`$. This four-way subdivision combines subdivisions of type $`\lambda (\mathrm{1}+ϵ,\mathrm{4})`$ and $`\lambda (\mathrm{3}ϵ,\mathrm{4}ϵ,\mathrm{4}ϵ)`$, so it must have a work factor between those two values. But by assumption those two values equal each other, so they also equal the work factor of the four-way subdivision, which is just $`\lambda (\mathrm{4},\mathrm{4},\mathrm{5},\mathrm{5})`$. We use the quantity $`\lambda (\mathrm{4},\mathrm{4},\mathrm{5},\mathrm{5})`$ frequently in the remainder of the paper, so we use $`\mathrm{\Lambda }`$ to denote this value. Theorem 1 immediately gives algorithms for some more well known problems, some of which we improve later. Of these, the least familiar is likely to be list $`k`$-coloring: given at each vertex of a graph a list of $`k`$ colors chosen from some larger set, find a coloring of the whole graph in which each vertex color is chosen from the corresponding list . ###### Corollary 2 We can solve the 3-coloring and 3-list coloring problems in time $`O(\mathrm{\Lambda }^n)`$, the 3-edge-coloring problem in time $`O(\mathrm{\Lambda }^m)`$, and the 3-SAT problem in time $`O(\mathrm{\Lambda }^t)`$, ###### Corollary 3 There is a randomized algorithm which finds the solution to any solvable $`(d,\mathrm{2})`$-CSP instance (with $`d>\mathrm{3}`$) in expected time $`O((\mathrm{0}.\mathrm{4}\mathrm{5}\mathrm{1}\mathrm{8}d)^n)`$. Proof: Randomly choose a subset of four values for each variable and apply our algorithm to the resulting $`(\mathrm{4},\mathrm{2})`$-CSP problem. Repeat with a new random choice until finding a solvable $`(\mathrm{4},\mathrm{2})`$-CSP instance. The random restriction of a variable has probability $`\mathrm{4}/d`$ of preserving solvability so the expected number of trials is $`(d/\mathrm{4})^n`$. Each trial takes time $`O(\mathrm{\Lambda }^{(\mathrm{2}ϵ)n})O(\mathrm{1}.\mathrm{8}\mathrm{0}\mathrm{7}\mathrm{2}^n)`$. The total expected time is therefore $`O((d/\mathrm{4})^n\mathrm{1}.\mathrm{8}\mathrm{0}\mathrm{7}\mathrm{2}^n)`$. ## 6 Vertex Coloring Simply by translating a 3-coloring problem into a $`(\mathrm{3},\mathrm{2})`$-CSP instance, as described above, we can test 3-colorability in time $`O(\mathrm{\Lambda }^n)`$. We now describe some methods to reduce this time bound even further. The basic idea is as follows: we find a small set of vertices $`SV(G)`$ with a large set $`N`$ of neighbors, and choose one of the $`\mathrm{3}^{\|S\|}`$ colorings for all vertices in $`S`$. For each such coloring, we translate the remaining problem to a $`(\mathrm{3},\mathrm{2})`$-CSP instance. The vertices in $`S`$ are already colored and need not be included in the $`(\mathrm{3},\mathrm{2})`$-CSP instance. The vertices in $`N`$ now have a colored neighbor, so for each such vertex at most two possible colors remain; therefore we can eliminate them from the $`(\mathrm{3},\mathrm{2})`$-CSP instance using Lemma 2. The remaining instance has $`k=\|V(G)SN\|`$ vertices, and can be solved in time $`O(\mathrm{\Lambda }^k)`$ by Theorem 1. The total time is thus $`O(\mathrm{3}^{\|S\|}\mathrm{\Lambda }^k)`$. By choosing $`S`$ appropriately we can make this quantity smaller than $`O(\mathrm{\Lambda }^n)`$. We can assume without loss of generality that all vertices in $`G`$ have degree three or more, since smaller degree vertices can be removed without changing 3-colorability. As a first cut at our algorithm, choose $`X`$ to be any set of vertices, no two adjacent or sharing a neighbor, and maximal with this property. Let $`Y`$ be the set of neighbors of $`X`$. We define a rooted forest $`F`$ covering $`G`$ as follows: let the roots of $`F`$ be the vertices in $`X`$, let each vertex in $`Y`$ be connected to its unique neighbor in $`X`$, and let each remaining vertex $`v`$ in $`G`$ be connected to some neighbor of $`v`$ in $`Y`$. (Such a neighbor must exist or $`v`$ could have been added to $`X`$). We let the set $`S`$ of vertices to be colored consist of all of $`X`$, together with each vertex in $`Y`$ having three or more children in $`F`$. We classify the subtrees of $`F`$ rooted at vertices in $`Y`$ as follows (Figure 15). If a vertex $`v`$ in $`Y`$ has no children, we call the subtree rooted at $`v`$ a club. If $`v`$ has one child, we call its subtree a stick. If it has two children, we call its subtree a fork. And if it has three or more children, we call its subtree a broom. We can now compute the total time of our algorithm by multiplying together a factor of $`\mathrm{3}`$ for each vertex in $`S`$ (that is, the roots of the trees of $`F`$ and of broom subtrees) and a factor of $`\mathrm{\Lambda }`$ for each leaf in a stick or fork. We define the cost of a vertex in a tree $`T`$ to be the product $`p`$ of such factors involving vertices of $`T`$, spread evenly among the vertices—if $`T`$ contains $`k`$ vertices the cost is $`p^{\mathrm{1}/k}`$. The total time of the algorithm will then be $`O(c^n)`$ where $`c`$ is the maximum cost of any vertex. It is not hard to show that this maximum is achieved in trees consisting of three forks (Figure 16), for which the cost is $`(\mathrm{3}(\mathrm{\Lambda })^\mathrm{6})^{\mathrm{1}/\mathrm{1}\mathrm{0}}\mathrm{1}.\mathrm{3}\mathrm{4}\mathrm{4}\mathrm{8}\mathrm{8}`$. Therefore we can three-color any graph in time $`O(\mathrm{1}.\mathrm{3}\mathrm{4}\mathrm{4}\mathrm{8}\mathrm{8}^n)`$. We can improve this somewhat with some more work. ### 6.1 Cycles of Degree-Three Vertices We begin by showing that we can assume that our graph has a special structure: the degree-three vertices do not form any cycles. For if they do form a cycle, we can remove it cheaply as follows. ###### Lemma 20 Let $`G`$ be a 3-coloring instance in which some cycle consists only of degree-three vertices. Then we can replace $`G`$ by smaller instances with work factor at most $`\lambda (\mathrm{5},\mathrm{6},\mathrm{7},\mathrm{8})\mathrm{1}.\mathrm{2}\mathrm{4}\mathrm{3}\mathrm{3}`$. Proof: Let the cycle $`C`$ consist of vertices $`v_\mathrm{1}`$, $`v_\mathrm{2}`$, $`\mathrm{}`$, $`v_k`$. We can assume without loss of generality that it has no chords, since otherwise we could find a shorter cycle in $`G`$; therefore each $`v_i`$ has a unique neighbor $`w_i`$ outside the cycle, although the $`w_i`$ need not be distinct from each other. Note that, if any $`w_i`$ and $`w_{i+\mathrm{1}}`$ are adjacent, then $`G`$ is 3-colorable iff $`GC`$ is; for, if we have a coloring of $`GC`$, then we can color $`C`$ by giving $`v_{i+\mathrm{1}}`$ the same color as $`w_i`$, and then proceeding to color the remaining cycle vertices in order $`v_{i+\mathrm{2}}`$, $`v_{i+\mathrm{3}}`$, $`\mathrm{}`$, $`v_k`$, $`v_\mathrm{1}`$, $`v_\mathrm{2}`$, $`\mathrm{}`$, $`v_i`$. Each successive vertex has only two previously-colored neighbors, so there remains at least one free color to use, until we return to $`v_i`$. When we color $`v_i`$, all three of its neighbors are colored, but two of them have the same color, so again there is a free color. As a consequence, if $`C`$ has even length, then $`G`$ is 3-colorable iff $`GC`$ is; for if some $`w_i`$ and $`w_{i+\mathrm{1}}`$ are given different colors, then the above argument colors $`C`$, while if all $`w_i`$ have the same color, then the other two colors can be used in alternation around $`C`$. The first remaining case is that $`k=\mathrm{3}`$ (Figure 17, left). Then we divide the problem into two smaller instances, by forcing $`w_\mathrm{1}`$ and $`w_\mathrm{2}`$ to have different colors in one instance (by adding an edge between them, Figure 17 top right) while forcing them to have the same color in the other instance (by collapsing the two vertices into a single supervertex, Figure 17 bottom right). If we add an edge between $`w_\mathrm{1}`$ and $`w_\mathrm{2}`$, we may remove $`C`$, reducing the problem size by three. If we give them the same color as each other, the instance is only colorable if $`v_\mathrm{3}`$ is also given the same color, so we can collapse $`v_\mathrm{3}`$ into the supervertex and remove the other two cycle vertices, reducing the problem size by four. Thus the work factor in this case is $`\lambda (\mathrm{3},\mathrm{4})\mathrm{1}.\mathrm{2}\mathrm{2}\mathrm{0}\mathrm{7}`$. If $`k`$ is odd and larger than three, we form three smaller instances, as shown in Figure 18. In the first, we add an edge between $`w_\mathrm{1}`$ and $`w_\mathrm{2}`$, and remove $`C`$, reducing the problem size by $`k`$. In the second, we collapse $`w_\mathrm{1}`$ and $`w_\mathrm{2}`$, add an edge between the new supervertex and $`w_\mathrm{3}`$, and again remove $`C`$, reducing the problem size by $`k+\mathrm{1}`$. In the third instance, we collapse $`w_\mathrm{1}`$, $`w_\mathrm{2}`$, and $`w_\mathrm{3}`$. This forces $`v_\mathrm{1}`$ and $`v_\mathrm{3}`$ to have the same color as each other, so we also collapse those two vertices into another supervertex and remove $`v_\mathrm{2}`$, reducing the problem size by four. For $`k\mathrm{7}`$ this gives work factor at most $`\lambda (\mathrm{4},\mathrm{7},\mathrm{8})\mathrm{1}.\mathrm{1}\mathrm{9}\mathrm{8}\mathrm{7}`$. For $`k=\mathrm{5}`$ the subproblem with $`n\mathrm{4}`$ vertices contains a triangle of degree-three vertices, and can be further subdivided into two subproblems of $`n\mathrm{7}`$ and $`n\mathrm{8}`$ vertices, giving the claimed work factor. Any degree-three vertices remaining after the application of this lemma must form components that are trees. As we now show, we can also limit the size of these trees. ###### Lemma 21 Let $`G`$ be a 3-coloring instance containing a connected subset of eight or more degree-three vertices. Then we can replace $`G`$ by smaller instances with work factor at most $`\lambda (\mathrm{2},\mathrm{5},\mathrm{6})\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{4}\mathrm{7}`$. Proof: Suppose the subset forms a $`k`$-vertex tree, and let $`v`$ be a vertex in this tree such that each subtree formed by removing $`v`$ has at most $`k/\mathrm{2}`$ vertices. Then, if $`G`$ is 3-colored, some two of the three neighbors of $`v`$ must be given the same color, so we can split the instance into three smaller instances, each of which collapses two of the three neighbors into a single supervertex. This collapse reduces the number of vertices by one, and allows the removal of $`v`$ (since after the collapse $`v`$ has degree two) and the subtree connected to the third vertex. Thus we achieve work factor $`\lambda (a,b,c)`$ where $`a+b+c=k+\mathrm{3}`$ and $`max\{a,b,c\}k/\mathrm{2}`$. The worst case is $`\lambda (\mathrm{2},\mathrm{5},\mathrm{6})`$, achieved when $`k=\mathrm{8}`$ and the tree is a path. ### 6.2 Planting Good Trees We define a bushy forest to be an unrooted forest within a given instance graph, such that each internal node has degree four or more (for an example, see the top three levels of Figure 21). A bushy forest is maximal if no internal node is adjacent to a vertex outside the forest, no leaf has three or more neighbors outside the forest, and no vertex outside the forest has four or more neighbors outside the forest. If a leaf $`v`$ does have three or more neighbors outside the forest, we could add those neighbors to the tree containing $`v`$, producing a bushy forest with more vertices. Similarly, if a vertex outside the forest has four or more neighbors outside the forest, we could extend the forest by adding another tree consisting of that vertex and its neighbors. As we now show, a maximal bushy forest must cover at least a constant fraction of a 3-coloring instance graph. ###### Lemma 22 Let $`G`$ be a graph in which all vertex degrees are three or more, and in which there is no cycle of degree-three vertices, let $`F`$ be a maximal bushy forest in $`G`$, and let $`r`$ denote the number of leaves in $`F`$. Then $`\|GF\|\mathrm{2}\mathrm{0}r/\mathrm{3}`$. Proof: Divide $`GF`$ into two subsets $`X`$ and $`Y`$, where $`X`$ consists of the vertices of degree four or more and $`Y`$ consists of the degree-three vertices. Let $`m_{A,B}`$ denote the number of edges connecting sets $`A`$ and $`B`$. Then each vertex in $`X`$ must have at least one edge connecting it to $`F`$, and at most three edges connecting it to $`Y`$, so $`m_{X,F}\|X\|`$ and $`m_{X,Y}\mathrm{3}\|X\|`$. Further, to avoid cycles, each connected component in $`Y`$ must form a tree, and if such a component has $`k`$ vertices, it must have $`k+\mathrm{2}`$ edges leaving it, and $`k\mathrm{7}`$ else we could apply Lemma 21. So, $`m_{Y,XF}\mathrm{9}\|Y\|/\mathrm{7}`$. If $`\mathrm{3}\|X\|\mathrm{9}\|Y\|/\mathrm{7}`$, $`m_{F,XY}=m_{F,X}+m_{F,XY}m_{X,Y}\mathrm{9}\|Y\|/\mathrm{7}\mathrm{2}\|X\|=\mathrm{3}\|XY\|/\mathrm{1}\mathrm{0}+(\mathrm{6}\mathrm{9}\|Y\|/\mathrm{7}\mathrm{0}\mathrm{2}\mathrm{3}\|X\|/\mathrm{1}\mathrm{0})\mathrm{3}\|XY\|/\mathrm{1}\mathrm{0}`$. And if $`\mathrm{3}\|X\|\mathrm{9}\|Y\|/\mathrm{7}`$, then again $`m_{F,XY}m_{F,X}\|X\|\mathrm{3}\|XY\|/\mathrm{1}\mathrm{0}`$. However, each leaf in $`F`$ has at most two edges outside $`F`$, or $`F`$ would not be maximal, so $`\|XY\|\mathrm{1}\mathrm{0}m_{F,XY}/\mathrm{3}\mathrm{2}\mathrm{0}r/\mathrm{3}`$. ### 6.3 Pruning Bad Trees After finding a maximal bushy forest $`F`$, we find a second forest $`H`$ in the remaining graph $`GF`$, as follows. Note that, due to the maximality of $`F`$, each vertex in $`GF`$ has at most three neighbors in $`GF`$. We first choose a maximal set $`T`$ of disjoint $`K_{\mathrm{1},\mathrm{3}}`$ subgraphs in $`GF`$. Then, we increase the size of $`T`$ as much as possible by operations in which we remove one $`K_{\mathrm{1},\mathrm{3}}`$ from $`T`$ and form two $`K_{\mathrm{1},\mathrm{3}}`$ subgraphs from the remaining vertices. Let $`X`$ denote the set of vertices in $`G(TF)`$ that are adjacent to vertices in $`F`$. By the maximality of $`F`$, each vertex in $`F`$ is adjacent to at most two vertices in $`X`$. Let $`Y=G(XTF)`$ denote the remaining vertices. By the maximality of $`T`$, each vertex in $`Y`$ is adjacent to at most two vertices in $`XY`$, and so must have a neighbor in $`T`$. Since $`GF`$ contains no degree-four vertices, each vertex in $`T`$ must have at most two neighbors in $`Y`$. As we now show, we can assign vertices in $`Y`$ to trees in $`T`$, extending each tree in $`T`$ to a tree of height at most two, in such a way that we do not form any tree with three forks, which would otherwise be the worst case for our algorithm. ###### Lemma 23 Let $`F`$, $`T`$, $`X`$, and $`Y`$ be as above. Then there exists a forest $`H`$ of height two trees with three branches each, such that the vertices of $`H`$ are exactly those of $`SY`$, such that each tree in $`H`$ has at most five grandchildren, and such that any tree with four or more grandchildren contains at least one vertex with degree four or more in $`G`$. Proof: We first show how to form a set $`H^{}`$ of non-disjoint trees in $`TY`$, and a set of weights on the grandchildren of these trees, such that each tree’s grandchildren have weight at most five. To do this, let each tree in $`H^{}`$ be formed by one of the $`K_{\mathrm{1},\mathrm{3}}`$ trees in $`T`$, together with all possible grandchildren in $`Y`$ that are adjacent to the $`K_{\mathrm{1},\mathrm{3}}`$ leaves. We assign each vertex in $`Y`$ unit weight, which we divide equally among the trees it belongs to. Then, suppose for a contradiction that some tree $`h`$ in $`H^{}`$ has grandchildren with total weight more than five. Then, its grandchildren must form three forks, and at least five of its six grandchildren must have unit weight; i.e., they belong only to tree $`h`$. Note that each vertex in $`Y`$ must have degree three, or we could have added it to the bushy forest, and all its neighbors must be in $`SY`$, or we could have added it to $`X`$. The unit weight grandchildren each have one neighbor in $`h`$ and two other neighbors in $`Y`$. These two other neighbors must be one each from the two other forks in $`h`$, for, if to the contrary some unit-weight grandchild $`v`$ does not have neighbors in both forks, we could have increased the number of trees in $`T`$ by removing $`h`$ and adding new trees rooted at $`v`$ and at the missed fork. Thus, these five grandchildren each connect to two other grandchildren, and (since no grandchild connects to three grandchildren) the six grandchildren together form a degree-two graph, that is, a union of cycles of degree-three vertices. But after applying Lemma 20 to $`G`$, it contains no such cycles. This contradiction implies that the weight of $`h`$ must be at most five. Similarly, if the weight of $`h`$ is more than three, it must have at least one fork, at least one unit-weight grandchild outside that fork, and at least one edge connecting that grandchild to a grandchild within the fork. This edge together with a path in $`h`$ forms a cycle, which must contain a high degree vertex. We are not quite done, because the assignment of grandchildren to trees in $`H^{}`$ is fractional and non-disjoint. To form the desired forest $`H`$, construct a network flow problem in which the flow source is connected to a node representing each tree $`tT`$ by an edge with capacity $`w(t)=\mathrm{5}`$ if $`t`$ contains a high degree vertex and capacity $`w(t)=\mathrm{3}`$ otherwise. The node corresponding to tree $`t`$ is connected by unit-capacity edges to nodes corresponding to the vertices in $`Y`$ that are adjacent to $`t`$, and each of these nodes is connected by a unit-capacity edge to a flow sink. Then the fractional weight system above defines a flow that saturates all edges into the flow sink and is therefore maximum (Figure 19, middle top). But any maximum flow problem with integer edge capacities has an integer solution (Figure 19, middle bottom). This solution must continue to saturate the sink edges, so each vertex in $`Y`$ will have one unit of flow to some tree $`t`$, and no flow to the other adjacent trees. Thus, the flow corresponds to an assignment of vertices in $`Y`$ to adjacent trees in $`T`$ such that each tree is assigned at most $`w(t)`$ vertices. We then simply let each tree in $`H`$ consist of a tree in $`T`$ together with its assigned vertices in $`Y`$ (Figure 19, bottom). ### 6.4 Improved Tree Coloring We now discuss how to color the trees in the height-two forest $`H`$ constructed in the previous subsection. As in the discussion at the start of this section, we color some vertices (typically just the root) of each tree in $`H`$, leave some vertices (typically the grandchildren) to be part of a later $`(\mathrm{3},\mathrm{2})`$-CSP instance, and average the costs over all the vertices in the tree. However, we average the costs in the following strange way: a cost of $`\mathrm{\Lambda }`$ is assigned to any vertex with degree four or higher in $`G`$, as if it was handled as part of the $`(\mathrm{3},\mathrm{2})`$-CSP instance. The remaining costs are then divided equally among the remaining vertices. ###### Lemma 24 Let $`T`$ be a tree with three children and at most five grandchildren. Then $`T`$ can be colored with cost per degree-three vertex at most $`(\mathrm{3}\mathrm{\Lambda }^\mathrm{3})^{\mathrm{1}/\mathrm{7}}\mathrm{1}.\mathrm{3}\mathrm{3}\mathrm{6}\mathrm{6}`$. Proof: First, suppose that $`T`$ has exactly five grandchildren. At least one vertex of $`T`$ has high degree. Two of the children $`x`$ and $`y`$ must be the roots of forks, while the third child $`z`$ is the root of a stick. We test each of the nine possible colorings of $`x`$ and $`y`$. In six of the cases, $`x`$ and $`y`$ are different, forcing the root to have one particular color (Figure 20, right). In these cases the only remaining vertex after translation to a $`(\mathrm{3},\mathrm{2})`$-CSP instance and application of Lemma 2 will be the child of $`z`$, so in each such case $`T`$ accumulates a further cost of $`\mathrm{\Lambda }`$. In the three cases in which $`x`$ and $`y`$ are colored the same (Figure 20, left), we must also take an additional factor of $`\mathrm{\Lambda }`$ for $`z`$ itself. One of these $`\mathrm{\Lambda }`$ factors goes to a high degree vertex, while the remaining work is split among the remaining eight vertices. The cost per vertex in this case is then at most $`(\mathrm{6}+\mathrm{3}\mathrm{\Lambda })^{\mathrm{1}/\mathrm{8}}\mathrm{1}.\mathrm{3}\mathrm{3}\mathrm{5}\mathrm{1}`$. If $`T`$ has fewer than five grandchildren, we choose a color for the root of the tree as described at the start of the section. The worst case occurs when the number of grandchildren is either three or four, and is $`(\mathrm{3}\mathrm{\Lambda }^\mathrm{3})^{\mathrm{1}/\mathrm{7}}\mathrm{1}.\mathrm{3}\mathrm{3}\mathrm{6}\mathrm{6}`$. ### 6.5 The Vertex Coloring Algorithm ###### Theorem 2 We can solve the 3-coloring problem in time $`O((\mathrm{2}^{\mathrm{3}/\mathrm{4}\mathrm{9}}\mathrm{3}^{\mathrm{4}/\mathrm{4}\mathrm{9}}\mathrm{\Lambda }^{\mathrm{2}\mathrm{4}/\mathrm{4}\mathrm{9}})^n)\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^n`$. Proof: As described in the preceding sections, we find a maximal bushy forest, then cover the remaining vertices by height-two trees. We choose colors for each internal vertex in the bushy forest, and for certain vertices in the height-two trees as described in Lemma 24. Vertices adjacent to these colored vertices are restricted to two colors, while the remaining vertices form a $`(\mathrm{3},\mathrm{2})`$-CSP instance and can be colored using our general $`(\mathrm{3},\mathrm{2})`$-CSP algorithm. Let $`p`$ denote the number of vertices that are roots in the bushy forest; $`q`$ denote the number of non-root internal vertices; $`r`$ denote the number of bushy forest leaves; $`s`$ denote the number of vertices adjacent to bushy forest leaves; and $`t`$ denote the number of remaining vertices, which must all be degree-three vertices in the height-two forest (Figure 21). Then the total time for the algorithm is at most $`\mathrm{3}^p\mathrm{2}^q\mathrm{\Lambda }^s(\mathrm{3}\mathrm{\Lambda }^\mathrm{3})^{t/\mathrm{7}}`$. We now consider which values of these parameters give the worst case for this time bound, subject to the constraints $`p,q,r,s,t\mathrm{0}`$, $`p+q+r+s+t=n`$, $`\mathrm{4}p+\mathrm{2}qr`$ (from the definition of a bushy forest), $`\mathrm{2}rs`$ (from the maximality of the forest), and $`\mathrm{2}\mathrm{0}r/\mathrm{3}s+t`$ (Lemma 22). We ignore the slightly tighter constraint $`p\mathrm{1}`$ since it only complicates the overall solution. Since the work per vertex in $`s`$ and $`t`$ is larger than that in the bushy forests, the time bound is maximized when $`s`$ and $`t`$ are as large as possible; that is, when $`s+t=\mathrm{2}\mathrm{0}r/\mathrm{3}`$. Further since the work per vertex in $`s`$ is larger than that in $`t`$, $`s`$ should be as large as possible; that is, $`s=\mathrm{2}r`$ and $`t=\mathrm{1}\mathrm{4}r/\mathrm{3}`$. Increasing $`p`$ or $`q`$ and correspondingly decreasing $`r`$, $`s`$, and $`t`$ only increases the time bound, since we pay a factor of 2 or more per vertex in $`p`$ and $`q`$ and at most $`\mathrm{\Lambda }`$ for the remaining vertices, so in the worst case the constraint $`\mathrm{4}p+\mathrm{2}qr`$ becomes an equality. It remains only to set the balance between parameters $`p`$ and $`q`$. There are two candidate solutions: one in which $`q=\mathrm{0}`$, so $`r=\mathrm{4}p`$, and one in which $`p=\mathrm{0}`$, so $`r=\mathrm{2}q`$. In the former case $`n=p+\mathrm{4}p+\mathrm{8}p+\mathrm{5}\mathrm{6}p/\mathrm{3}=\mathrm{9}\mathrm{5}p/\mathrm{3}`$ and the time bound is $`\mathrm{3}^p\mathrm{\Lambda }^{\mathrm{8}p}(\mathrm{3}\mathrm{\Lambda }^\mathrm{3})^{\mathrm{8}p/\mathrm{3}}=\mathrm{3}^{\mathrm{1}\mathrm{1}p/\mathrm{3}}\mathrm{\Lambda }^{\mathrm{1}\mathrm{6}p}\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{7}^n`$. In the latter case $`n=q+\mathrm{2}q+\mathrm{4}q+\mathrm{2}\mathrm{8}q/\mathrm{3}=\mathrm{4}\mathrm{9}q/\mathrm{3}`$ and the time bound is $`\mathrm{2}^q\mathrm{\Lambda }^{\mathrm{4}q}(\mathrm{3}\mathrm{\Lambda }^\mathrm{3})^{\mathrm{4}q/\mathrm{3}}=\mathrm{2}^q\mathrm{3}^{\mathrm{4}q/\mathrm{3}}\mathrm{\Lambda }^{\mathrm{8}q}\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^n`$. ## 7 Edge Coloring We now describe an algorithm for finding edge colorings of undirected graphs, using at most three colors, if such colorings exist. We can assume without loss of generality that the graph has vertex degree at most three. Then $`m\mathrm{3}n/\mathrm{2}`$, so by applying our vertex coloring algorithm to the line graph of $`G`$ we could achieve time bound $`\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^{\mathrm{3}n/\mathrm{2}}\mathrm{1}.\mathrm{5}\mathrm{3}\mathrm{1}\mathrm{9}^n`$. Just as we improved our vertex coloring algorithm by performing some reductions in the vertex coloring model before treating the problem as a $`(\mathrm{3},\mathrm{2})`$-CSP instance, we improve this edge coloring bound by performing some reductions in the edge coloring model before treating the problem as a vertex coloring instance. The main idea is to solve a problem intermediate in generality between 3-edge-coloring and 3-vertex-coloring: 3-edge-coloring with some added constraints that certain pairs of edges should not be the same color. ###### Lemma 25 Suppose a constrained 3-edge-coloring instance contains an unconstrained edge connecting two degree-three vertices. Then the instance can be replaced by two smaller instances with three fewer edges and two fewer vertices each. Proof: Let the given edge be $`(w,x)`$, and let its four neighbors be $`(u,w)`$, $`(v,w)`$, $`(x,y)`$, and $`(x,z)`$. Then $`(w,x)`$ can be colored only if its four neighbors together use two of the three colors, which forces these neighbors to be matched into equally colored pairs in one of two ways. Thus, we can replace the instance by two smaller instances: one in which we replace the five edges by the two edges $`(u,y)`$ and $`(v,z)`$, and one in which we replace the five edges by the two edges $`(u,z)`$ and $`(v,y)`$; in each case we add a constraint between the two new edges. The reduction operation described in Lemma 25 is depicted in Figure 22. We let $`m_\mathrm{3}`$ denote the number of edges with three neighbors in an unconstrained 3-edge-coloring instance, and $`m_\mathrm{4}`$ denote the number of edges with four neighbors. Edges with fewer neighbors can be removed at no cost, so we can assume without loss of generality that $`m=m_\mathrm{3}+m_\mathrm{4}`$. ###### Lemma 26 In an unconstrained 3-edge-coloring instance, we can find in polynomial time a set $`S`$ of $`m_\mathrm{4}/\mathrm{3}`$ edges such that Lemma 25 can be applied independently to each edge in $`S`$. Proof: Use a maximum matching algorithm in the graph induced by the edges with four neighbors. If the graph is 3-colorable, the resulting matching must contain at least $`m_\mathrm{4}/\mathrm{3}`$ edges. Applying Lemma 25 to an edge in a matching neither constrains any other edge in the matching, nor causes the remaining edges to stop being a matching. ###### Lemma 27 $`m_\mathrm{3}=\frac{\mathrm{6}}{\mathrm{5}}n\frac{\mathrm{4}}{\mathrm{5}}m_\mathrm{4}`$. Proof: Assign a charge of $`\mathrm{6}/\mathrm{5}`$ to each vertex of the graph, and redistribute this charge equally to each incident edge. Further assign an additional $`\mathrm{1}/\mathrm{5}`$ charge to each four-neighbor edge. Then each edge receives a unit charge, so $`m_\mathrm{3}+m_\mathrm{4}=m=(\mathrm{6}/\mathrm{5})n+(\mathrm{1}/\mathrm{5})m_\mathrm{4}`$. Subtracting $`m_\mathrm{4}`$ from both sides yields the result. ###### Theorem 3 We can 3-edge-color any 3-edge-colorable graph, in time $`O(\mathrm{2}^{n/\mathrm{2}})`$. Proof: We apply Lemma 26, resulting in a set of $`\mathrm{2}^{m_\mathrm{4}/\mathrm{3}}`$ constrained 3-edge-coloring problems each having only $`m_\mathrm{3}`$ edges. We then treat these remaining problems as 3-vertex-coloring problems on the corresponding line graphs, augmented by additional edges representing the constraints added by Lemma 25. The time for this algorithm is thus at most $`O(\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^{m_\mathrm{3}}\mathrm{2}^{m_\mathrm{4}/\mathrm{3}})`$. By Lemma 27, we can rewrite this bound as $`O(\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^{\mathrm{6}n/\mathrm{5}}(\mathrm{2}^{\mathrm{1}/\mathrm{3}}\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^{\mathrm{4}/\mathrm{5}})^{m_\mathrm{4}})`$. Since $`\mathrm{2}^{\mathrm{1}/\mathrm{3}}\mathrm{1}.\mathrm{3}\mathrm{2}\mathrm{8}\mathrm{9}^{\mathrm{4}/\mathrm{5}}>\mathrm{1}`$, this time bound is maximized when $`m_\mathrm{4}`$ is maximized, which occurs when $`m_\mathrm{4}=\mathrm{3}n/\mathrm{2}`$ and $`m_\mathrm{3}=\mathrm{0}`$. For this value, all the work occurs within Lemma 26, and gives the stated time bound. #### Acknowledgments A preliminary version of this paper was presented at the 36th IEEE Symp. Foundations of Comp. Sci., 1995. The first author thanks Russell Impagliazzo and Richard Lipton for bringing this problem to his attention. Both authors thank Laszlo Lovasz for helpful discussions.
warning/0006/nucl-th0006017.html	ar5iv	text	# Mixing of the 𝑓₀ and 𝑎₀ scalar mesons in threshold photoproduction ## I Introduction The $`f_0`$ and $`a_0`$ scalar mesons present a well-known puzzle for which several interesting, albeit controversial proposals have been made ranging from quark-antiquark makeup, to four-quark, and to $`\overline{K}K`$ molecular structure. Our goal is not to review these suggestions, but rather to investigate if new information on the makeup of these scalar mesons can be gleaned from threshold $`f_0(980)`$ and $`a_0(980)`$ photoproduction. Our hope is that threshold production of $`0^+`$ $`K\overline{K}`$ and two pion states that arise from the decay of scalar mesons will shed light on the nature of the $`f_0`$ and $`a_0`$ scalar mesons. For that purpose, we examine a potentially important final state isospin mixing effect. The $`f_0`$ has been observed in a photoproduction experiment at Fermilab . Measurements of exclusive $`f_0`$ photo/electroproduction are being performed at Jefferson Laboratory . Extensive theoretical studies of $`f_0`$ photoproduction have been presented in recent papers . Our study differs from these works in that we assume $`f_0`$ and $`a_0`$ mesons are produced directly and play the role of doorway states, while in these mesons enter via final state interaction in $`K\overline{K}`$ and $`\pi \pi `$ channels. The two descriptions are complementary and further studies are needed to clarify their interconnection. Another novel feature of our present work is the inclusion of possibly important final state $`f_0a_0`$ mixing. This mixing is induced by scalar meson transitions into and out of $`K^+K^{}`$ and $`K^0\overline{K}^0`$ intermediate states, which generates $`f_0a_0`$ mixing because of the 8 MeV splitting between the $`K^+K^{}`$ and $`K^0\overline{K}^0`$ thresholds. Thus scalar meson mixing arises from the difference between the neutral and charged kaon masses, e.g. from isospin breaking. This effect is fully included into our treatment. Calculations presented below rely on a very restricted number of parameters and most of which are fixed using recent experimental data on the $`\varphi \gamma f_0/a_0`$ reaction . We study the reaction $$\gamma +pp^{}+f_0/a_0p^{}+m_1m_2,$$ (1) using the $`f_0`$ and the $`a_0`$ as the possible doorway for subsequent decay meson production. Here $`m_1m_2`$ denotes $`\pi \pi `$ or $`K\overline{K}`$ pairs, or the $`\pi \eta `$ system. A thorough theoretical investigation of this reaction should include at least three essential points: a) the general structure of the invariant scalar meson photoproduction amplitude, including expressions for the cross sections and spin observables; b) effects of the final state interaction, including $`K\overline{K}`$ threshold phenomena; c) a dynamical model for the reaction mechanism, e.g., an effective Lagrangian and/or a set of leading diagrams. In our paper we concentrate on points (a) and (b), which are more universal and less model dependent than point (c) above. The structure of the basic production amplitude and spin observables presented below do not rely on any explicit reaction mechanisms (other than the doorway assumption) and are applicable to the photoproduction of any scalar meson. The final state interaction (FSI) form factor we present later is constructed in a model independent way based on a theory of resonance mixture. Only a few parameters are needed which on one hand are sensitive to the nature of the $`f_0/a_0`$ mesons, and on the other can be extracted from recent experimental data on the $`\phi \gamma f_0/a_0`$ reaction. Concerning the dynamical point (c), the dominant reaction mechanism is expected to depend upon the kinematical conditions of the explicit experiment. For example, at low photon momenta $`s`$–channel resonances might give a substantial contribution , while forward photoproduction at higher energies might be dominated by $`t`$–channel, $`\rho `$ and $`\omega `$ meson exchanges . An alternate, very promising approach making use of chiral Lagrangians was proposed recently . Although we do not propose an explicit reaction mechanism, our driving idea is complementary to that of Refs. . We take the $`f_0/a_0`$ to be “doorway states” for the reaction. Namely, we assume that $`f_0/a_0`$ scalar mesons are produced “directly” and then propagate under the strong influence of two nearby $`K\overline{K}`$ thresholds before decaying into $`K\overline{K},\pi \pi ,`$ or $`\pi \eta .`$ This propagator, or FSI form factor, has a simple form and is applicable to any reaction involving $`f_0/a_0.`$ We begin by analyzing the structure of the photoproduction of scalar mesons using this $`0^+`$ doorway model, then we formulate a description of the final state interaction. The general behavior of the cross-section as a function of the invariant mass of the kaon or pion pair is then generated to gauge the importance of the mixing near threshold and for comparison with experiment. Speculations about the role of scalar-meson nucleon P-waves are presented, especially as they relate to the possible observation of spin observables. ## II Diagrams and cross sections Before focusing on $`f_0/a_0`$ photoproduction in reaction ( 1), we comment on possible background effects. The amplitude for two-meson photoproduction $`\gamma pp^{}m_1m_2`$ can be represented symbolically by the series of graphs depicted in Fig. 1. Although all diagrams entering into the background term $`B_\alpha `$ <sup></sup><sup></sup>* The subscript $`\alpha `$ denotes the $`\pi \pi ,`$ or $`K\overline{K}`$ channels. (Fig. 2) might be explicitly calculated (some of them were calculated in ), our main point is that $`B_\alpha `$ is a smooth function of the energy in the vicinity of the $`K\overline{K}`$ threshold; whereas, dramatic energy dependent effects in the invariant mass distribution of the $`m_1m_2`$ pair (see below) are generated mainly by the first two diagrams in Fig. 1. We neglect the pion loop diagram shown in Fig. 3 (see )for two reasons. One reason is that near the $`K\overline{K}`$ threshold the pion loop by itself, i.e., without adding a resonant rescattering vertex, has smooth energy behavior and has to be included into the renormalized vertex of the $`f_0`$ production. The second reason for neglecting pion loops is based on Weinberg’s power counting theorem , which states that pion loops are suppressed compared to the tree diagrams depicted in Fig. 1. According to these arguments the pion loops are of the same order as higher terms in the effective Lagrangian. However, the above reasoning is far from being a rigorous proof of the direct production dominance. The theoretical and experimental study of this problem is at an early stage and we hope that comparison of complementary approaches will advance the investigation. Consider now the graphs on the top of Fig. 1. We first consider only the $`f_0`$-meson; $`a_0`$ production will be included later. Then these diagrams under consideration are described by the following equation $$T_\alpha =V_{\alpha f}G_fF_f+B_\alpha ,$$ (2) where $`F_f`$ is the $`f_0`$ photoproduction amplitude, $`G_f`$ is the $`f_0`$ propagator (also called the FSI form factor), and $`V_{\alpha f}`$ is the amplitude for $`f_0`$ decay into the final meson pair channel $`\alpha `$ ($`\alpha \pi \pi ,K\overline{K}`$), $`B_\alpha `$ denotes the background terms for channel $`\alpha `$ which we presume to be weakly energy dependent. As shown below, the form factor $`G_f`$ contains dynamical information on the nature of the $`f_0`$ meson. This form factor also reflects the interplay of the $`f_0`$ with the nearby $`K\overline{K}`$ threshold. The simple kinematical fact that $`K^+K^{}`$ and $`K^0\overline{K}^0`$ thresholds are split by 8 MeV induces significant $`f_0/a_0`$ mixing, as first noticed by Achasov et al. . Taking account of this $`f_0/a_0`$ mixing, Eq. (2) is now extended to include a sum over intermediate $`f_0`$ and $`a_0`$ states $$T_\alpha =\underset{i,k}{}V_{\alpha i}G_{ik}F_k+B_\alpha ,$$ (3) with $`i,k=f_0,a_0`$, i.e., the form factor $`G_{ik}`$ becomes a $`2\times 2`$ matrix in $`f_0/a_0`$ space. As pointed out earlier by Stodolsky measurement of the invariant mass distribution of the $`(m_1+m_2)`$ system in reaction ( 1) enables one to determine the nonorthogonality of the decaying states; i.e. the probability for $`f_0a_0`$ transitions. The cross section for the reaction $`\gamma pp^{}m_1m_2`$ (see Fig. 1), where the meson pair is in channel $`\alpha ,`$ is $$\sigma _\alpha =\frac{m^2}{(2\pi )^3}\frac{1}{\lambda ^{1/2}(s,0,m^2)}𝑑E_{m_1}𝑑E_{m_2}\frac{d\mathrm{\Omega }_{m_1}}{4\pi }\frac{d\phi _{m_2}}{2\pi }\|T_\alpha \|^2,$$ (4) where $`\lambda (x,y,z)(xyz)^24yz,s=(k+p)^2=(p^{}+q_{m_1}+q_{m_2})^2,m`$ is the proton mass, $`\mathrm{\Omega }_{m_1}`$ is the angle of decay meson 1 with respect to the incident photon beam $`\stackrel{}{k}`$ and $`\mathrm{\Omega }_{m_2}=(\mathrm{cos}\theta _{m_1m_2},\phi _{m_2})`$ is the angle of decay meson 2 with respect to $`\mathrm{\Omega }_{m_1}`$. The decay mesons are coplanar as described by $`\phi _{m_2}.`$ To observe the $`f_0/a_0`$ meson and to investigate its properties use is made of the double differential cross section $$\frac{d\sigma _\alpha }{dtdm_x}=\frac{m^2}{(2\pi )^3}\frac{\sqrt{\lambda (m_x^2,m_{m_1}^2,m_{m_2}^2)}}{4m_x\lambda (s,0,m^2)}\frac{1}{4\pi }𝑑\mathrm{\Omega }^{}\|T_\alpha \|^2,$$ (5) or of the invariant mass $`m_x`$ distribution $$\frac{d\sigma _\alpha }{dm_x}=\frac{m^2}{(2\pi )^3}\frac{\sqrt{\lambda (s,m_x^2,m^2)\lambda (m_x^2,m_{m_1}^2,m_{m_2}^2)}}{4sm_x(sm^2)}\frac{1}{2}_1^{+1}d\mathrm{cos}\theta \frac{1}{4\pi }𝑑\mathrm{\Omega }^{}\|T_\alpha \|^2,$$ (6) where $`t=(pp^{})^2,m_x=s_1^{1/2},s_1=(p_{m_1}+p_{m_2})^2,`$ and $`\mathrm{\Omega }^{}`$ is the angle between the relative momentum of the two decay mesons in their c.m. system. The angle $`\theta `$ refers to the scalar meson production angle. Proton spinors are normalized according to $`\overline{u}u=1.`$ We will use this last expression to display the cross-section as a function of the invariant mass $`m_x`$, once we have described the amplitude $`T_\alpha `$ of Eq. (3). The initial and final helicity summations are suppressed above in $`\|T_\alpha \|^2.`$ For unpolarized beam, target and recoil proton experiments, it is therefore understood that one sums over final and averages over initial helicities. However, if one has a polarized beam or polarized target, the associated cross sections can be expressed, using the bold assumption that $`B_\alpha `$ interference can be neglected, by: $$\|T_\alpha \|^2=V_{\alpha i}G_{ik}Tr[(F_k\rho _IF_k^{}^{})]G_{i^{}k^{}}^{}V_{\alpha i^{}}^{},$$ (7) where $`\rho _I`$ is the density matrix which describes the spin state of the initial beam and target. Since the decay $`V_{\alpha i}`$ and the propagator $`G_{ik}`$ are both independent of helicities, the inner term above $`𝒮_{kk^{}}=Tr[(F_k\rho _IF_k^{}^{})]`$ includes a trace over the helicity quantum numbers in the usual way. Note that $`𝒮_{kk^{}}^{}=𝒮_{k^{}k}`$ which assures real observables and that $`𝒮`$ is Hermitian in the $`f_0/a_0`$ channel space, which should not be confused with the Hermiticity of the density matrix in spin-space. The function $`𝒮_{kk^{}}`$ is a coupled-channels version of the usual ensemble average expression. For a polarized recoil proton experiment, $`𝒮_{kk^{}}`$ generalizes to $`𝒮_{kk^{}}=Tr[\rho _F(F_k\rho _IF_k^{}^{})],`$ where $`\rho _F`$ describes the final recoil polarization measurement. With this expression, it is possible to extract the effect of having a polarized beam and/or target on the cross-section. It is of interest to see if such measurements, especially as meson-nucleon P-waves appear, can shed light on the structure of the scalar mesons. We are now ready to examine the structure of scalar meson photoproduction in order to specify the scalar meson production amplitude $`F.`$ ## III Invariant Amplitudes and Spin Observables for the Scalar Meson The spin structure of the scalar meson photoproduction amplitude $`F`$ and the corresponding spin observables are discussed in this section. Again $`k`$, $`p`$ and $`p^{}`$ denote the 4-momenta of the photon and the initial and final protons, respectively; while the scalar meson’s 4-momentum is denoted by $`q.`$ The decomposition of the scalar meson photoproduction amplitude as invariant amplitudes has the following (CGLN )form $$F=\overline{u}(p^{})\underset{j=1}{\overset{4}{}}M_jA_j(s,t,u)u(p).$$ (8) We follow the $`\overline{u}u=1`$ normalization conditions of . For the scalar meson, the four invariant amplitudes $`M_j`$ are $$M_1=\epsilon \text{/}k\text{/},$$ $$M_2=2(\epsilon p)(kp^{})2(\epsilon p^{})(kp),$$ $$M_3=\epsilon \text{/}(kp)k\text{/}(\epsilon p),$$ (9) $$M_4=\epsilon \text{/}(kp^{})k\text{/}(\epsilon p^{}),$$ where $`\epsilon _\mu `$ is the photon polarization 4-vector. The four amplitudes $`M_i`$ differ from the $`0^{}`$ meson photoproduction only in the omission of the $`\gamma _5`$ pseudoscalar factor. The amplitude $`F`$ can be expressed in terms of the two–component spinors $`\chi `$. For this purpose we write $$F=\sqrt{\frac{\omega (p)\omega (p^{})}{4m^2}}<\chi (p^{})\|R\|\chi (p)>,$$ (10) where $`\omega (p)=E_p+m`$. Substitution of the invariant amplitudes Eq. (10) into Eq. (8) leads to the following CGLN-type representation for the amplitude $`R`$ $$R=i\stackrel{}{\epsilon }(\stackrel{}{\sigma }\times \widehat{k})R_1+(\stackrel{}{\epsilon }\stackrel{}{\sigma })(\stackrel{}{\sigma }\widehat{q})R_2+i(\stackrel{}{\epsilon }\widehat{q})(\widehat{q}\times \widehat{k})\stackrel{}{\sigma }R_3+(\stackrel{}{\epsilon }\widehat{q})R_4,$$ (11) where $`\widehat{k},\widehat{q}`$ are unit vectors. The four amplitudes $`R_j`$ are related to the four amplitudes $`A_j`$ of Eq. (8) via $$R_1=(\sqrt{s}m)\left\{A_1+\frac{1}{2}(\sqrt{s}m)A_3+\frac{(kp^{})}{\sqrt{s}+m}A_4\right\}$$ $$R_2=\frac{\|\stackrel{}{q}\|}{E_p^{}+m}\left\{(\sqrt{s}m)A_1+\frac{1}{2}(sm^2)A_3+(kp^{})A_4\right\}$$ (12) $$R_3=\|\stackrel{}{q}\|^2(\frac{\sqrt{s}m}{E_p^{}+m})\left\{(\sqrt{s}m)A_2A_4\right\}$$ $$R_4=(\sqrt{s}m)\|\stackrel{}{q}\|\{A_2(\sqrt{s}+m)[1\frac{\stackrel{}{k}\stackrel{}{q}}{(E_p+m)(E_p^{}+m)}]+$$ $$+A_4[1+(\frac{\sqrt{s}+m}{\sqrt{s}m})\frac{\stackrel{}{k}\stackrel{}{q}}{(E_p+m)(E_p^{}+m)}]\}.$$ The $`(kp^{})`$ factor denotes a 4-vector product. To recast (12) into a fully relativistic invariant form, the following simple kinematical equations may be used $$2(\stackrel{}{k}\stackrel{}{q})=t\mu ^2+\frac{1}{2s}(sm^2)(sm^2+\mu ^2),$$ $$2kp^{}=m^2u,$$ $$\|\stackrel{}{q}\|=\frac{1}{2\sqrt{s}}\lambda ^{1/2}(s,\mu ^2,m^2),$$ (13) $$E_p=\frac{1}{2\sqrt{s}}(s+m^2),E_p^{}=\frac{1}{2\sqrt{s}}(s+m^2\mu ^2),$$ where $`\mu `$ is the mass of either $`f_0`$ or $`a_0.`$ Note that an alternative form of the amplitude (12) may be found in Ref. . Note that at the scalar meson production threshold all of the $`R_i`$ amplitudes vanish except for $`R_1.`$ We conclude that at the scalar meson production threshold the operator $`F`$ is given by $$F=\frac{[\omega (p)\omega (p^{})]^{1/2}}{2m}R_1i\stackrel{}{\epsilon }(\stackrel{}{\sigma }\times \widehat{k}).$$ (14) The value of the coefficient $`R_1`$ is discussed below in Sec. 5. As meson-nucleon P-waves turn on, the associated $`J=\frac{1}{2}^{},J=\frac{3}{2}^{}`$ amplitudes contribute with an initial linear dependence on the momentum $`\|\stackrel{}{q}\|;`$ hence, it is likely that the terms $`R_2`$ and $`R_4`$ will appear along with their operators. That unfolding of P-waves has implications concerning which spin observables assume nonzero values as the energy rises beyond the threshold value of $`E_\gamma (lab)=1.5`$ GeV. ## IV The Form Factor of the scalar Mesons near the $`K\overline{K}`$ Threshold In this section, we consider the FSI form factors $`G_f`$(pure $`0^+`$ propagation) and $`G_{ik}`$(coupled $`0^+`$ propagation) introduced in (2) and (3). These form factors describe the propagator of an unstable particle in the energy region overlapping some of the decay thresholds ($`K^+K^{}`$ and $`K^0\overline{K}^0`$). We are mainly interested in the coupled-channel form factor $`G_{ik},`$ which includes $`f_0a_0`$ mixing, but to simplify the discussion, we start with the simple one-channel propagator $`G_f,`$ which describes the $`f_0`$-meson and its decay channels. Later we generalize that discussion to the coupled-channels $`G_{ik}`$ case. There exists an overwhelming number of approaches to describe an unstable particle and particle mixtures. We follow the general phenomenological approach developed by Stodolsky and Kobzarev, Nikolaev and Okun , and in a slightly modified form presented in the book of Terent’ev . This approach has already been applied to the $`f_0/a_0`$ system in . We now briefly outline the derivation of the basic equations for propagation of an unstable particle; see for details. First consider $`f_0`$ and its decay channels ($`a_0`$ will be incorporated later) and let us introduce a set of state vectors $`i>`$ describing the scalar meson. One of these states, often denoted as $`f>,`$ is a discrete state; continuum states are also included in the set $`i>.`$ The discrete state couples to the continuum states and thereby acquires a width and a shift in mass; e.g., it becomes an unstable state. Let the whole scalar meson system be described by the Hamiltonian $`H=H_0+V`$, where $`H_0`$ has a multichannel continuous spectrum $`H_0i>=E_ii>`$ and a discrete state $`H_0f>=E_ff>.`$ The interaction V is responsible for the transitions between the above channels, so that with V “turned on,” $`f>`$ becomes a resonance. Consider the transition amplitude $$A_{if}(t)<if;t>,$$ (15) where $`f;t>=\mathrm{exp}(i(H_0+V)t)f>`$. “Diagonal” transitions $`A_{ff}(t)`$ are also included in this definition. The amplitude $`A_{if}(t)`$ satisfies the equation $$i\frac{}{t}A_{if}(t)=E_iA_{if}(t)+\underset{j}{}V_{ij}A_{jf}(t),$$ (16) where $`V_{ij}=<iVj>,`$ and the initial condition is $`A_{if}(0)=\delta _{if}`$. The equivalent integral equation reads $$A_{if}(t)=e^{iE_it}A_{if}(0)i\underset{j}{}_0^t𝑑t^{}V_{ij}e^{iE_i(tt^{})}A_{jf}(t^{}).$$ (17) Next we introduce the propagator for the interacting scalar meson system in the energy representation $$G_{if}(E+i\delta )=_0^{\mathrm{}}𝑑te^{i(E+i\delta )t}A_{if}(t).$$ (18) Applying a time integration similar to Eq. (18) to all of Eq. (17) yields $$G_{if}(E+i\delta )=\frac{iA_{if}(0)}{EE_i+i\delta }+\underset{j}{}V_{ij}\frac{G_{jf}(E+i\delta )}{EE_i+i\delta }.$$ (19) For $`if`$ taking account of the initial condition, one gets $$G_{if}(E+i\delta )=V_{if}\frac{G_{ff}(E+i\delta )}{EE_i+i\delta }+\underset{jf}{}V_{ij}\frac{G_{jf}(E+i\delta )}{EE_i+i\delta },$$ (20) while for $`i=f`$ the equation has the form $$G_{ff}(E+i\delta )=\frac{i}{EE_f+i\delta }+V_{ff}\frac{G_{ff}(E+i\delta )}{EE_f+i\delta }+\underset{jf}{}V_{fj}\frac{G_{jf}(E+i\delta )}{EE_f+i\delta }.$$ (21) Now we return to (20) and assume that $`V`$ only connects $`\|f>`$ with its decay channels, while direct coupling between decay channels is absent, i.e., $`V_{ij}=0`$ in Eq. (20) (this assumption may be dropped without changing the results significantly see Ref. ). Then we take $`G_{if}`$ given by the left hand side of Eq. (20) with $`V_{ij}=0`$ on the right hand side, change the index $`i`$ in $`G_{if}`$ into $`j`$ and substitute this $`G_{jf}`$ into (21). Thus we obtain $$\left(EE_fV_{ff}\underset{jf}{}V_{fj}\frac{1}{EE_j+i\delta }V_{jf}+i\delta \right)G_{ff}(E+i\delta )=i.$$ (22) This is the result for a pure $`f_0`$ meson case. The physical meaning of this propagator is that there is a shift in energy of the interacting meson due to a self interaction, plus a complex contribution from transitions to intermediate continuum states. Some of the flux into the intermediate state returns to the discrete state $`\|f>,`$ some flows away, thereby creating a shift in the width as well as the real part of the energy. Note, we often designate the interacting propagator $`G_{ff}`$ simply as $`G_f.`$ We now examine the role of the scalar mesons coupling to pion and kaon pairs in the intermediate states. Some remarks are now in order. Since we are interested in $`K\overline{K}`$ channels with thresholds close to the $`f_0`$ mass the non-relativistic derivation and in particular nonrelativistic kinematics used above are justified; the relativistic generalization is straightforward. Keeping in mind that $`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$ thresholds are far away from the $`f_0`$ position, we rewrite Eq. (22) in terms of “renormalized” eigenvalues instead of the “bare” ones. By that we mean that $`E_f`$ is redefined to include both the self-energy contribution $`V_{ff}`$ and the part of the sum in Eq. (22) that arises from channels other than the $`K\overline{K}`$ channels. Hence, effects of the $`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$ are incorporated by redefining or renormalizing $`E_f.`$ In the vicinity of the $`K\overline{K}`$ threshold, the terms arising from intermediate pion pairs are almost energy independent. They shift $`E_f`$ by the following real and imaginary parts $$\underset{\pi \pi }{}V_{f\pi \pi }\frac{1}{EE_{\pi \pi }+i\delta }V_{\pi \pi f}=M_fi\frac{\mathrm{\Gamma }_f}{2}.$$ (23) The width $`\mathrm{\Gamma }_f`$ is thus generated by the imaginary part of the transitions to the pion pair continuum. For the $`a_0`$-meson propagator analogous real and imaginary energy shifts arise from $`\eta \pi `$ intermediate states. The sum in (23) implies both a sum over different channels ($`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$) and integration over the energies $`E_{\pi \pi }`$ in each channel. For large values of $`E_{\pi \pi }`$ the integral diverges and standard renormalization has to be performed. In the simplified approach presented here we may even argue that the integral is convergent (or cut off) due to the matrix elements $`V_{f\pi \pi }`$, $`V_{\pi \pi f}`$. We keep the same notation $`E_f`$ for the real part of the $`f_0`$ energy renormalized in the above sense; i.e., after $`V_{ff}`$ and $`M_f`$ are absorbed into it. Then the $`f_0`$ propagator takes the following simple form $$G_f(E+i\delta )=\frac{i}{EE_f+i\frac{\mathrm{\Gamma }f}{2}_{K\overline{K}}V_{fK\overline{K}}\frac{1}{EE_{K\overline{K}}+i\delta }V_{K\overline{K}f}},$$ (24) where summation is now only over $`K^+K^{}`$ and $`K^0\overline{K}^0`$ channels and the sum also implies integration over energy. Being interested in the $`K\overline{K}`$ near threshold energy region, one cannot neglect the splitting between $`K^+K^{}`$ and $`K^0\overline{K}^0`$ thresholds equal to 2 ($`m_{K^0}m_{K^\pm })`$ = 8 MeV. This mass difference induces isospin violating $`f_0a_0`$ mixing. Such “kinematical” isospin violation was carefully studied earlier using effective range theory . For the $`f_0a_0`$ system the effect was probably first pointed out in , where it was shown that mixing is enhanced, i.e., is of the order of $`[(m_{K^0}m_{K^\pm }]/m_{K^0}]^{\frac{1}{2}},`$ instead of $`[(m_{K^0}m_{K^\pm }]/m_{K^0}]`$ as might be naively expected. As we shall see, this enhancement also follows in our approach and motivates us to include this effect in scalar meson production. Isospin violating $`f_0a_0`$ mixing has also been studied earlier by T. Barnes. We will not repeat the derivation including the $`a_0`$ meson since it proceeds along the same lines. The propagator now becomes a $`2\times 2`$ matrix in the $`f_0a_0`$ basis. Instead of (22) and (24), one gets $$\underset{m=f,a}{}\left[(EE_n+i\frac{\mathrm{\Gamma }_n}{2})\delta _{nm}\underset{K\overline{K}}{}V_{nK\overline{K}}\frac{1}{EE_{K\overline{K}}+i\delta }V_{K\overline{K}m}\right]G_{mk}=i\delta _{nk},$$ (25) where $`n,k=f,a.`$ The $`a_0`$-meson entering into this equation has been also “renormalized,” this time the $`\pi \eta `$ channel plays the role of the $`\pi \pi `$ channel for $`f_0`$. Our next task is to present explicit expressions for the sums over $`K\overline{K}`$ entering into (24) and (25). This problem was considered in in detail including the Coulomb interaction in the $`K^+K^{}`$ channel. The Coulomb interaction results in spectacular effects including the formation of the $`K^+K^{}`$ atom, but the typical energy scale for these phenomena is of the order of few KeV which requires an experimental resolution that is probably inaccessible in the near future. Neglecting the Coulomb interaction, the above sums are reduced to expressions of the type $$\underset{K\overline{K}}{}V_{nK\overline{K}}Y(m_x,m_K)V_{K\overline{K}m},$$ where we replaced $`E`$ by $`m_x`$ to be consistent with expressions (5-6) for the cross sections; here, $`Y(m_x,m_K)`$ are integrals of the type $$Y(m_x,m_K)=\frac{1}{2m_x^3}\frac{d^3p}{(2\pi )^3}\frac{1}{m_x2m_K\frac{p^2}{m_K}+i\delta },$$ (26) with $`m_K=m_{K^\pm },m_{K^0}`$. As already mentioned, the integral (26) is formally divergent which is however neither dangerous or important as soon as we focus on the $`m_x`$ region close to the $`K\overline{K}`$ threshold. Close to the $`K\overline{K}`$ threshold the leading energy dependent contribution is $$Y(m_x,m_K)=i\frac{1}{32\pi m_x}\sqrt{\frac{m_x2m_K}{m_K}+i0},$$ (27) where $`+i0`$ means that for $`m_x<2m_K`$ the square root acquires a positive imaginary part. From ( 25) and ( 27), it follows that for $`m_x>2m_K`$ the decay width into the $`K\overline{K}`$ channel is given by $`(n=f,a)`$ $$\mathrm{\Gamma }_{nK\overline{K}}=\frac{V_{nK\overline{K}}^2}{16\pi m_x}\sqrt{\frac{m_x2m_K}{m_K}}.$$ (28) Comparison of ( 28) with the parameterization used in shows that $`V_{nK\overline{K}}=g_{nK\overline{K}},`$ where $`g_{nK\overline{K}}`$ are the coupling constants used in . Recall that Eq. (25) is written in the basis of $`f_0`$ and $`a_0`$ states having definite isospins, $`I_{f_0}=1,I_{a_0}=0.`$ Therefore, $$V_{K^+K^{}f}=V_{K^0\overline{K}^0f}V_{K^+K^{}a}=V_{K^0\overline{K}^0a}.$$ (29) Next we introduce the notation $$𝒟\frac{V_{K^+K^{}f}^2}{32\pi m_x},\zeta \frac{V_{K^+K^{}a}}{V_{K^+K^{}f}},$$ (30) where $`\zeta `$ can be complex with phase $`\varphi ,`$ see later. In this notation, the explicit form of Eq. (25) can be expressed as a matrix in $`f_0a_0`$ space as: $`i\widehat{G}^1=\left(\begin{array}{cc}m_xE_f+i\frac{\mathrm{\Gamma }_f}{2}& 0\\ 0& m_xE_a+i\frac{\mathrm{\Gamma }_a}{2}\end{array}\right)`$ (33) $`+i𝒟(\begin{array}{cc}1& \zeta \\ \zeta ^{}& \|\zeta \|^2\end{array}))\sqrt{{\displaystyle \frac{m_x2m_{K^+}}{m_{K^+}}}+i0}`$ (36) $`+i𝒟\left(\begin{array}{cc}1& \zeta \\ \zeta ^{}& \|\zeta \|^2\end{array}\right)\sqrt{{\displaystyle \frac{m_x2m_{K^0}}{m_{K^0}}}+i0},`$ (39) Here we see that the isospin violating $`f_0a_0`$ mixing is really proportional to $`\sqrt{\frac{m_x2m_K}{m_K}}`$ and is indeed enhanced as was pointed out in Ref. . We also see that far beyond the region $`\|m_x2m_K\|8MeV`$ of the $`K\overline{K}`$ thresholds, the nondiagonal elements cancel each other and thereby extinguish the $`f_0a_0`$ mixing. Away from the threshold region one should also take into account corrections to ( 27), i.e., include the next terms in the expression for the renormalized kaon loop. Now we have at our disposal all quantities needed to calculate the near threshold amplitudes and cross sections using equations presented in Sec. 2. For example, the amplitude for the process $`\gamma +pp^{}+f_0p^{}+\pi \pi `$ reads $$T_{\pi \pi }=V_{\pi \pi f}G_{ff}F_f+V_{\pi \pi f}G_{fa}F_a+B_{\pi \pi }.$$ (40) Here we see the doorway production of the $`f_0`$, followed by its subsequent propagation via $`G_{ff}`$ and then its decay $`V_{f\pi \pi }.`$ The second term includes the doorway production of the $`a_0,`$ followed by its transition into a $`f_0`$ via isospin violation and then the $`f_0`$ decays into a pion pair. Other pion pair processes are included in the background term $`B_{\pi \pi }.`$ For final kaon production the corresponding result is: $$T_{K^+K^{}}=V_{K^+K^{}f}G_{ff}F_f+V_{K^+K^{}a}G_{aa}F_a+V_{K^+K^{}f}G_{fa}F_a+V_{K^+K^{}a}G_{af}F_f+B_{K^+K^{}}.$$ (41) For final $`\pi \eta `$ production the corresponding result is: $$T_{\pi \eta }=V_{\pi \eta a}G_{aa}F_a+V_{\pi \eta a}G_{af}F_f+B_{\pi \eta }.$$ (42) The values of the parameters are discussed in the next section. ## V Parameters of the Model The theoretical investigation of the light scalar mesons $`f_0`$ and $`a_0`$ has a long history. A most thoughtful study of $`f_0/a_0`$ has been performed by the Novosibirsk group (see and references therein). Although our approach differs from other authors, including the Novosibirsk group, the key parameters are similar in various treatments. Thus, in order to fix the parameters we shall relate them to those used by the Novosibirsk group. This would still leave the set of parameters incomplete, namely the photoproduction amplitude $`F_k`$ and the background term $`B_\alpha `$ ( see Eq.3) have no analogues in the Novosibirsk set and to get at least an educated guess about them, we resort to the recent paper . We stress that the characteristics of $`f_0/a_0`$ are at present far from being perfectly established. For example, the Novosibirsk group presents several solutions with some of the parameters ranging over rather wide limits . In order to determine the parameters of the model we use the recent experimental data presented in . These data do not allow us to fix all the parameters and some parameters will be set from different consideration (see below). The model under consideration can be described by several equivalent sets of parameters. We used the following set: $`E_f,E_a,\mathrm{\Gamma }_f,\mathrm{\Gamma }_a,V_{fK^+K^{}}`$, and $`\zeta `$. The quantities $`V_{fK^0\overline{K}^0}`$, $`V_{aK^0\overline{K}^0}`$ and $`𝒟`$ are then expressed according to ( 29) and ( 30), while $`V_{f\pi \pi }`$ and $`V_{a\pi \eta }`$ are determined via $`\mathrm{\Gamma }_f`$ and $`\mathrm{\Gamma }_a`$ according to the equation $$\mathrm{\Gamma }_{jm_1m_2}=\frac{\|V_{jm_1m_2}\|^2}{16\pi m_x}\rho _{m_1m_2},$$ (43) where $`j=f,a`$, and $`\rho _{m_1m_2}`$ is the two-body phase space. At this point we remind the reader that since we treat the $`K\overline{K}`$ channels explicitly, their contribution does not enter into $`E_j`$ and $`\mathrm{\Gamma }_j(j=f,a)`$. We also mention that the “visible” widths of $`f_0/a_0`$ are smaller than the widths $`\mathrm{\Gamma }_f`$ and $`\mathrm{\Gamma }_a`$ . Finally, we recall that our parameters $`V_{jm_1m_2}`$ are equivalent to $`g_{jm_1m_2}`$ used by the Novosibirsk group . In we find the following parameters relevant for our purposes (experimental errors are omitted): $$\frac{\|V_{f\pi \pi }\|^2}{4\pi }=0.44GeV^2,R=\frac{\|V_{fK\overline{K}}\|^2}{\|V_{f\pi \pi }\|^2}=3.77.$$ (44) Using ( 43) and isotopic relations $$\|V_{f\pi ^+\pi ^{}}\|^2=\frac{2}{3}\|V_{f\pi \pi }\|^2,\|V_{fK^+K^{}}\|^2=\frac{1}{2}\|V_{fK\overline{K}}\|^2,$$ (45) we get from ( 44): $$\mathrm{\Gamma }_f=0.1GeV,V_{fK^+K^{}}=3.23GeV.$$ (46) Next setting in ( 39) $`m_x`$ equal to the physical mass of the $`f_0`$-meson, i. e. $`m_x=0.98`$ GeV, we find $`E_f=0.947`$ GeV. Unfortunately, we can not rely on well-established experimental data for the parameters of the $`a_0`$ meson. Therefore we assume $`\mathrm{\Gamma }_a\mathrm{\Gamma }_f`$ and $`E_aE_f`$ (these assumptions are in line with the solutions proposed in .) From $`\mathrm{\Gamma }_a`$ we get $`V_{a\pi \eta }=3.08`$ GeV. Finally we set $`\zeta =1,`$ which is true in both four-quark and molecular models of $`f_0/a_0`$ . Thus we arrive at the following parameter set (set A): $$E_f=E_a=0.947GeV,\mathrm{\Gamma }_a=\mathrm{\Gamma }_f=0.1GeV,V_{fK^+K^{}}=3.23GeV,\zeta =1.$$ (47) Note that the value $`\|V_{fK^+K^{}}\|^2/4\pi =0.83`$ GeV<sup>2</sup> lies in between the values typical for four-quark and molecular models of $`f_0`$-meson . In order to display the effect of $`f_0a_0`$ mixing, we consider also set B whose only difference from set A is that in B the splitting of the $`K^0\overline{K}^0`$ and $`K^+K^{}`$ thresholds is ignored and the kaon mass is taken equal to the average value $`m_K=\frac{1}{2}(m_{K^0}+m_{K^\pm })`$ and hence mixing is thereby turned off. This results in cancellation of the non-diagonal terms in Eq. 39. Within our approach we are not in a position to determine the value of the photoproduction amplitude $`F_j(j=f,a).`$ Its invariant structure is given by Eqs. 9 and 11, but the magnitudes of the invariant amplitudes remain unknown. In a recent paper the $`f_0/a_0`$ photoproduction has been considered in the near-threshold region, namely at $`k=1.7GeV.`$ Close to threshold the dominant contribution stems from the amplitude $`R_1`$ of Eq. 11. Its value can be considered a free parameter, but we prefer to take for $`R_1`$ the value proposed in This threshold value for $`R_1`$ is basically the pseudoscalar photoproduction Kroll-Ruderman limit scaled by a factor $`k/(2M)`$ due to the switch to a scalar meson The reasons for adopting this value for $`R_1`$ is on one hand because it enables a direct comparison of our results with that of and on the other hand, as already argued in Sec. 2, we may consider the pion loops included into the renormalized vertex of $`f_0`$ photoproduction while in the approach of the two pion photoproduction is the driving mechanism that determines the value of $`R_1.`$ The calculation or even estimate of the background term $`B_\alpha `$ in 2 is beyond the scope of this article. Some of the diagrams entering into $`B_\alpha `$ are depicted in Fig. 2 and in general may be calculated. In Ref. the background for $`d\sigma _{\pi ^+\pi ^{}}/dm_x`$ (see 6) was estimated to be around $`55\mu b/GeV.`$ In Fig. 4, we present the invariant mass distribution, $`d\sigma (K^+K^{})/dm_x,`$ for $`K^+K^{}`$ pair production calculated according to Eqs. 2.5, 4.17.and 4.19. In order to display the role of isospin breaking $`f_0a_0`$ mixing, we plot on the same figure the curve obtained with set B parameters, i.e., with $`f_0a_0`$ mixing switched off (which is achieved by using the same mass for the charged and neutral kaons). The two curves are different, in line with “enhanced” mixing in the sense described above. (The scalar meson mixing can be as large as a 70 % enhancement.) The structure in the invariant mass distribution for $`K^+K^{}`$ pair production in this near threshold region arises from the isospin breaking difference in mass between charged and neutral kaons. The $`a_0f_0`$ interference pattern turns out to be very sensitive to the phase of the parameter $`\zeta `$ (see (4.16)). The sensitivity to that phase is clearly displayed in Fig. 4. One should however keep in mind that in both $`q\overline{q}`$ and $`q^2\overline{q}^2`$ constituent quark models this parameter is predicted to be real . In Fig.5, the $`\pi ^+\pi ^{}`$ invariant mass distribution $`d\sigma (\pi ^+\pi ^{})/dm_x`$ is presented. We now use Eqs. 2.5, 4.17.and 4.18. Here the deviation from a simple Breit-Wigner resonance structure due to the influence of the $`K\overline{K}`$ thresholds is clearly seen. We do not insist on the absolute values of the cross sections in Figs. 4 and 5, because of backgound effects, but we do believe the change due to mixing is of some predictive value. ## VI Conclusions The analysis presented herein illustrates that threshold photoproduction can provide insight into the nature of the $`f_0`$ and $`a_0`$ scalar mesons. Valuable information may be obtained from the $`K\overline{K}`$ mass distribution in the threshold region if measured with a few MeV resolution. In our approach, this distribution has been described in a model-independent way, once the doorway idea is invoked. Most parameters can be deduced from $`\varphi \gamma (a_0+f_0)K\overline{K}`$ experiments. The invariant amplitude formalism for the photoproduction of scalar mesons presented in Sec. III, enables us to consider higher energy regions which are accessible at JLAB. Going to higher energies increases the number of CGLN amplitude parameters by involving all of the $`R_i`$ terms in (3.4). Also, with increasing energy more partial waves come into play giving rise to nonzero spin observables. The role of spin observables and how they evolve with increasing energy will be discussed in a future paper. ###### Acknowledgements. We wish to thank Dr. S. Bashinsky and Profs. W. Kloet, J. Thompson, S. Eidelman, and E.P. Solodov for their helpful suggestions. We also thank Professor R. Schumacher for information about scalar mesons. B.K. wishes to express appreciation for warm hospitality and financial support during his visit to the University of Pittsburgh and for support by RFFI grants 97-02-16406 and 00-02-17836, as well as from the U.S. National Science Foundation.
warning/0006/cond-mat0006010.html	ar5iv	text	# Statistical Mechanics of Recurrent Neural Networks I. Statics ## 1 Introduction Statistical mechanics deals with large systems of stochastically interacting microscopic elements (particles, atomic magnets, polymers, etc.). The strategy of statistical mechanics is to abandon any ambition to solve models of such systems at the microscopic level of individual elements, but to use the microscopic laws to calculate equations describing the behaviour of a suitably chosen set of macroscopic observables. The toolbox of statistical mechanics consists of methods to perform this reduction from the microscopic to a macroscopic level, which are all based on efficient ways to do the bookkeeping of probabilities. The experience and intuition that has been built up over the last century tells us what to expect, and serves as a guide in finding the macroscopic observables and in seeing the difference between relevant mathematical subtleties and irrelevant ones. As in any statistical theory, clean and transparent mathematical laws can be expected to emerge only for large (preferably infinitely large) systems. In this limit one often encounters phase transitions, i.e. drastic changes in the system’s macroscopic behaviour at specific values of global control parameters. Recurrent neural networks, i.e. neural networks with synaptic feedback loops, appear to meet the criteria for statistical mechanics to apply, provided we indeed restrict ourselves to large systems. Here the microscopic stochastic dynamical variables are the firing states of the neurons or their membrane potentials, and one is mostly interested in quantities such as average state correlations and global information processing quality, which are indeed measured by macroscopic observables. In contrast to layered networks, one cannot simply write down the values of successive neuron states for models of recurrent neural networks; here they must be solved from (mostly stochastic) coupled dynamic equations. Under special conditions (‘detailed balance’), which usually translate into the requirement of synaptic symmetry, the stochastic process of evolving neuron states leads towards an equilibrium situation where the microscopic state probabilities are known, and where the techniques of equilibrium statistical mechanics can be applied in one form or another. The equilibrium distribution found, however, will not always be of the conventional Boltzmann form. For non-symmetric networks, where the asymptotic (stationary) statistics are not known, dynamical techniques from non-equilibrium statistical mechanics are the only tools available for analysis. The ‘natural’ set of macroscopic quantities (or ‘order parameters’) to be calculated can be defined in practice as the smallest set which will obey closed deterministic equations in the limit of an infinitely large network. Being high-dimensional non-linear systems with extensive feedback, the dynamics of recurrent neural networks are generally dominated by a wealth of attractors (fixed-point attractors, limit-cycles, or even more exotic types), and the practical use of recurrent neural networks (in both biology and engineering) lies in the potential for creation and manipulation of these attractors through adaptation of the network parameters (synapses and thresholds). Input fed into a recurrent neural network usually serves to induce a specific initial configuration (or firing pattern) of the neurons, which serves as a cue, and the ‘output’ is given by the (static or dynamic) attractor which has been triggered by this cue. The most familiar types of recurrent neural network models, where the idea of creating and manipulating attractors has been worked out and applied explicitly, are the so-called attractor neural networks for associative memory, designed to store and retrieve information in the form of neuronal firing patterns and/or sequences of neuronal firing patterns. Each pattern to be stored is represented as a microscopic state vector. One then constructs synapses and thresholds such that the dominant attractors of the network are precisely the pattern vectors (in the case of static recall), or where, alternatively, they are trajectories in which the patterns are successively generated microscopic system states. From an initial configuration (the ‘cue’, or input pattern to be recognised) the system is allowed to evolve in time autonomously, and the final state (or trajectory) reached can be interpreted as the pattern (or pattern sequence) recognized by network from the input (see figure 1). For such programmes to work one clearly needs recurrent neural networks with extensive ‘ergodicity breaking’: the state vector will during the course of the dynamics (at least on finite time-scales) have to be confined to a restricted region of state space (an ‘ergodic component’), the location of which is to depend strongly on the initial conditions. Hence our interest will mainly be in systems with many attractors. This, in turn, has implications at a theoretical/mathematical level: solving models of recurrent neural networks with extensively many attractors requires advanced tools from disordered systems theory, such as replica theory (statics) and generating functional analysis (dynamics). It will turn out that a crucial issue is whether or not the synapses are symmetric. Firstly, synaptic asymmetry is found to rule out microscopic equilibrium, which has implications for the mathematical techniques which are available: studying models of recurrent networks with non-symmetric synapses requires solving the dynamics, even if one is only interested in the stationary state. Secondly, the degree of synaptic asymmetry turns out to be a deciding factor in determining to what extent the dynamics will be glassy, i.e. extremely slow and non-trivial, close to saturation (where one has an extensive number of attractors). In this paper (on statics) and its sequel (on dynamics) I will discuss only the statistical mechanical analysis of neuronal firing processes in recurrent networks with static synapses, i.e. network operation as opposed to network learning. I will also restrict myself to networks with either full or randomly diluted connectivity, the area in which the main progress has been made during the last few decades. Apart from these restrictions, the text aims to be reasonably comprehensive and self-contained. Even within the confined area of the operation of recurrent neural networks a truly impressive amount has been achieved, and many of the constraints on mathematical models which were once thought to be essential for retaining solvability but which were regrettable from a biological point of view (such as synaptic symmetry, binary neuron states, instantaneous neuronal communication, a small number of attractors, etc.) have by now been removed with success. At the beginning of the new millennium we know much more about the dynamics and statics of recurrent neural networks than ever before. I aim to cover in a more or less unified manner the most important models and techniques which have been launched over the years, ranging from simple symmetric and non-symmetric networks with only a finite number of attractors, to the more complicated ones with an extensive number, and I will explain in detail the techniques which have been designed and used to solve them. In the present paper I will first discuss and solve various members of the simplest class of models: those where all synapses are the same. Then I turn to the Hopfield model, which is the archetypical model to describe the functioning of symmetric neural networks as associative memories (away from saturation, where the number of attractors is finite), and to a coupled oscillator model storing phase patterns (again away from saturation). Next I will discuss a model with Gaussian synapses, where the number of attractors diverges, in order to introduce the so-called replica method, followed by a section on the solution of the Hopfield model near saturation. I close this paper with a guide to further references and an assessment of the past and future deliverables of the equilibrium statistical mechanical analysis of recurrent neural networks. ## 2 Definitions & Properties of Microscopic Laws In this section I define the most common microscopic models for recurrent neural networks, I show how one can derive the corresponding descriptions of the stochastic evolution in terms of evolving state probabilities, and I discuss some fundamental statistical mechanical properties. ### 2.1 Stochastic Dynamics of Neuronal Firing States Microscopic Definitions for Binary Neurons. The simplest non-trivial definition of a recurrent neural network is that where $`N`$ binary neurons $`\sigma _i\{1,1\}`$ (in which the states ‘1’ and ‘-1’ represent firing and rest, respectively) respond iteratively and synchronously to post-synaptic potentials (or local fields) $`h_i(𝝈)`$, with $`𝝈=(\sigma _1,\mathrm{},\sigma _N)`$. The fields are assumed to depend linearly on the instantaneous neuron states: $$\mathrm{𝑃𝑎𝑟𝑎𝑙𝑙𝑒𝑙}:\sigma _i(\mathrm{}+1)=\mathrm{sgn}[h_i(𝝈(\mathrm{}))+T\eta _i(\mathrm{})]h_i(𝝈)=\underset{j}{}J_{ij}\sigma _j+\theta _i$$ (1) The stochasticity is in the independent random numbers $`\eta _i(\mathrm{})\mathrm{}`$ (representing threshold noise), which are all drawn according to some distribution $`w(\eta )`$. The parameter $`T`$ is introduced to control the amount of noise. For $`T=0`$ the process (1) is deterministic: $`\sigma _i(\mathrm{}+1)=\mathrm{sgn}[h_i(𝝈(\mathrm{}))]`$. The opposite extreme is choosing $`T=\mathrm{}`$, here the system evolution is fully random. The external fields $`\theta _i`$ represent neural thresholds and/or external stimuli, $`J_{ij}`$ represents the synaptic efficacy at the junction $`ji`$ ($`J_{ij}>0`$ implies excitation, $`J_{ij}<0`$ inhibition). Alternatively we could decide that at each iteration step $`\mathrm{}`$ only a single randomly drawn neuron $`\sigma _i_{\mathrm{}}`$ is to undergo an update of the type (1): $$\mathrm{𝑆𝑒𝑞𝑢𝑒𝑛𝑡𝑖𝑎𝑙}:\begin{array}{cc}ii_{\mathrm{}}:\hfill & \sigma _i(\mathrm{}+1)=\sigma _i(\mathrm{})\hfill \\ i=i_{\mathrm{}}:\hfill & \sigma _i(\mathrm{}+1)=\mathrm{sgn}\left[h_i(𝝈(\mathrm{}))+T\eta _i(\mathrm{})\right]\text{}\hfill \end{array}$$ (2) with the local fields as in (1). The stochasticity is now both in the independent random numbers $`\eta _i(\mathrm{})`$ (the threshold noise) and in the site $`i_{\mathrm{}}`$ to be updated, drawn randomly from the set $`\{1,\mathrm{},N\}`$. For simplicity we assume $`w(\eta )=w(\eta )`$, and define $$g[z]=2_0^zd\eta w(\eta ):g[z]=g[z],\underset{z\pm \mathrm{}}{lim}g[z]=\pm 1,\frac{d}{dz}g[z]0$$ Popular choices for the threshold noise distributions are $$w(\eta )=(2\pi )^{\frac{1}{2}}e^{\frac{1}{2}\eta ^2}:g[z]=\mathrm{Erf}[z/\sqrt{2}],w(\eta )=\frac{1}{2}[1\mathrm{tanh}^2(\eta )]:g[z]=\mathrm{tanh}(z)$$ From Stochastic Equations to Evolving Probabilities. From the microscopic equations (1,2), which are suitable for numerical simulations, we can derive an equivalent but mathematically more convenient description in terms of microscopic state probabilities $`p_{\mathrm{}}(𝝈)`$. Equations (1,2) state that, if the system state $`𝝈(\mathrm{})`$ is given, a neuron $`i`$ to be updated will obey $$\mathrm{Prob}\left[\sigma _i(\mathrm{}+1)\right]=\frac{1}{2}\left[1+\sigma _i(\mathrm{}+1)g[\beta h_i(𝝈(\mathrm{}))]\right]$$ (3) with $`\beta =T^1`$. In the case (1) this rule applies to all neurons, and thus we simply get $`p_{\mathrm{}+1}(𝝈)=_{i=1}^N\frac{1}{2}\left[1+\sigma _ig[\beta h_i(𝝈(\mathrm{}))]\right]`$. If, on the other hand, instead of $`𝝈(\mathrm{})`$ only the probability distribution $`p_{\mathrm{}}(𝝈)`$ is given, this expression for $`p_{\mathrm{}+1}(𝝈)`$ is to be averaged over the possible states at time $`\mathrm{}`$: $$\mathrm{𝑃𝑎𝑟𝑎𝑙𝑙𝑒𝑙}:p_{\mathrm{}+1}(𝝈)=\underset{𝝈^{}}{}W[𝝈;𝝈^{}]p_{\mathrm{}}(𝝈^{})W[𝝈;𝝈^{}]=\underset{i=1}{\overset{N}{}}\frac{1}{2}[1+\sigma _ig[\beta h_i(𝝈^{})]]$$ (4) This is the standard representation of a Markov chain. Also the sequential process (2) can be formulated in terms of probabilities, but here expression (3) applies only to the randomly drawn candidate $`i_{\mathrm{}}`$. After averaging over all possible realisations of the sites $`i_{\mathrm{}}`$ we obtain: $$p_{\mathrm{}+1}(𝝈)=\frac{1}{N}\underset{i}{}\left\{[\underset{ji}{}\delta _{\sigma _j,\sigma _j(\mathrm{})}]\frac{1}{2}\left[1+\sigma _ig[\beta h_i(𝝈(\mathrm{}))]\right]\right\}$$ (with the Kronecker symbol: $`\delta _{ij}=1`$ if $`i=j`$, $`\delta _{ij}=0`$ otherwise). If, instead of $`𝝈(\mathrm{})`$, the probabilities $`p_{\mathrm{}}(𝝈)`$ are given, this expression is to be averaged over the possible states at time $`\mathrm{}`$, with the result: $$p_{\mathrm{}+1}(𝝈)=\frac{1}{N}\underset{i}{}\frac{1}{2}\left[1+\sigma _ig[\beta h_i(𝝈)]\right]p_{\mathrm{}}(𝝈)+\frac{1}{N}\underset{i}{}\frac{1}{2}\left[1+\sigma _ig[\beta h_i(F_i𝝈)]\right]p_{\mathrm{}}(F_i𝝈)$$ with the state-flip operators $`F_i\mathrm{\Phi }(𝝈)=\mathrm{\Phi }(\sigma _1,\mathrm{},\sigma _{i1},\sigma _i,\sigma _{i+1},\mathrm{},\sigma _N)`$. This equation can again be written in the standard form $`p_{\mathrm{}+1}(𝝈)=_𝝈^{}W[𝝈;𝝈^{}]p_{\mathrm{}}(𝝈^{})`$, but now with the transition matrix $$\mathrm{𝑆𝑒𝑞𝑢𝑒𝑛𝑡𝑖𝑎𝑙}:W[𝝈;𝝈^{}]=\delta _{𝝈,𝝈^{}}+\frac{1}{N}\underset{i}{}\{w_i(F_i𝝈)\delta _{𝝈,F_i𝝈^{}}w_i(𝝈)\delta _{𝝈,𝝈^{}}\}$$ (5) where $`\delta _{𝝈,𝝈^{}}=_i\delta _{\sigma _i,\sigma _i^{}}`$ and $$w_i(𝝈)=\frac{1}{2}\left[1\sigma _i\mathrm{tanh}\left[\beta h_i(𝝈)\right]\right]$$ (6) Note that, as soon as $`T>0`$, the two transition matrices $`W[𝝈;𝝈^{}]`$ in (4,5) both describe ergodic systems: from any initial state $`𝝈^{}`$ one can reach any final state $`𝝈`$ with nonzero probability in a finite number of steps (being one in the parallel case, and $`N`$ in the sequential case). It now follows from the standard theory of stochastic processes (see e.g. ) that in both cases the system evolves towards a unique stationary distribution $`p_{\mathrm{}}(𝝈)`$, where all probabilities $`p_{\mathrm{}}(𝝈)`$ are non-zero. From Discrete to Continuous Times. The above processes have the (mathematically and biologically) less appealing property that time is measured in discrete units. For the sequential case we will now assume that the duration of each of the iteration steps is a continuous random number (for parallel dynamics this would make little sense, since all updates would still be made in full synchrony). The statistics of the durations are described by a function $`\pi _{\mathrm{}}(t)`$, defined as the probability that at time $`t`$ precisely $`\mathrm{}`$ updates have been made. Upon denoting the previous discrete-time probabilities as $`\widehat{p}_{\mathrm{}}(𝝈)`$, our new process (which now includes the randomness in step duration) will be described by $$p_t(𝝈)=\underset{\mathrm{}0}{}\pi _{\mathrm{}}(t)\widehat{p}_{\mathrm{}}(𝝈)=\underset{\mathrm{}0}{}\pi _{\mathrm{}}(t)\underset{𝝈^{}}{}W^{\mathrm{}}[𝝈;𝝈^{}]p_0(𝝈^{})$$ and time has become a continuous variable. For $`\pi _{\mathrm{}}(t)`$ we make the Poisson choice $`\pi _{\mathrm{}}(t)=\frac{1}{\mathrm{}!}(\frac{t}{\mathrm{\Delta }})^{\mathrm{}}e^{t/\mathrm{\Delta }}`$. From $`\mathrm{}_\pi =t/\mathrm{\Delta }`$ and $`\mathrm{}^2_\pi =t/\mathrm{\Delta }+t^2/\mathrm{\Delta }^2`$ it follows that $`\mathrm{\Delta }`$ is the average duration of an iteration step, and that the relative deviation in $`\mathrm{}`$ at a given $`t`$ vanishes for $`\mathrm{\Delta }0`$ as $`\sqrt{\mathrm{}^2_\pi \mathrm{}_\pi ^2}/\mathrm{}_\pi =\sqrt{\mathrm{\Delta }/t}`$. The nice properties of the Poisson distribution under temporal derivation allow us to derive: $$\mathrm{\Delta }\frac{d}{dt}p_t(𝝈)=\underset{𝝈^{}}{}W[𝝈;𝝈^{}]p_t(𝝈^{})p_t(𝝈)$$ For sequential dynamics we choose $`\mathrm{\Delta }=\frac{1}{N}`$ so that, as in the parallel case, in one time unit each neuron will on average be updated once. The master equation corresponding to (5) acquires the form $$\frac{d}{dt}p_t(𝝈)=\underset{i}{}\left\{w_i(F_i𝝈)p_t(F_i𝝈)w_i(𝝈)p_t(𝝈)\right\}$$ (7) The $`w_i(𝝈)`$ (6) now play the role of transition rates. The choice $`\mathrm{\Delta }=\frac{1}{N}`$ implies $`\sqrt{\mathrm{}^2_\pi \mathrm{}_\pi ^2}/\mathrm{}_\pi =\sqrt{1/Nt}`$, so we will still for $`N\mathrm{}`$ no longer have uncertainty in where we are on the $`t`$ axis. Microscopic Definitions for Continuous Neurons. Alternatively, we could start with continuous neuronal variables $`\sigma _i`$ (representing e.g. firing frequencies or oscillator phases), where $`i=1,\mathrm{},N`$, and with stochastic equations of the form $$\sigma _i(t+\mathrm{\Delta })=\sigma _i(t)+\mathrm{\Delta }f_i(𝝈(t))+\sqrt{2T\mathrm{\Delta }}\xi _i(t)$$ (8) Here we have introduced (as yet unspecified) deterministic state-dependent forces $`f_i(𝝈)`$, and uncorrelated Gaussian distributed random forces $`\xi _i(t)`$ (the noise), with $`\xi _i(t)=0`$ and $`\xi _i(t)\xi _j(t^{})=\delta _{ij}\delta _{t,t^{}}`$. As before, the parameter $`T`$ controls the amount of noise in the system, ranging from $`T=0`$ (deterministic dynamics) to $`T=\mathrm{}`$ (completely random dynamics). If we take the limit $`\mathrm{\Delta }0`$ in (8) we find a Langevin equation (with a continuous time variable): $$\frac{d}{dt}\sigma _i(t)=f_i(𝝈(t))+\eta _i(t)$$ (9) This equation acquires its meaning only as the limit $`\mathrm{\Delta }0`$ of (8). The moments of the new noise variables $`\eta _i(t)=\xi _i(t)\sqrt{2T/\mathrm{\Delta }}`$ in (9) are given by $`\eta _i(t)=0`$ and $`\eta _i(t)\eta _j(t^{})=2T\delta _{ij}\delta (tt^{})`$. This can be derived from the moments of the $`\xi _i(t)`$. For instance: $$\eta _i(t)\eta _j(t^{})=\underset{\mathrm{\Delta }0}{lim}\frac{2T}{\mathrm{\Delta }}\xi _i(t)\xi _j(t^{})=2T\delta _{ij}\underset{\mathrm{\Delta }0}{lim}\frac{1}{\mathrm{\Delta }}\delta _{t,t^{}}=2TC\delta _{ij}\delta (tt^{})$$ The constant $`C`$ is found by summing over $`t^{}`$, before taking the limit $`\mathrm{\Delta }0`$, in the above equation: $$𝑑t^{}\eta _i(t)\eta _j(t^{})=\underset{\mathrm{\Delta }0}{lim}2T\underset{t^{}=\mathrm{}}{\overset{\mathrm{}}{}}\xi _i(t)\xi _j(t^{})=2T\delta _{ij}\underset{\mathrm{\Delta }0}{lim}\underset{t^{}=\mathrm{}}{\overset{\mathrm{}}{}}\delta _{t,t^{}}=2T\delta _{ij}$$ Thus $`C=1`$, which indeed implies $`\eta _i(t)\eta _j(t^{})=2T\delta _{ij}\delta (tt^{})`$. More directly, one can also calculate the moment generating function $$e^{i{\scriptscriptstyle 𝑑t_i\psi _i(t)\eta _i(t)}}=\underset{\mathrm{\Delta }0}{lim}\underset{i,t}{}\frac{dz}{\sqrt{2\pi }}e^{\frac{1}{2}z^2+iz\psi _i(t)\sqrt{2T\mathrm{\Delta }}}=\underset{\mathrm{\Delta }0}{lim}\underset{i,t}{}e^{T\mathrm{\Delta }\psi _i^2(t)}=e^{T{\scriptscriptstyle 𝑑t_i\psi _i^2(t)}}$$ (10) From Stochastic Equations to Evolving Probabilities. A mathematically more convenient description of the process (9) is provided by the Fokker-Planck equation for the microscopic state probability density $`p_t(𝝈)=\delta [𝝈𝝈(t)]`$, which we will now derive. For the discrete-time process (8) we expand the $`\delta `$-distribution in the definition of $`p_{t+\mathrm{\Delta }}(𝝈)`$ (in a distributional sense): $$p_{t+\mathrm{\Delta }}(𝝈)p_t(𝝈)=\delta \left[𝝈𝝈(t)\mathrm{\Delta }𝒇(𝝈(t))\sqrt{2T\mathrm{\Delta }}𝝃(t)\right]\delta [𝝈𝝈(t)]$$ $$=\underset{i}{}\frac{}{\sigma _i}\delta [𝝈𝝈(t)]\left[\mathrm{\Delta }f_i(𝝈(t))+\sqrt{2T\mathrm{\Delta }}\xi _i(t)\right]+T\mathrm{\Delta }\underset{ij}{}\frac{^2}{\sigma _i\sigma _j}\delta [𝝈𝝈(t)]\xi _i(t)\xi _j(t)+𝒪(\mathrm{\Delta }^{\frac{3}{2}})$$ The variables $`𝝈(t)`$ depend only on noise variables $`\xi _j(t^{})`$ with $`t^{}<t`$, so that for any function $`A`$: $`A[𝝈(t)]\xi _i(t)=A[𝝈(t)]\xi _i(t)=0`$, and $`A[𝝈(t)]\xi _i(t)\xi _j(t)=\delta _{ij}A[𝝈(t)]`$. As a consequence: $$\frac{1}{\mathrm{\Delta }}\left[p_{t+\mathrm{\Delta }}(𝝈)p_t(𝝈)\right]=\underset{i}{}\frac{}{\sigma _i}\delta [𝝈𝝈(t)]f_i(𝝈(t))+T\underset{i}{}\frac{^2}{\sigma _i^2}\delta [𝝈𝝈(t)]+𝒪(\mathrm{\Delta }^{\frac{1}{2}})$$ $$=\underset{i}{}\frac{}{\sigma _i}\left[p_t(𝝈)f_i(𝝈)\right]+T\underset{i}{}\frac{^2}{\sigma _i^2}p_t(𝝈)+𝒪(\mathrm{\Delta }^{\frac{1}{2}})$$ By taking the limit $`\mathrm{\Delta }0`$ we then arrive at the Fokker-Planck equation: $$\frac{d}{dt}p_t(𝝈)=\underset{i}{}\frac{}{\sigma _i}\left[p_t(𝝈)f_i(𝝈)\right]+T\underset{i}{}\frac{^2}{\sigma _i^2}p_t(𝝈)$$ (11) Examples: Graded Response Neurons and Coupled Oscillators. In the case of graded response neurons the continuous variable $`\sigma _i`$ represents the membrane potential of neuron $`i`$, and (in their simplest form) the deterministic forces are given by $`f_i(𝝈)=_jJ_{ij}\mathrm{tanh}[\gamma \sigma _j]\sigma _i+\theta _i`$, with $`\gamma >0`$ and with the $`\theta _i`$ representing injected currents. Conventional notation is restored by putting $`\sigma _iu_i`$. Thus equation (9) specialises to $$\frac{d}{dt}u_i(t)=\underset{j}{}J_{ij}\mathrm{tanh}[\gamma u_j(t)]u_i(t)+\theta _i+\eta _i(t)$$ (12) One often chooses $`T=0`$ (i.e. $`\eta _i(t)=0`$), the rationale being that threshold noise is already assumed to have been incorporated via the non-linearity in (12). In our second example the variables $`\sigma _i`$ represent the phases of coupled neural oscillators, with forces of the form $`f_i(𝝈)=_jJ_{ij}\mathrm{sin}(\sigma _j\sigma _i)+\omega _i`$. Individual synapses $`J_{ij}`$ now try to enforce either pair-wise synchronisation ($`J_{ij}>0`$) or pair-wise anti-synchronisation ($`J_{ij}<0`$), and the $`\omega _i`$ represent the natural frequencies of the individual oscillators. Conventional notation dictates $`\sigma _i\varphi _i`$, giving $$\frac{d}{dt}\varphi _i(t)=\omega _i+\underset{j}{}J_{ij}\mathrm{sin}[\varphi _j(t)\varphi _i(t)]+\eta _i(t)$$ (13) ### 2.2 Synaptic Symmetry & Lyapunov Functions Noise-free Symmetric Networks of Binary Neurons. In the deterministic limit $`T0`$ the rules (1) for networks of synchronously evolving binary neurons reduce to the deterministic map $$\sigma _i(\mathrm{}+1)=\mathrm{sgn}\left[h_i(𝝈(\mathrm{}))\right]$$ (14) It turns out that for systems with symmetric interactions, $`J_{ij}=J_{ji}`$ for all $`(ij)`$, one can construct a Lyapunov function, i.e. a function of $`𝝈`$ which during the dynamics decreases monotonically and is bounded from below (see e.g. ): $$\mathrm{𝐵𝑖𝑛𝑎𝑟𝑦}\&\mathrm{𝑃𝑎𝑟𝑎𝑙𝑙𝑒𝑙}:L[𝝈]=\underset{i}{}\|h_i(𝝈)\|\underset{i}{}\sigma _i\theta _i$$ (15) Clearly $`L_i[_j\|J_{ij}\|+\|\theta _i\|]_i\|\theta _i\|`$. During iteration of (14) we find: $$L[𝝈(\mathrm{}+1)]L[𝝈(\mathrm{})]=\underset{i}{}\|h_i(𝝈(\mathrm{}+1))\|+\underset{i}{}\sigma _i(\mathrm{}+1)[\underset{j}{}J_{ij}\sigma _j(\mathrm{})+\theta _i]\underset{i}{}\theta _i\left[\sigma _i(\mathrm{}+1)\sigma _i(\mathrm{})\right]$$ $$=\underset{i}{}\|h_i(𝝈(\mathrm{}+1))\|+\underset{i}{}\sigma _i(\mathrm{})h_i(𝝈(\mathrm{}+1))=\underset{i}{}\|h_i(𝝈(\mathrm{}+1))\|\left[1\sigma _i(\mathrm{}+2)\sigma _i(\mathrm{})\right]0$$ (where we used (14) and $`J_{ij}=J_{ji}`$). So $`L`$ decreases monotonically until a stage is reached where $`\sigma _i(\mathrm{}+2)=\sigma _i(\mathrm{})`$ for all $`i`$. Thus, with symmetric interactions this system will in the deterministic limit always end up in a limit cycle with period $`2`$. A similar result is found for networks with binary neurons and sequential dynamics. In the limit $`T0`$ the rules (2) reduce to the map $$\sigma _i(\mathrm{}+1)=\delta _{i,i_{\mathrm{}}}\mathrm{sgn}\left[h_i(𝝈(\mathrm{}))\right]+[1\delta _{i,i_{\mathrm{}}}]\sigma _i(\mathrm{})$$ (16) (in which we still have randomness in the choice of site to be updated). For systems with symmetric interactions and without self-interactions, i.e. $`J_{ii}=0`$ for all $`i`$, we again find a Lyapunov function: $$\mathrm{𝐵𝑖𝑛𝑎𝑟𝑦}\&\mathrm{𝑆𝑒𝑞𝑢𝑒𝑛𝑡𝑖𝑎𝑙}:L[𝝈]=\frac{1}{2}\underset{ij}{}\sigma _iJ_{ij}\sigma _j\underset{i}{}\sigma _i\theta _i$$ (17) This quantity is bounded from below: $`L\frac{1}{2}_{ij}\|J_{ij}\|_i\|\theta _i\|`$. Upon calling the site $`i_{\mathrm{}}`$ selected for update at step $`\mathrm{}`$ simply $`i`$, the change in $`L`$ during iteration of (16) can be written as: $$L[𝝈(\mathrm{}+1)]L[𝝈(\mathrm{})]=\theta _i[\sigma _i(\mathrm{}+1)\sigma _i(\mathrm{})]\frac{1}{2}\underset{k}{}J_{ik}[\sigma _i(\mathrm{}+1)\sigma _k(\mathrm{}+1)\sigma _i(\mathrm{})\sigma _k(\mathrm{})]$$ $$\frac{1}{2}\underset{j}{}J_{ji}[\sigma _j(\mathrm{}+1)\sigma _i(\mathrm{}+1)\sigma _j(\mathrm{})\sigma _i(\mathrm{})]$$ $$=[\sigma _i(\mathrm{})\sigma _i(\mathrm{}+1)][\underset{j}{}J_{ij}\sigma _j(\mathrm{})+\theta _i]=\|h_i(𝝈(\mathrm{}))\|\left[1\sigma _i(\mathrm{})\sigma _i(\mathrm{}+1)\right]0$$ Here we used (16), $`J_{ij}=J_{ji}`$, and absence of self-interactions. Thus $`L`$ decreases monotonically until $`\sigma _i(t+1)=\sigma _i(t)`$ for all $`i`$. With symmetric synapses, but without diagonal terms, the sequentially evolving binary neurons system will in the deterministic limit always end up in a stationary state. Noise-free Symmetric Networks of Continuous Neurons. One can derive similar results for models with continuous variables. Firstly, in the deterministic limit the graded response equations (12) simplify to $$\frac{d}{dt}u_i(t)=\underset{j}{}J_{ij}\mathrm{tanh}[\gamma u_j(t)]u_i(t)+\theta _i$$ (18) Symmetric networks again admit a Lyapunov function (there is no need to eliminate self-interactions): $$\mathrm{𝐺𝑟𝑎𝑑𝑒𝑑}\mathrm{𝑅𝑒𝑠𝑝𝑜𝑛𝑠𝑒}:L[𝒖]=\frac{1}{2}\underset{ij}{}J_{ij}\mathrm{tanh}[\gamma u_i]\mathrm{tanh}[\gamma u_j]+\underset{i}{}[\gamma _0^{u_i}dvv[1\mathrm{tanh}^2[\gamma v]]\theta _i\mathrm{tanh}[\gamma u_i]]$$ Clearly $`L\frac{1}{2}_{ij}\|J_{ij}\|_i\|\theta _i\|`$ (the term in $`L[𝒖]`$ with the integral is non-negative). During the noise-free dynamics (18) one can use the identity $`L/u_i=\gamma [1\mathrm{tanh}^2[\gamma u_i]](du_i/dt)`$, valid only when $`J_{ij}=J_{ji}`$, to derive $$\frac{d}{dt}L=\underset{i}{}\frac{L}{u_i}\frac{du_i}{dt}=\gamma \underset{i}{}[1\mathrm{tanh}^2[\gamma u_i]][\frac{d}{dt}u_i]^20$$ Again $`L`$ is found to decrease monotonically, until $`du_i/dt=0`$ for all $`i`$, i.e. until we are at a fixed-point. Finally, the coupled oscillator equations (13) reduce in the noise-free limit to $$\frac{d}{dt}\varphi _i(t)=\omega _i+\underset{j}{}J_{ij}\mathrm{sin}[\varphi _j(t)\varphi _i(t)]$$ (19) Note that self-interactions $`J_{ii}`$ always drop out automatically. For symmetric oscillator networks, a construction of the type followed for the graded response equations would lead us to propose $$\mathrm{𝐶𝑜𝑢𝑝𝑙𝑒𝑑}\mathrm{𝑂𝑠𝑐𝑖𝑙𝑙𝑎𝑡𝑜𝑟𝑠}:L[\mathit{\varphi }]=\frac{1}{2}\underset{ij}{}J_{ij}\mathrm{cos}[\varphi _i\varphi _j]\underset{i}{}\omega _i\varphi _i$$ (20) This function indeed decreases monotonically, due to $`L/\varphi _i=d\varphi _i/dt`$ : $$\frac{d}{dt}L=\underset{i}{}\frac{L}{\varphi _i}\frac{d\varphi _i}{dt}=\underset{i}{}[\frac{d}{dt}\varphi _i]^20$$ In fact (19) describes gradient descent on the surface $`L[\mathit{\varphi }]`$. However, due to the term with the natural frequencies $`\omega _i`$ the function $`L[\mathit{\varphi }]`$ is not bounded, so it cannot be a Lyapunov function. This could have been expected; when $`J_{ij}=0`$ for all $`(i,j)`$, for instance, one finds continually increasing phases $`\varphi _i(t)=\varphi _i(0)+\omega _it`$. Removing the $`\omega _i`$, in contrast, gives the bound $`L_j\|J_{ij}\|`$. Now the system must go to a fixed-point. In the special case $`\omega _i=\omega `$ ($`N`$ identical natural frequencies) we can transform away the $`\omega _i`$ by putting $`\varphi (t)=\stackrel{~}{\varphi }_i(t)+\omega t`$, and find the relative phases $`\stackrel{~}{\varphi }_i`$ to go to a fixed-point. ### 2.3 Detailed Balance & Equilibrium Statistical Mechanics Detailed Balance for Binary Networks. The results obtained above indicate that networks with symmetric synapses are a special class. We now show how synaptic symmetry is closely related to the detailed balance property, and derive a number of consequences. An ergodic Markov chain of the form (4,5), i.e. $$p_{\mathrm{}+1}(𝝈)=\underset{𝝈^{}}{}W[𝝈;𝝈^{}]p_{\mathrm{}}(𝝈^{})$$ (21) is said to obey detailed balance if its (unique) stationary solution $`p_{\mathrm{}}(𝝈)`$ has the property $$W[𝝈;𝝈^{}]p_{\mathrm{}}(𝝈^{})=W[𝝈^{};𝝈]p_{\mathrm{}}(𝝈)\mathrm{for}\mathrm{all}𝝈,𝝈^{}$$ (22) All $`p_{\mathrm{}}(𝝈)`$ which satisfy (22) are stationary solutions of (21), this is easily verified by substitution. The converse is not true. Detailed balance states that, in addition to $`p_{\mathrm{}}(𝝈)`$ being stationary, one has equilibrium: there is no net probability current between any two microscopic system states. It is not a trivial matter to investigate systematically for which choices of the threshold noise distribution $`w(\eta )`$ and the synaptic matrix $`\{J_{ij}\}`$ detailed balance holds. It can be shown that, apart from trivial cases (e.g. systems with self-interactions only) a Gaussian distribution $`w(\eta )`$ will not support detailed balance. Here we will work out details only for the choice $`w(\eta )=\frac{1}{2}[1\mathrm{tanh}^2(\eta )]`$, and for $`T>0`$ (where both discrete systems are ergodic). For parallel dynamics the transition matrix is given in (4), now with $`g[z]=\mathrm{tanh}[z]`$, and the detailed balance condition (22) becomes $$\frac{e^{\beta _i\sigma _ih_i(𝝈^{})}p_{\mathrm{}}(𝝈^{})}{_i\mathrm{cosh}[\beta h_i(𝝈^{})]}=\frac{e^{\beta _i\sigma _i^{}h_i(𝝈)}p_{\mathrm{}}(𝝈)}{_i\mathrm{cosh}[\beta h_i(𝝈)]}\mathrm{for}\mathrm{all}𝝈,𝝈^{}$$ (23) All $`p_{\mathrm{}}(𝝈)`$ are non-zero (ergodicity), so we may safely put $`p_{\mathrm{}}(𝝈)=e^{\beta [_i\theta _i\sigma _i+K(𝝈)]}_i\mathrm{cosh}[\beta h_i(𝝈)]`$, which, in combination with definition (1) simplifies the detailed balance condition to: $$K(𝝈)K(𝝈^{})=\underset{ij}{}\sigma _i\left[J_{ij}J_{ji}\right]\sigma _j^{}\mathrm{for}\mathrm{all}𝝈,𝝈^{}$$ (24) Averaging (24) over all possible $`𝝈^{}`$ gives $`K(𝝈)=K(𝝈^{})_𝝈^{}`$ for all $`𝝈`$, i.e. $`K`$ is a constant, whose value follows from normalising $`p_{\mathrm{}}(𝝈)`$. So, if detailed balance holds the equilibrium distribution must be: $$p_{\mathrm{eq}}(𝝈)e^{\beta _i\theta _i\sigma _i}\underset{i}{}\mathrm{cosh}[\beta h_i(𝝈)]$$ (25) For symmetric systems detailed balance indeed holds: (25) solves (23), since $`K(𝝈)=K`$ solves the reduced problem (24). For non-symmetric systems, however, there can be no equilibrium. For $`K(𝝈)=K`$ the condition (24) becomes $`_{ij}\sigma _i\left[J_{ij}J_{ji}\right]\sigma _j^{}=0`$ for all $`𝝈,𝝈^{}\{1,1\}^N`$. For $`N2`$ the vector pairs $`(𝝈,𝝈^{})`$ span the space of all $`N\times N`$ matrices, so $`J_{ij}J_{ji}`$ must be zero. For $`N=1`$ there simply exists no non-symmetric synaptic matrix. In conclusion: for binary networks with parallel dynamics, interaction symmetry implies detailed balance, and vice versa. For sequential dynamics, with $`w(\eta )=\frac{1}{2}[1\mathrm{tanh}^2(\eta )]`$, the transition matrix is given by (5) and the detailed balance condition (22) simplifies to $$\frac{e^{\beta \sigma _ih_i(F_i𝝈)}p_{\mathrm{}}(F_i𝝈)}{\mathrm{cosh}\left[\beta h_i(F_i𝝈)\right]}=\frac{e^{\beta \sigma _ih_i(𝝈)}p_{\mathrm{}}(𝝈)}{\mathrm{cosh}\left[\beta h_i(𝝈)\right]}\mathrm{for}\mathrm{all}𝝈\mathrm{and}\mathrm{all}i$$ Self-interactions $`J_{ii}`$, inducing $`h_i(F_i𝝈)h_i(𝝈)`$, complicate matters. Therefore we first consider systems where all $`J_{ii}=0`$. All stationary probabilities $`p_{\mathrm{}}(𝝈)`$ being non-zero (ergodicity), we may write: $$p_{\mathrm{}}(𝝈)=e^{\beta [_i\theta _i\sigma _i+\frac{1}{2}_{ij}\sigma _iJ_{ij}\sigma _j+K(𝝈)]}$$ (26) Using relations like $`_{kl}J_{kl}F_i(\sigma _k\sigma _l)=_{kl}J_{kl}\sigma _k\sigma _l2\sigma _i_{ki}\left[J_{ik}+J_{ki}\right]\sigma _k`$ we can simplify the detailed balance condition to $`K(F_i𝝈)K(𝝈)=\sigma _i_{ki}\left[J_{ik}J_{ki}\right]\sigma _k`$ for all $`𝝈`$ and all $`i`$. If to this expression we apply the general identity $`\left[1F_i\right]f(𝝈)=2\sigma _i\sigma _if(𝝈)_{\sigma _i}`$ we find for $`ij`$: $$[F_j1][F_i1]K(𝝈)=2\sigma _i\sigma _j\left[J_{ij}J_{ji}\right]\mathrm{for}\mathrm{all}𝝈\mathrm{and}\mathrm{all}ij$$ The left-hand side is symmetric under permutation of the pair $`(i,j)`$, which implies that the interaction matrix must also be symmetric: $`J_{ij}=J_{ji}`$ for all $`(i,j)`$. We now find the trivial solution $`K(𝝈)=K`$ (constant), detailed balance holds and the corresponding equilibrium distribution is $$p_{\mathrm{eq}}(𝝈)e^{\beta H(𝝈)}H(𝝈)=\frac{1}{2}\underset{ij}{}\sigma _iJ_{ij}\sigma _j\underset{i}{}\theta _i\sigma _i$$ (27) In conclusion: for binary networks with sequential dynamics, but without self-interactions, interaction symmetry implies detailed balance, and vice versa. In the case of self-interactions the situation is more complicated. However, here one can still show that non-symmetric models with detailed balance must be pathological, since the requirements can be met only for very specific choices for the $`\{J_{ij}\}`$. Detailed Balance for Networks with Continuous Neurons. Let us finally turn to the question of when we find microscopic equilibrium (stationarity without probability currents) in continuous models described by a Fokker-Planck equation (11). Note that (11) can be seen as a continuity equation for the density of a conserved quantity: $`\frac{d}{dt}p_t(𝝈)+_i\frac{}{\sigma _i}J_i(𝝈,t)=0`$. The components $`J_i(𝝈,t)`$ of the current density are given by $$J_i(𝝈,t)=[f_i(𝝈)T\frac{}{\sigma _i}]p_t(𝝈)$$ Stationary distributions $`p_{\mathrm{}}(𝝈)`$ are those which give $`_i\frac{}{\sigma _i}J_i(𝝈,\mathrm{})=0`$ (divergence-free currents). Detailed balance implies the stronger statement $`J_i(𝝈,\mathrm{})=0`$ for all $`i`$ (zero currents), so $`f_i(𝝈)=T\mathrm{log}p_{\mathrm{}}(𝝈)/\sigma _i`$, or $$f_i(𝝈)=H(𝝈)/\sigma _i,p_{\mathrm{}}(𝝈)e^{\beta H(𝝈)}$$ (28) for some $`H(𝝈)`$, i.e. the forces $`f_i(𝝈)`$ must be conservative. However, one can have conservative forces without a normalisable equilibrium distribution. Just take $`H(𝝈)=0`$, i.e. $`f_i(𝝈,t)=0`$: here we have $`p_{\mathrm{eq}}(𝝈)=C`$, which is not normalisable for $`𝝈\mathrm{}^N`$. For this particular case equation (11) is solved easily: $`p_t(𝝈)=[4\pi Tt]^{N/2}𝑑𝝈^{}p_0(𝝈^{})e^{[𝝈𝝈^{}]^2/4Tt}`$, so the limit $`lim_t\mathrm{}p_t(𝝈)`$ indeed does not exist. One can prove the following (see e.g. ). If the forces are conservative and if $`p_{\mathrm{}}(𝝈)e^{\beta H(𝝈)}`$ is normalisable, then it is the unique stationary solution of the Fokker-Planck equation, to which the system converges for all initial distributions $`p_0L^1[\mathrm{}^N]`$ which obey $`_\mathrm{}^N𝑑𝝈e^{\beta H(𝝈)}p_0^2(𝝈)<\mathrm{}`$. Assessing when our two particular model examples of graded response neurons or coupled oscillators obey detailed balance has thus been reduced mainly to checking whether the associated deterministic forces $`f_i(𝝈)`$ are conservative. Note that conservative forces must obey $$\mathrm{for}\mathrm{all}𝝈,\mathrm{for}\mathrm{all}ij:f_i(𝝈)/\sigma _jf_j(𝝈)/\sigma _i=0$$ (29) In the graded response equations (18) the deterministic forces are $`f_i(𝒖)=_jJ_{ij}\mathrm{tanh}[\gamma u_j]u_i+\theta _i`$. Here $`f_i(𝒖)/u_jf_j(𝒖)/u_i=\gamma \{J_{ij}[1\mathrm{tanh}^2[\gamma u_j]J_{ji}[1\mathrm{tanh}^2[\gamma u_i]\}`$. At $`𝒖=\mathrm{𝟎}`$ this reduces to $`J_{ij}J_{ji}`$, i.e. the interaction matrix must be symmetric. For symmetric matrices we find away from $`𝒖=\mathrm{𝟎}`$: $`f_i(𝒖)/u_jf_j(𝒖)/u_i=\gamma J_{ij}\{\mathrm{tanh}^2[\gamma u_i]\mathrm{tanh}^2[\gamma u_j]\}`$. The only way for this to be zero for any $`𝒖`$ is by having $`J_{ij}=0`$ for all $`ij`$, i.e. all neurons are disconnected (in this trivial case the system (18) does indeed obey detailed balance). Network models of interacting graded-response neurons of the type (18) apparently never reach equilibrium, they will always violate detailed balance and exhibit microscopic probability currents. In the case of coupled oscillators (13), where the deterministic forces are $`f_i(\mathit{\varphi })=_jJ_{ij}\mathrm{sin}[\varphi _j\varphi _i]+\omega _i`$ one finds the left-hand side of condition (29) to give $`f_i(\mathit{\varphi })/\varphi _jf_j(\mathit{\varphi })/\varphi _i=[J_{ij}J_{ji}]\mathrm{cos}[\varphi _j\varphi _i]`$. Requiring this to be zero for any $`\mathit{\varphi }`$ gives the condition $`J_{ij}=J_{ji}`$ for any $`ij`$. We have already seen that symmetric oscillator networks indeed have conservative forces: $`f_i(\mathit{\varphi })=H(\mathit{\varphi })/\varphi _i`$, with $`H(\mathit{\varphi })=\frac{1}{2}_{ij}J_{ij}\mathrm{cos}[\varphi _i\varphi _j]_i\omega _i\varphi _i`$. If in addition we choose all $`\omega _i=0`$ the function $`H(𝝈)`$ will also be bounded from below, and, although $`p_{\mathrm{}}(\mathit{\varphi })e^{\beta H(\mathit{\varphi })}`$ is still not normalisable on $`\mathit{\varphi }\mathrm{}^N`$, the full $`2\pi `$-periodicity of the function $`H(𝝈)`$ now allows us to identify $`\varphi _i+2\pi \varphi _i`$ for all $`i`$, so that now $`\mathit{\varphi }[\pi ,\pi ]^N`$ and $`𝑑\mathit{\varphi }e^{\beta H(\mathit{\varphi })}`$ does exist. Thus symmetric coupled oscillator networks with zero natural frequencies obey detailed balance. In the case of non-zero natural frequencies, in contrast, detailed balance does not hold. Equilibrium Statistical Mechanics. The above results establish the link with equilibrium statistical mechanics (see e.g. ). For binary systems with symmetric synapses (in the sequential case: without self-interactions) and with threshold noise distributions of the form $`w(\eta )=\frac{1}{2}[1\mathrm{tanh}^2(\eta )]`$, detailed balance holds and we know the equilibrium distributions. For sequential dynamics it has the Boltzmann form (27) and we can apply standard equilibrium statistical mechanics. The parameter $`\beta `$ can formally be identified with the inverse ‘temperature’ in equilibrium, $`\beta =T^1`$, and the function $`H(𝝈)`$ is the usual Ising spin Hamiltonian. In particular we can define the partition function $`Z`$ and the free energy $`F`$: $$p_{\mathrm{eq}}(𝝈)=\frac{1}{Z}e^{\beta H(𝝈)}H(𝝈)=\frac{1}{2}\underset{ij}{}\sigma _iJ_{ij}\sigma _j\underset{i}{}\theta _i\sigma _i$$ (30) $$Z=\underset{𝝈}{}e^{\beta H(𝝈)}F=\beta ^1\mathrm{log}Z$$ (31) The free energy can be used as the generating function for equilibrium averages. Taking derivatives with respect to external fields $`\theta _i`$ and interactions $`J_{ij}`$, for instance, produces $`\sigma _i=F/\theta _i`$ and $`\sigma _i\sigma _j=F/J_{ij}`$, whereas equilibrium averages of arbitrary state variable $`f(𝝈)`$ can be obtained by adding suitable generating terms to the Hamiltonian: $`H(𝝈)H(𝝈)+\lambda f(𝝈)`$, $`f=lim_{\lambda 0}F/\lambda `$. In the parallel case (25) we can again formally write the equilibrium probability distribution in the Boltzmann form and define a corresponding partition function $`\stackrel{~}{Z}`$ and a free energy $`\stackrel{~}{F}`$: $$p_{\mathrm{eq}}(𝝈)=\frac{1}{Z}e^{\beta \stackrel{~}{H}(𝝈)}\stackrel{~}{H}(𝝈)=\underset{i}{}\theta _i\sigma _i\frac{1}{\beta }\underset{i}{}\mathrm{log}2\mathrm{cosh}[\beta h_i(𝝈)]$$ (32) $$\stackrel{~}{Z}=\underset{𝝈}{}e^{\beta \stackrel{~}{H}(𝝈)}\stackrel{~}{F}=\beta ^1\mathrm{log}\stackrel{~}{Z}$$ (33) which again serve to generate averages: $`\stackrel{~}{H}(𝝈)\stackrel{~}{H}(𝝈)+\lambda f(𝝈)`$, $`f=lim_{\lambda 0}\stackrel{~}{F}/\lambda `$. However, standard thermodynamic relations involving derivation with respect to $`\beta `$ need no longer be valid, and derivation with respect to fields or interactions generates different types of averages, such as $$\stackrel{~}{F}/\theta _i=\sigma _i+\mathrm{tanh}[\beta h_i(𝝈)]\stackrel{~}{F}/J_{ii}=\sigma _i\mathrm{tanh}[\beta h_i(𝝈)]$$ $$ij:\stackrel{~}{F}/J_{ij}=\sigma _i\mathrm{tanh}[\beta h_j(𝝈)]+\sigma _j\mathrm{tanh}[\beta h_i(𝝈)]$$ One can use $`\sigma _i=\mathrm{tanh}[\beta h_i(𝝈)]`$, which can be derived directly from the equilibrium equation $`p_{\mathrm{eq}}(𝝈)=_𝝈^{}W[𝝈;𝝈^{}]p_{\mathrm{eq}}(𝝈)`$, to simplify the first of these identities. A connected network of graded-response neurons can never be in an equilibrium state, so our only model example with continuous neuronal variables for which we can set up the equilibrium statistical mechanics formalism is the system of coupled oscillators (13) with symmetric synapses and absent (or uniform) natural frequencies $`\omega _i`$. If we define the phases as $`\varphi _i[\pi ,\pi ]`$ we have again an equilibrium distribution of the Boltzmann form, and we can define the standard thermodynamic quantities: $$p_{\mathrm{eq}}(\mathit{\varphi })=\frac{1}{Z}e^{\beta H(\mathit{\varphi })}H(\mathit{\varphi })=\frac{1}{2}\underset{ij}{}J_{ij}\mathrm{cos}[\varphi _i\varphi _j]$$ (34) $$Z=_\pi ^\pi \mathrm{}_\pi ^\pi 𝑑\mathit{\varphi }e^{\beta H(\mathit{\varphi })}F=\beta ^1\mathrm{log}Z$$ (35) These generate equilibrium averages in the usual manner. For instance $`\mathrm{cos}[\varphi _i\varphi _j]=F/J_{ij}`$, whereas averages of arbitrary state variables $`f(\mathit{\varphi })`$ follow, as before, upon introducing suitable generating terms: $`H(\mathit{\varphi })H(\mathit{\varphi })+\lambda f(\mathit{\varphi })`$, $`f=lim_{\lambda 0}F/\lambda `$. In this chapter we restrict ourselves to symmetric networks which obey detailed balance, so that we know the equilibrium probability distribution and equilibrium statistical mechanics applies. In the case of sequential dynamics we will accordingly not allow for the presence of self-interactions. ## 3 Simple Recurrent Networks with Binary Neurons ### 3.1 Networks with Uniform Synapses We now turn to a simple toy model to show how equilibrium statistical mechanics is used for solving neural network models, and to illustrate similarities and differences between the different dynamics types. We choose uniform infinite-range synapses and zero external fields, and calculate the free energy for the binary systems (1,2), parallel and sequential, and with threshold noise distribution $`w(\eta )=\frac{1}{2}[1\mathrm{tanh}^2(\eta )]`$: $$J_{ij}=J_{ji}=J/N(ij),J_{ii}=\theta _i=0\mathrm{for}\mathrm{all}i$$ The free energy is an extensive object, $`lim_N\mathrm{}F/N`$ is finite. For the models (1,2) we now obtain: $$\mathrm{𝐵𝑖𝑛𝑎𝑟𝑦}\&\mathrm{𝑆𝑒𝑞𝑢𝑒𝑛𝑡𝑖𝑎𝑙}:\underset{N\mathrm{}}{lim}F/N=\underset{N\mathrm{}}{lim}(\beta N)^1\mathrm{log}\underset{𝝈}{}e^{\beta N\left[\frac{1}{2}Jm^2(𝝈)\right]}$$ $$\mathrm{𝐵𝑖𝑛𝑎𝑟𝑦}\&\mathrm{𝑃𝑎𝑟𝑎𝑙𝑙𝑒𝑙}:\underset{N\mathrm{}}{lim}\stackrel{~}{F}/N=\underset{N\mathrm{}}{lim}(\beta N)^1\mathrm{log}\underset{𝝈}{}e^{N\left[\mathrm{log}2\mathrm{cosh}[\beta Jm(𝝈)]\right]}$$ with the average activity $`m(𝝈)=\frac{1}{N}_k\sigma _k`$. We have to count the number of states $`𝝈`$ with a prescribed average activity $`m=2n/N1`$ ($`n`$ is the number of neurons $`i`$ with $`\sigma _i=1`$), in expressions of the form $$\frac{1}{N}\mathrm{log}\underset{𝝈}{}e^{NU\left[m(𝝈)\right]}=\frac{1}{N}\mathrm{log}\underset{n=0}{\overset{N}{}}(\begin{array}{c}N\\ n\end{array})e^{NU[2n/N1]}=\frac{1}{N}\mathrm{log}_1^1𝑑me^{N\left[\mathrm{log}2c^{}(m)+U[m]\right]}$$ $$\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}\underset{𝝈}{}e^{NU\left[m(𝝈)\right]}=\mathrm{log}2+\underset{m[1,1]}{\mathrm{max}}\left\{U[m]c^{}(m)\right\}$$ with the entropic function $`c^{}(m)=\frac{1}{2}(1+m)\mathrm{log}(1+m)+\frac{1}{2}(1m)\mathrm{log}(1m)`$. In order to get there we used Stirling’s formula to obtain the leading term of the factorials (only terms which are exponential in $`N`$ survive the limit $`N\mathrm{}`$), we converted (for $`N\mathrm{}`$) the summation over $`n`$ into an integration over $`m=2n/N1[1,1]`$, and we carried out the integral over $`m`$ via saddle-point integration (see e.g. ). This leads to a saddle-point problem whose solution gives the free energies: $$\underset{N\mathrm{}}{lim}F/N=\underset{m[1,1]}{\mathrm{min}}f_{\mathrm{seq}}(m)\beta f_{\mathrm{seq}}(m)=c^{}(m)\mathrm{log}2\frac{1}{2}\beta Jm^2$$ (36) $$\underset{N\mathrm{}}{lim}\stackrel{~}{F}/N=\underset{m[1,1]}{\mathrm{min}}f_{\mathrm{par}}(m)\beta f_{\mathrm{par}}(m)=c^{}(m)2\mathrm{log}2\mathrm{log}\mathrm{cosh}[\beta Jm]$$ (37) The functions to be minimised are shown in figure 2. The equations from which to solve the minima are easily obtained by differentiation, using $`\frac{d}{dm}c^{}(m)=\mathrm{tanh}^1(m)`$. For sequential dynamics we find $$\mathrm{𝐵𝑖𝑛𝑎𝑟𝑦}\&\mathrm{𝑆𝑒𝑞𝑢𝑒𝑛𝑡𝑖𝑎𝑙}:m=\mathrm{tanh}[\beta Jm]$$ (38) (the so-called Curie-Weiss law). For parallel dynamics we find $$m=\mathrm{tanh}\left[\beta J\mathrm{tanh}[\beta Jm]\right]$$ One finds that the solutions of the latter equation again obey a Curie-Weiss law. The definition $`\widehat{m}=\mathrm{tanh}[\beta \|J\|m]`$ transforms it into the coupled equations $`m=\mathrm{tanh}[\beta \|J\|\widehat{m}]`$ and $`\widehat{m}=\mathrm{tanh}[\beta \|J\|m]`$, from which we derive $`0[m\widehat{m}]^2=[m\widehat{m}]\left[\mathrm{tanh}[\beta \|J\|\widehat{m}]\mathrm{tanh}[\beta \|J\|m]\right]0`$. Since $`\mathrm{tanh}[\beta \|J\|m]`$ is a monotonically increasing function of $`m`$, this implies $`\widehat{m}=m`$, so $$\mathrm{𝐵𝑖𝑛𝑎𝑟𝑦}\&\mathrm{𝑃𝑎𝑟𝑎𝑙𝑙𝑒𝑙}:m=\mathrm{tanh}[\beta \|J\|m]$$ (39) Our study of the toy models has thus been reduced to analysing the non-linear equations (38) and (39). If $`J0`$ (excitation) the two types of dynamics lead to the same behaviour. At high noise levels, $`T>J`$, both minimisation problems are solved by $`m=0`$ (see figure 2), describing a disorganised (paramagnetic) state. This can be seen upon writing the right-hand side of (38) in integral form: $$m^2=m\mathrm{tanh}[\beta Jm]=\beta Jm^2_0^1𝑑z[1\mathrm{tanh}^2[\beta Jmz]]\beta Jm^2$$ So $`m^2[1\beta J]0`$, which gives $`m=0`$ as soon as $`\beta J<1`$. A phase transition occurs at $`T=J`$ (a bifurcation of non-trivial solutions of (38)), and for $`T<J`$ the equations for $`m`$ are solved by the two non-zero solutions of (38), describing a state where either all neurons tend to be firing ($`m>0`$) or where they tend to be quiet ($`m<0`$). This becomes clear when we expand (38) for small $`m`$: $`m=\beta Jm+𝒪(m^3)`$, so precisely at $`\beta J=1`$ one finds a de-stabilisation of the trivial solution $`m=0`$, together with the creation of (two) stable non-trivial ones (see also figure 2). Furthermore, using the identity $`c^{}(\mathrm{tanh}x)=x\mathrm{tanh}x\mathrm{log}\mathrm{cosh}x`$, we obtain from (36,37) the relation $`lim_N\mathrm{}\stackrel{~}{F}/N=2lim_N\mathrm{}F/N`$. For $`J<0`$ (inhibition), however, the two types of dynamics give quite different results. For sequential dynamics the relevant minimum is located at $`m=0`$ (the paramagnetic state). For parallel dynamics, the minimisation problem is invariant under $`JJ`$, so the behaviour is again of the Curie-Weiss type (see figure 2 and equation (39)), with a paramagnetic state for $`T>\|J\|`$, a phase transition at $`T=\|J\|`$, and order for $`T<\|J\|`$. This difference between the two types of dynamics for $`J<0`$ is explained by studying dynamics. As we will see in a subsequent chapter, for the present (toy) model in the limit $`N\mathrm{}`$ the average activity evolves in time according to the deterministic laws $$\frac{d}{dt}m=\mathrm{tanh}[\beta Jm]mm(t+1)=\mathrm{tanh}[\beta Jm(t)]$$ for sequential and parallel dynamics, respectively. For $`J<0`$ the sequential system always decays towards the trivial state $`m=0`$, whereas for sufficiently large $`\beta `$ the parallel system enters the stable limit-cycle $`m(t)=M_\beta (1)^t`$ (where $`M_\beta `$ is the non-zero solution of (39)). The concepts of ‘distance’ and ‘local minima’ are quite different for the two dynamics types; in contrast to the sequential case, parallel dynamics allows the system to make the transition $`mm`$ in equilibrium. ### 3.2 Phenomenology of Hopfield Models The Ideas Behind the Hopfield Model. The Hopfield model is a network of binary neurons of the type (1,2), with threshold noise $`w(\eta )=\frac{1}{2}[1\mathrm{tanh}^2(\eta )]`$, and with a specific recipe for the synapses $`J_{ij}`$ aimed at storing patterns, motivated by suggestions made in the late nineteen-forties . The original model was in fact defined more narrowly, as the zero noise limit of the system (2), but the term has since then been accepted to cover a larger network class. Let us first consider the simplest case and try to store a single pattern $`𝝃\{1,1\}^N`$ in noise-less infinite-range binary networks. Appealing candidates for interactions and thresholds would be $`J_{ij}=\xi _i\xi _j`$ and $`\theta _i=0`$ (for sequential dynamics we put $`J_{ii}=0`$ for all $`i`$). With this choice the Lyapunov function (17) becomes: $$L_{\mathrm{seq}}[𝝈]=\frac{1}{2}N\frac{1}{2}[\underset{i}{}\xi _i\sigma _i]^2$$ It will have to decrease monotonically during the dynamics, from which we immediately deduce $$\underset{i}{}\xi _i\sigma _i(0)>0:𝝈(\mathrm{})=𝝃,\underset{i}{}\xi _i\sigma _i(0)<0:𝝈(\mathrm{})=𝝃$$ This system indeed reconstructs dynamically the original pattern $`𝝃`$ from an input vector $`𝝈(0)`$, at least for sequential dynamics. However, en passant we have created an additional attractor: the state $`𝝃`$. This property is shared by all binary models in which the external fields are zero, where the Hamiltonians $`H(𝝈)`$ (30) and $`\stackrel{~}{H}(𝝈)`$ (32) are invariant under an overall sign change $`𝝈𝝈`$. A second feature common to several (but not all) attractor neural networks is that each initial state will lead to pattern reconstruction, even nonsensical (random) ones. The Hopfield model is obtained by generalising the previous simple one-pattern recipe to the case of an arbitrary number $`p`$ of binary patterns $`𝝃^\mu =(\xi _1^\mu ,\mathrm{},\xi _N^\mu )\{1,1\}^N`$: $$J_{ij}=\frac{1}{N}\underset{\mu =1}{\overset{p}{}}\xi _i^\mu \xi _j^\mu ,\theta _i=0\mathrm{for}\mathrm{all}i(\mathrm{sequential}\mathrm{dynamics}:J_{ii}0\mathrm{for}\mathrm{all}i)$$ (40) The prefactor $`N^1`$ has been inserted to ensure that the limit $`N\mathrm{}`$ will exist in future expressions. The process of interest is that where, triggered by correlation between the initial state and a stored pattern $`𝝃^\lambda `$, the state vector $`𝝈`$ evolves towards $`𝝃^\lambda `$. If this happens, pattern $`𝝃^\lambda `$ is said to be recalled. The similarity between a state vector and the stored patterns is measured by so-called overlaps $$m_\mu (𝝈)=\frac{1}{N}\underset{i}{}\xi _i^\mu \sigma _i$$ (41) Numerical simulations illustrate the functioning of the Hopfield model as an associative memory, and the description of the recall process in terms of overlaps. Our simulated system is an $`N=841`$ Hopfield model, in which $`p=10`$ patterns have been stored (see figure 3) according to prescription (40). The two-dimensional arrangement of the neurons in this example is just a guide to the eye; since the network is fully connected the physical location of the neurons is irrelevant. The dynamics is as given by (2), with $`T=0.1`$. In figure 4 we first show (left column) the result of letting the system evolve in time from an initial state, which is a noisy version of one of the stored patterns (here $`40\%`$ of the neuronal states $`\sigma _i`$ where corrupted, according to $`\sigma _i\sigma _i`$). The top left row of graphs shows snapshots of the microscopic state as the system evolves in time. The bottom left row shows the values of the $`p=10`$ overlaps $`m_\mu `$, as defined in (41), as functions of time; the one which evolves towards the value 1 corresponds to the pattern being reconstructed. The right column of figure 4 shows a similar experiment, but here the initial state is drawn at random. The system subsequently evolves towards a mixture of the stored patterns, which is found to be very stable, due to the fact that the patterns involved (see figure 3) are significantly correlated. It will be clear that, although the idea of information storage via the creation of attractors does work, the choice (40) for the synapses is still too simple to be optimal; in addition to the desired states $`𝝃^\mu `$ and their mirror images $`𝝃^\mu `$, even more unwanted spurious attractors are created. Yet this model will already push the analysis to the limits, as soon as we allow for the storage of an extensive number of patterns $`𝝃^\mu `$. Issues Related to Saturation: Storage Capacity & Non-Trivial Dynamics. In our previous simulation example the loading of the network was modest; a total of $`\frac{1}{2}N(N1)=353,220`$ synapses were used to store just $`pN=8,410`$ bits of information. Let us now investigate the behaviour of the network when the number of patterns scales with the system size as $`p=\alpha N`$ ($`\alpha >0`$); now for large $`N`$ the number of bits stored per synapse will be $`pN/\frac{1}{2}N(N1)2\alpha `$. This is called the saturation regime. Again numerical simulations, but now with finite $`\alpha `$, illustrate the main features and complications of recall dynamics in the saturation regime. In our example the dynamics are given by (1) (parallel updates), with $`T=0.1`$ and threshold noise distribution $`w(\eta )=\frac{1}{2}[1\mathrm{tanh}^2(\eta )]`$; the patterns are chosen randomly. Figure 5 shows the result of measuring in such simulations the two quantities $$m=m_1(𝝈)r=\alpha ^1\underset{\mu >1}{}m_\mu ^2(𝝈)$$ (42) following initial states which are correlated with pattern $`𝝃^1`$ only. For large $`N`$ we can distinguish structural overlaps, where $`m_\mu (𝝈)=𝒪(1)`$, from accidental ones, where $`m_\mu (𝝈)=𝒪(N^{\frac{1}{2}})`$ (as for a randomly drawn $`𝝈`$). Overlaps with non-nominated patterns are seen to remain $`𝒪(N^{\frac{1}{2}})`$, i.e. $`r(t)=𝒪(1)`$. We observe competition between pattern recall ($`m1`$) and interference of non-nominated patterns ($`m0`$, with $`r`$ increasing), and a profound slowing down of the process for non-recall trajectories. The initial overlap (the ‘cue’) needed to trigger recall is found to increase with increasing $`\alpha `$ (the loading) and increasing $`T`$ (the noise). Further numerical experimentation, with random patterns, reveals that at any noise level $`T`$ there is a critical storage level $`\alpha _c(T)`$ above which recall is impossible, with an absolute upper limit of $`\alpha _c=\mathrm{max}_T\alpha _c(T)=\alpha _c(0)0.139`$. The competing forces at work are easily recognised when working out the local fields (1), using (40): $$h_i(𝝈)=\xi _i^1m_1(𝝈)+\frac{1}{N}\underset{\mu >1}{}\xi _i^\mu \underset{ji}{}\xi _j^\mu \sigma _j+𝒪(N^1)$$ (43) The first term in (43) drives $`𝝈`$ towards pattern $`𝝃^1`$ as soon as $`m_1(𝝈)>0`$. The second terms represent interference, caused by correlations between $`𝝈`$ and non-nominated patterns. One easily shows (to be demonstrated later) that for $`N\mathrm{}`$ the fluctuations in the values of the recall overlap $`m`$ will vanish, and that for the present types of initial states and threshold noise the overlap $`m`$ will obey $$m(t+1)=𝑑zP_t(z)\mathrm{tanh}[\beta (m(t)+z)]P_t(z)=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{i}{}\delta [z\frac{1}{N}\underset{\mu >1}{}\xi _i^1\xi _i^\mu \underset{ji}{}\xi _j^\mu \sigma _j(t)]$$ (44) If all $`\sigma _i(0)`$ are drawn independently, $`\mathrm{Prob}[\sigma _i(0)=\pm \xi _i^1]=\frac{1}{2}[1\pm m(0)]`$, the central limit theorem states that $`P_0(z)`$ is Gaussian. One easily derives $`z_0=0`$ and $`z^2_0=\alpha `$, so at $`t=0`$ equation (44) gives $$m(1)=\frac{dz}{\sqrt{2\pi }}e^{\frac{1}{2}z^2}\mathrm{tanh}[\beta (m(0)+z\sqrt{\alpha })]$$ (45) The above ideas, and equation (45) in particular, go back to . For times $`t>0`$, however, the independence of the states $`\sigma _i`$ need no longer hold. As a simple approximation one could just assume that the $`\sigma _i`$ remain uncorrelated at all times, i.e. $`\mathrm{Prob}[\sigma _i(t)=\pm \xi _i^1]=\frac{1}{2}[1\pm m(t)]`$ for all $`t0`$, such that the argument given for $`t=0`$ would hold generally, and where (for randomly drawn patterns) the mapping (45) would describe the overlap evolution at all times: $$m(t+1)=\frac{dz}{\sqrt{2\pi }}e^{\frac{1}{2}z^2}\mathrm{tanh}[\beta (m(t)+z\sqrt{\alpha })]$$ (46) This equation, however, must be generally incorrect. Firstly, figure 5 already shows that knowledge of $`m(t)`$ only does not yet permit prediction of $`m(t+1)`$. Secondly, upon working out its bifurcation properties one finds that equation (46) predicts a storage capacity of $`\alpha _c=2/\pi 0.637`$, which is no way near to what is actually being observed. We will see in the paper on dynamics that only for certain types of extremely diluted networks (where most of the synapses are cut) equation (46) is indeed correct on finite times; in these networks the time it takes for correlations between neuron states to build up diverges with $`N`$, so that correlations are simply not yet noticable on finite times. For fully connected Hopfield networks storing random patterns near saturation, i.e. with $`\alpha >0`$, the complicated correlations building up between the microscopic variables in the course of the dynamics generate an interference noise distribution which is intrinsically non-Gaussian, see e.g. figure 6. This leads to a highly non-trivial dynamics which is fundamentally different from that in the $`lim_N\mathrm{}p/N=0`$ regime. Solving models of recurrent neural networks in the saturation regime boils down to calculating this non-Gaussian noise distribution, which requires advanced mathematical techniques (in statics and dynamics), and constitutes the main challenge to the theorist. The simplest way to evade this challenge is to study situations where the interference noise is either trivial (as with asymmetric extremely diluted models) or where it vanishes, which happens in fully connected networks when $`\alpha =lim_N\mathrm{}p/N=0`$ (as with finite $`p`$). The latter $`\alpha =0`$ regime is the one we will explore first. ### 3.3 Analysis of Hopfield Models Away From Saturation Equilibrium Order Parameter Equations. A binary Hopfield network with parameters given by (40) obeys detailed balance, and the Hamiltonian $`H(𝝈)`$ (30) (corresponding to sequential dynamics) and the pseudo-Hamiltonian $`\stackrel{~}{H}(𝝈)`$ (32) (corresponding to parallel dynamics) become $$H(𝝈)=\frac{1}{2}N\underset{\mu =1}{\overset{p}{}}m_\mu ^2(𝝈)+\frac{1}{2}p\stackrel{~}{H}(𝝈)=\frac{1}{\beta }\underset{i}{}\mathrm{log}2\mathrm{cosh}[\beta \underset{\mu =1}{\overset{p}{}}\xi _i^\mu m_\mu (𝝈)]$$ (47) with the overlaps (41). Solving the statics implies calculating the free energies $`F`$ and $`\stackrel{~}{F}`$: $$F=\frac{1}{\beta }\mathrm{log}\underset{𝝈}{}e^{\beta H(𝝈)}\stackrel{~}{F}=\frac{1}{\beta }\mathrm{log}\underset{𝝈}{}e^{\beta \stackrel{~}{H}(𝝈)}$$ Upon introducing the short-hand notation $`𝒎=(m_1,\mathrm{},m_p)`$ and $`𝝃_i=(\xi _i^1,\mathrm{},\xi _i^p)`$, both free energies can be expressed in terms of the density of states $`𝒟(𝒎)=2^N_𝝈\delta [𝒎𝒎(𝝈)]`$: $$F/N=\frac{1}{\beta }\mathrm{log}2\frac{1}{\beta N}\mathrm{log}𝑑𝒎𝒟(𝒎)e^{\frac{1}{2}\beta N𝒎^2}+\frac{p}{2N}$$ (48) $$\stackrel{~}{F}/N=\frac{1}{\beta }\mathrm{log}2\frac{1}{\beta N}\mathrm{log}𝑑𝒎𝒟(𝒎)e^{_{i=1}^N\mathrm{log}2\mathrm{cosh}\left[\beta 𝝃_i𝒎\right]}$$ (49) (note: $`𝑑𝒎\delta [𝒎𝒎(𝝈)]=1`$). In order to proceed we need to specify how the number of patterns $`p`$ scales with the system size $`N`$. In this section we will follow (equilibrium analysis following sequential dynamics) and (equilibrium analysis following parallel dynamics), and assume $`p`$ to be finite. One can now easily calculate the leading contribution to the density of states, using the integral representation of the $`\delta `$-function and keeping in mind that according to (48,49) only terms exponential in $`N`$ will retain statistical relevance for $`N\mathrm{}`$: $$\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}𝒟(𝒎)=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}𝑑𝒙e^{iN𝒙𝒎}e^{i_{i=1}^N\sigma _i𝝃_i𝒙}_𝝈$$ $$=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}𝑑𝒙e^{N[i𝒙𝒎+\mathrm{log}\mathrm{cos}[𝝃𝒙]_𝝃]}$$ with the abbreviation $`\mathrm{\Phi }(𝝃)_𝝃=lim_N\mathrm{}\frac{1}{N}_{i=1}^N\mathrm{\Phi }(𝝃_i)`$. The leading contribution to both free energies can be expressed as a finite-dimensional integral, for large $`N`$ dominated by that saddle-point (extremum) for which the extensive exponent is real and maximal: $$\underset{N\mathrm{}}{lim}F/N=\frac{1}{\beta N}\mathrm{log}𝑑𝒎𝑑𝒙e^{N\beta f(𝒎,𝒙)}=\mathrm{extr}_{𝒙,𝒎}f(𝒎,𝒙)$$ $$\underset{N\mathrm{}}{lim}\stackrel{~}{F}/N=\frac{1}{\beta N}\mathrm{log}𝑑𝒎𝑑𝒙e^{N\beta \stackrel{~}{f}(𝒎,𝒙)}=\mathrm{extr}_{𝒙,𝒎}\stackrel{~}{f}(𝒎,𝒙)$$ with $$\begin{array}{cc}f(𝒎,𝒙)\hfill & =\frac{1}{2}𝒎^2i𝒙𝒎\beta ^1\mathrm{log}2\mathrm{cos}\left[\beta 𝝃𝒙\right]_𝝃\hfill \\ \stackrel{~}{f}(𝒎,𝒙)\hfill & =\beta ^1\mathrm{log}2\mathrm{cosh}\left[\beta 𝝃𝒎\right]_𝝃i𝒙𝒎\beta ^1\mathrm{log}2\mathrm{cos}\left[\beta 𝝃𝒙\right]_𝝃\hfill \end{array}$$ The saddle-point equations for $`f`$ and $`\stackrel{~}{f}`$ are given by: $$\begin{array}{ccccc}f:\hfill & & 𝒙=i𝒎,\hfill & & i𝒎=𝝃\mathrm{tan}\left[\beta 𝝃𝒙\right]_𝝃\hfill \\ \stackrel{~}{f}:\hfill & & 𝒙=i𝝃\mathrm{tanh}\left[\beta 𝝃𝒎\right]_𝝃,\hfill & & i𝒎=𝝃\mathrm{tan}\left[\beta 𝝃𝒙\right]_𝝃\hfill \end{array}$$ In saddle-points $`𝒙`$ turns out to be purely imaginary. However, after a shift of the integration contours, putting $`𝒙=i𝒙^{}(𝒎)+𝒚`$ (where $`i𝒙^{}(𝒎)`$ is the imaginary saddle-point, and where $`𝒚\mathrm{}^p`$) we can eliminate $`𝒙`$ in favor of $`𝒚\mathrm{}^p`$ which does have a real saddle-point, by construction.<sup>1</sup><sup>1</sup>1Our functions to be integrated have no poles, but strictly speaking we still have to verify that the integration segments linking the original integration regime to the shifted one will not contribute to the integrals. This is generally a tedious and distracting task, which is often skipped. For simple models, however (e.g. networks with uniform synapses), the verification can be carried out properly, and all is found to be safe. We then obtain<sup>2</sup><sup>2</sup>2Here we used the equation $`f(𝒎,𝒙)/𝒎=\mathrm{𝟎}`$ to express $`𝒙`$ in terms of $`𝒎`$, because this is simpler. Strictly speaking we should have used $`f(𝒎,𝒙)/𝒙=\mathrm{𝟎}`$ for this purpose; our short-cut could in principle generate additional solutions. In the present model, however, we can check explicitly that this is not the case. Also, in view of the imaginary saddle-point $`𝒙`$, we can not be certain that, upon elimination of $`𝒙`$, the relevant saddle-point of the remaining function $`f(𝒎)`$ must be a minimum. This will have to be checked, for instance by inspection of the $`T\mathrm{}`$ limit. $$\begin{array}{ccc}\mathrm{𝑆𝑒𝑞𝑢𝑒𝑛𝑡𝑖𝑎𝑙}\mathrm{𝐷𝑦𝑛𝑎𝑚𝑖𝑐𝑠}:\hfill & & 𝒎=𝝃\mathrm{tanh}[\beta 𝝃𝒎]_𝝃\hfill \\ \mathrm{𝑃𝑎𝑟𝑎𝑙𝑙𝑒𝑙}\mathrm{𝐷𝑦𝑛𝑎𝑚𝑖𝑐𝑠}:\hfill & & 𝒎=𝝃\mathrm{tanh}[\beta 𝝃[𝝃^{}\mathrm{tanh}[\beta 𝝃^{}𝒎]_𝝃^{}]]_𝝃\hfill \end{array}$$ (compare to e.g. (38,39)). The solutions of the above two equations will in general be identical. To see this, let us denote $`\widehat{𝒎}=𝝃\mathrm{tanh}\left[\beta 𝝃𝒎\right]_𝝃`$, with which the saddle point equation for $`\stackrel{~}{f}`$ decouples into: $$𝒎=𝝃\mathrm{tanh}\left[\beta 𝝃\widehat{𝒎}\right]_𝝃\widehat{𝒎}=𝝃\mathrm{tanh}\left[\beta 𝝃𝒎\right]_𝝃$$ so $$\left[𝒎\widehat{𝒎}\right]^2=\left[(𝝃𝒎)(𝝃\widehat{𝒎})\right]\left[\mathrm{tanh}(\beta 𝝃\widehat{𝒎})\mathrm{tanh}(\beta 𝝃𝒎)\right]_𝝃$$ Since $`tanh`$ is a monotonicaly increasing function, we must have $`\left[𝒎\widehat{𝒎}\right]𝝃=0`$ for each $`𝝃`$ that contributes to the averages $`\mathrm{}_𝝃`$. For all choices of patterns where the covariance matrix $`C_{\mu \nu }=\xi _\mu \xi _\nu _𝝃`$ is positive definite, we thus obtain $`𝒎=\widehat{𝒎}`$. The final result is: for both types of dynamics (sequential and parallel) the overlap order parameters in equilibrium are given by the solution $`𝒎^{}`$ of $$𝒎=𝝃\mathrm{tanh}\left[\beta 𝝃𝒎\right]_𝝃$$ (50) which minimises<sup>3</sup><sup>3</sup>3We here indeed know the relevant saddle-point to be a minimum: the only solution of the saddle-point equations at high temperatures, $`𝒎=\mathrm{𝟎}`$, is seen to minimise $`f(𝒎)`$, since $`f(𝒎)+\beta ^1\mathrm{log}2=\frac{1}{2}𝒎^2(1\beta )+𝒪(\beta ^3)`$. $$f(𝒎)=\frac{1}{2}𝒎^2\frac{1}{\beta }\mathrm{log}2\mathrm{cosh}\left[\beta 𝝃𝒎\right]_𝝃$$ (51) The free energies of the ergodic compoments are $`lim_N\mathrm{}F/N=f(𝒎^{})`$ and $`lim_N\mathrm{}\stackrel{~}{F}/N=2f(𝒎^{})`$. Adding generating terms of the form $`HH+\lambda g[𝒎(𝝈)]`$ to the Hamiltonians allows us identify $`g[𝒎(𝝈)]_{\mathrm{eq}}=lim_{\lambda 0}F/\lambda =g[𝒎^{}]`$. Thus, in equilibrium the fluctuations in the overlap order parameters $`𝒎(𝝈)`$ (41) vanish for $`N\mathrm{}`$. Their deterministic values are simply given by $`𝒎^{}`$. Note that in the case of sequential dynamics we could also have used linearisation with Gaussian integrals (as used previously for coupled oscillators with uniform synapses) to arrive at this solution, with $`p`$ auxiliary integrations, but that for parallel dynamics this would not have been possible. Analysis of Order Parameter Equations: Pure States & Mixture States. We will restrict our further discussion to the case of randomly drawn patterns, so $$\mathrm{\Phi }(𝝃)_𝝃=2^p\underset{𝝃\{1,1\}^p}{}\mathrm{\Phi }(𝝃),\xi _\mu _𝝃=0,\xi _\mu \xi _\nu _𝝃=\delta _{\mu \nu }$$ (generalisation to correlated patterns is in principle straightforward). We first establish an upper bound for the temperature for where non-trivial solutions $`𝒎^{}`$ could exist, by writing (50) in integral form: $$m_\mu =\beta \xi _\mu (𝝃𝒎)_0^1𝑑\lambda [1\mathrm{tanh}^2[\beta \lambda 𝝃𝒎]]_𝝃$$ from which we deduce $$0=𝒎^2\beta (𝝃𝒎)^2_0^1𝑑\lambda [1\mathrm{tanh}^2[\beta \lambda 𝝃𝒎]]_𝝃𝒎^2\beta (𝝃𝒎)^2_𝝃=𝒎^2(1\beta )$$ For $`T>1`$ the only solution of (50) is the paramagnetic state $`𝒎=0`$, which gives for the free energy per neuron $`T\mathrm{log}2`$ and $`2T\mathrm{log}2`$ (for sequential and parallel dynamics, respectively). At $`T=1`$ a phase transition occurs, which follows from expanding (50) for small $`\|𝒎\|`$ in powers of $`\tau =\beta 1`$: $$m_\mu =(1+\tau )m_\mu \frac{1}{3}\underset{\nu \rho \lambda }{}m_\nu m_\rho m_\lambda \xi _\mu \xi _\nu \xi _\rho \xi _\lambda _𝝃+𝒪(𝒎^5,\tau 𝒎^3)=m_\mu [1+\tau 𝒎^2+\frac{2}{3}m_\mu ^2]+𝒪(𝒎^5,\tau 𝒎^3)$$ The new saddle-point scales as $`m_\mu =\stackrel{~}{m}_\mu \tau ^{1/2}+𝒪(\tau ^{3/2})`$, with for each $`\mu `$: $`\stackrel{~}{m}_\mu =0`$ or $`0=1\stackrel{~}{𝒎}^2+\frac{2}{3}\stackrel{~}{m}_\mu ^2`$. The solutions are of the form $`\stackrel{~}{m}_\mu \{\stackrel{~}{m},0,\stackrel{~}{m}\}`$. If we denote with $`n`$ the number of non-zero components in the vector $`\stackrel{~}{𝒎}`$, we derive from the above identities: $`\stackrel{~}{m}_\mu =0`$ or $`\stackrel{~}{m}_\mu =\pm \sqrt{3}/\sqrt{3n2}`$. These saddle-points are called mixture states, since they correspond to microscopic configurations correlated equally with a finite number $`n`$ of the stored patterns (or their negatives). Without loss of generality we can always perform gauge transformations on the set of stored patterns (permutations and reflections), such that the mixture states acquire the form $$𝒎=m_n(\stackrel{n\mathrm{times}}{\stackrel{}{1,\mathrm{},1}},\stackrel{pn\mathrm{times}}{\stackrel{}{0,\mathrm{},0}})m_n=[\frac{3}{3n2}]^{\frac{1}{2}}(\beta 1)^{1/2}+\mathrm{}$$ (52) These states are in fact saddle-points of the surface $`f(𝒎)`$ (51) for any finite temperature, as can be verified by substituting (52) as an ansatz into (50): $$\mu n:m_n=\xi _\mu \mathrm{tanh}[\beta m_n\underset{\nu n}{}\xi _\nu ]_𝝃\mu >n:0=\xi _\mu \mathrm{tanh}[\beta m_n\underset{\nu n}{}\xi _\nu ]_𝝃$$ The second equation is automatically satisfied since the average factorises. The first equation leads to a condition determining the amplitude $`m_n`$ of the mixture states: $$m_n=[\frac{1}{n}\underset{\mu n}{}\xi _\mu ]\mathrm{tanh}[\beta m_n\underset{\nu n}{}\xi _\nu ]_𝝃$$ (53) The corresponding values of $`f(𝒎)`$, to be denoted by $`f_n`$, are $$f_n=\frac{1}{2}nm_n^2\frac{1}{\beta }\mathrm{log}2\mathrm{cosh}[\beta m_n\underset{\nu n}{}\xi _\nu ]_𝝃$$ (54) The relevant question at this stage is whether or not these saddle-points correspond to local minima of the surface $`f(𝒎)`$ (51). The second derivative of $`f(𝒎)`$ is given by $$\frac{^2f(𝒎)}{m_\mu m_\nu }=\delta _{\mu \nu }\beta \xi _\mu \xi _\nu \left[1\mathrm{tanh}^2\left[\beta 𝝃𝒎\right]\right]_𝝃$$ (55) (a local minimum corresponds to a positive definite second derivative). In the trivial saddle-point $`𝒎=0`$ this gives simply $`\delta _{\mu \nu }(1\beta )`$, so at $`T=1`$ this state destabilises. In a mixture state of the type (52) the second derivative becomes: $$D_{\mu \nu }^{(n)}=\delta _{\mu \nu }\beta \xi _\mu \xi _\nu [1\mathrm{tanh}^2[\beta m_n\underset{\rho n}{}\xi _\rho ]]_𝝃$$ Due to the symmetries in the problem the spectrum of the matrix $`D^{(n)}`$ can be calculated. One finds the following eigenspaces, with $`Q=\mathrm{tanh}^2[\beta m_n_{\rho n}\xi _\rho ]_𝝃`$ and $`R=\xi _1\xi _2\mathrm{tanh}^2[\beta m_n_{\rho n}\xi _\rho ]_𝝃`$: $$\begin{array}{ccc}& \mathrm{Eigenspace}:\hfill & \mathrm{Eigenvalue}:\hfill \\ I:\hfill & 𝒙=(0,\mathrm{},0,x_{n+1},\mathrm{},x_p)\hfill & 1\beta [1Q]\hfill \\ II:\hfill & 𝒙=(1,\mathrm{},1,0,\mathrm{},0)\hfill & 1\beta [1Q+(1n)R]\hfill \\ III:\hfill & 𝒙=(x_1,\mathrm{},x_n,0,\mathrm{},0),_\mu x_\mu =0\hfill & 1\beta [1Q+R]\hfill \end{array}$$ Eigenspace $`III`$ and the quantity $`R`$ only come into play for $`n>1`$. To find the smallest eigenvalue we need to know the sign of $`R`$. With the abbreviation $`M_𝝃=_{\rho n}\xi _\rho `$ we find: $$\begin{array}{cc}n(n1)R\hfill & =M_𝝃^2\mathrm{tanh}^2[\beta m_nM_𝝃]_𝝃n\mathrm{tanh}^2[\beta m_nM_𝝃]_𝝃\hfill \\ & =[M_𝝃^2M_𝝃^{}^2_𝝃^{}]\mathrm{tanh}^2[\beta m_n\|M_𝝃\|]_𝝃\hfill \\ & =[M_𝝃^2M_𝝃^{}^2_𝝃^{}]\left\{\mathrm{tanh}^2[\beta m_n\sqrt{M_𝝃^2}]\mathrm{tanh}^2[\beta m_n\sqrt{M_𝝃^{}^2_𝝃^{}}]\right\}_𝝃0\hfill \end{array}$$ We may now identify the conditions for an $`n`$-mixture state to be a local minimum of $`f(𝒎)`$. For $`n=1`$ the relevant eigenvalue is $`I`$, now the quantity $`Q`$ simplifies considerably. For $`n>1`$ the relevant eigenvalue is $`III`$, here we can combine $`Q`$ and $`R`$ into one single average: $$\begin{array}{cc}n=1:\hfill & 1\beta [1\mathrm{tanh}^2[\beta m_1]]>0\hfill \\ n=2:\hfill & 1\beta >0\hfill \\ n3:\hfill & 1\beta [1\mathrm{tanh}^2[\beta m_n_{\rho =3}^n\xi _\rho ]_𝝃]>0\hfill \end{array}$$ The $`n=1`$ states, correlated with one pattern only, are the desired solutions. They are stable for all $`T<1`$, since partial differentiation with respect to $`\beta `$ of the $`n=1`$ amplitude equation (53) gives $$m_1=\mathrm{tanh}[\beta m_1]1\beta [1\mathrm{tanh}^2[\beta m_1]]=m_1[1\mathrm{tanh}^2[\beta m_1]](m_1/\beta )^1$$ (clearly $`\mathrm{sgn}[m_1]=\mathrm{sgn}[m_1/\beta ]`$). The $`n=2`$ mixtures are always unstable. For $`n3`$ we have to solve the amplitude equations (53) numerically to evaluate their stability. The result is shown in figure 7, together with the corresponding ‘free energies’ $`f_n`$ (54). It turns out that only for odd $`n`$ will there be a critical temperature below which the $`n`$-mixture states are local minima of $`f(𝒎)`$. From figure 7 we can also conclude that, in terms of the network functioning as an associative memory, noise is actually beneficial in the sense that it can be used to eliminate the unwanted $`n>1`$ ergodic components (while retaining the relevant ones: the pure $`n=1`$ states). In fact the overlap equations (50) do also allow for stable solutions different from the $`n`$-mixture states discussed here. They are in turn found to be continuously bifurcating mixtures of the mixture states. However, for random (or uncorrelated) patterns they come into existence only near $`T=0`$ and play a marginal role; phase space is dominated by the odd $`n`$-mixture states. We have now solved the model in equilibrium for finite $`p`$ and $`N\mathrm{}`$. Most of the relevant information on when and to what extent stored random patterns will be recalled is summarised in figure 7. For non-random patterns one simply has to study the bifurcation properties of equation (50) for the new pattern statistics at hand; this is only qualitatively different from the random pattern analysis explained above. The occurrence of multiple saddle-points corresponding to local minima of the free energy signals ergodicity breaking. Although among these only the global minimum will correspond to the thermodynamic equilibrium state, the non-global minima correspond to true ergodic components, i.e. on finite time-scales they will be just as relevant as the global minimum. ## 4 Simple Recurrent Networks of Coupled Oscillators ### 4.1 Coupled Oscillators with Uniform Synapses Models with continuous variables involve integration over states, rather than summation. For a coupled oscillator network (13) with uniform synapses $`J_{ij}=J/N`$ and zero frequencies $`\omega _i=0`$ (which is a simple version of the model in ) we obtain for the free energy per oscillator: $$\underset{N\mathrm{}}{lim}F/N=\underset{N\mathrm{}}{lim}\frac{1}{\beta N}\mathrm{log}_\pi ^\pi \mathrm{}_\pi ^\pi 𝑑\mathit{\varphi }e^{(\beta J/2N)\left[[_i\mathrm{cos}(\varphi _i)]^2+[_i\mathrm{sin}(\varphi _i)]^2\right]}$$ We would now have to ‘count’ microscopic states with prescribed average cosines and sines. A faster route exploits auxiliary Gaussian integrals, via the identity $$e^{\frac{1}{2}y^2}=Dze^{yz}$$ (56) with the short-hand $`Dx=(2\pi )^{\frac{1}{2}}e^{\frac{1}{2}x^2}dx`$ (this alternative would also have been open to us in the binary case; my aim in this section is to explain both methods): $$\underset{N\mathrm{}}{lim}F/N=\underset{N\mathrm{}}{lim}\frac{1}{\beta N}\mathrm{log}_\pi ^\pi \mathrm{}_\pi ^\pi 𝑑\mathit{\varphi }DxDye^{\sqrt{\beta J/N}\left[x_i\mathrm{cos}(\varphi _i)+y_i\mathrm{sin}(\varphi _i)\right]}$$ $$=\underset{N\mathrm{}}{lim}\frac{1}{\beta N}\mathrm{log}DxDy\left[_\pi ^\pi 𝑑\varphi e^{\mathrm{cos}(\varphi )\sqrt{\beta J(x^2+y^2)/N}}\right]^N$$ $$=\underset{N\mathrm{}}{lim}\frac{1}{\beta N}\mathrm{log}_0^{\mathrm{}}𝑑qqe^{\frac{1}{2}N\beta \|J\|q^2}\left[_\pi ^\pi 𝑑\varphi e^{\beta \|J\|q\mathrm{cos}(\varphi )\sqrt{\mathrm{sgn}(J)}}\right]^N$$ where we have transformed to polar coordinates, $`(x,y)=q\sqrt{\beta \|J\|N}(\mathrm{cos}\theta ,\mathrm{sin}\theta )`$, and where we have already eliminated (constant) terms which will not survive the limit $`N\mathrm{}`$. Thus, saddle-point integration gives us, quite similar to the previous cases (36,37): $$\underset{N\mathrm{}}{lim}F/N=\underset{q0}{\mathrm{min}}f(q)\begin{array}{ccc}J>0:\hfill & & \beta f(q)=\frac{1}{2}\beta \|J\|q^2\mathrm{log}[2\pi I_0(\beta \|J\|q)]\hfill \\ J<0:\hfill & & \beta f(q)=\frac{1}{2}\beta \|J\|q^2\mathrm{log}[2\pi I_0(i\beta \|J\|q)]\hfill \end{array}$$ (57) in which the $`I_n(z)`$ are the Bessel functions (see e.g. ). The function $`f(q)`$ is shown in figure 8. The equations from which to solve the minima are obtained by differentiation, using $`\frac{d}{dz}I_0(z)=I_1(z)`$: $$J>0:q=\frac{I_1(\beta \|J\|q)}{I_0(\beta \|J\|q)}J<0:q=i\frac{I_1(i\beta \|J\|q)}{I_0(i\beta \|J\|q)}$$ (58) Again, in both cases the problem has been reduced to studying a single non-linear equation. The physical meaning of the solution follows from the identity $`2F/J=N^1_{ij}\mathrm{cos}(\varphi _i\varphi _j)`$: $$\underset{N\mathrm{}}{lim}[\frac{1}{N}\underset{i}{}\mathrm{cos}(\varphi _i)]^2+\underset{N\mathrm{}}{lim}[\frac{1}{N}\underset{i}{}\mathrm{sin}(\varphi _i)]^2=\mathrm{sgn}(J)q^2$$ From this equation it also follows that $`q1`$. Note: since $`f(q)/q=0`$ at the minimum, one only needs to consider the explicit derivative of $`f(q)`$ with respect to $`J`$. If the synapses induce anti-synchronisation, $`J<0`$, the only solution of (58) (and the minimum in (57)) is the trivial state $`q=0`$. This also follows immediately from the equation which gave the physical meaning of $`q`$. For synchronising forces, $`J>0`$, on the other hand, we again find the trivial solution at high noise levels, but a globally synchronised state with $`q>0`$ at low noise levels. Here a phase transition occurs at $`T=\frac{1}{2}J`$ (a bifurcation of non-trivial solutions of (58)), and for $`T<\frac{1}{2}J`$ the minimum of (57) is found at two non-zero values for $`q`$. The critical noise level is again found upon expanding the saddle-point equation, using $`I_0(z)=1+𝒪(z^2)`$ and $`I_1(z)=\frac{1}{2}z+𝒪(z^3)`$: $`q=\frac{1}{2}\beta Jq+𝒪(q^3)`$. Precisely at $`\beta J=2`$ one finds a de-stabilisation of the trivial solution $`q=0`$, together with the creation of (two) stable non-trivial ones (see figure 8). Note that, in view of (57), we are only interested in non-negative values of $`q`$. One can prove, using the properties of the Bessel functions, that there are no other (discontinuous) bifurcations of non-trivial solutions of the saddle-point equation. Note, finally, that the absence of a state with global anti-synchronisation for $`J<0`$ has the same origin as the absence of an anti-ferromagnetic state for $`J<0`$ in the previous models with binary neurons. Due to the long-range nature of the synapses $`J_{ij}=J/N`$ such states simply cannot exist: whereas any set of oscillators can be in a fully synchronised state, if two oscillators are in anti-synchrony it is already impossible for a third to be simultaneously in anti-synchrony with the first two (since anti-synchrony with one implies synchrony with the other). ### 4.2 Coupled Oscillator Attractor Networks Intuition & Definitions. Let us now turn to an alternative realisation of information storage in a recurrent network based upon the creation of attractors. We will solve models of coupled neural oscillators of the type (13), with zero natural frequencies (since we wish to use equilibrium techniques), in which real-valued patterns are stored as stable configurations of oscillator phases, following . Let us, however, first find out how to store a single pattern $`𝝃[\pi ,\pi ]^N`$ in a noise-less infinite-range oscillator network. For simplicity we will draw each component $`\xi _i`$ independently at random from $`[\pi ,\pi ]`$, with uniform probability density. This allows us to use asymptotic properties such as $`\|N^1_je^{i\mathrm{}\xi _j}\|=𝒪(N^{\frac{1}{2}})`$ for any integer $`\mathrm{}`$. A sensible choice for the synapses would be $`J_{ij}=\mathrm{cos}[\xi _i\xi _j]`$. To see this we work out the corresponding Lyapunov function (20): $$L[\mathit{\varphi }]=\frac{1}{2N^2}\underset{ij}{}\mathrm{cos}[\xi _i\xi _j]\mathrm{cos}[\varphi _i\varphi _j]L[𝝃]=\frac{1}{2N^2}\underset{ij}{}\mathrm{cos}^2[\xi _i\xi _j]=\frac{1}{4}+𝒪(N^{\frac{1}{2}})$$ (the factors of $`N`$ have been inserted to achieve appropriate scaling in the $`N\mathrm{}`$ limit). The function $`L[\mathit{\varphi }]`$, which is obviously bounded from below, must decrease monotonically during the dynamics. To find out whether the state $`𝝃`$ is a stable fixed-point of the dynamics we have to calculate $`L`$ and derivatives of $`L`$ at $`\mathit{\varphi }=𝝃`$: $$\frac{L}{\varphi _i}\|_𝝃=\frac{1}{2N^2}\underset{j}{}\mathrm{sin}[2(\xi _i\xi _j)]\frac{^2L}{\varphi _i^2}\|_𝝃=\frac{1}{N^2}\underset{j}{}\mathrm{cos}^2[\xi _i\xi _j]ij:\frac{^2L}{\varphi _i\varphi _j}\|_𝝃=\frac{1}{N^2}\mathrm{cos}^2[\xi _i\xi _j]$$ Clearly $`lim_N\mathrm{}L[𝝃]=\frac{1}{4}`$. Putting $`\mathit{\varphi }=𝝃+\mathrm{\Delta }\mathit{\varphi }`$, with $`\mathrm{\Delta }\varphi _i=𝒪(N^0)`$, we find $$L[𝝃+\mathrm{\Delta }\mathit{\varphi }]L[𝝃]=\underset{i}{}\mathrm{\Delta }\varphi _i\frac{L}{\varphi _i}\|_𝝃+\frac{1}{2}\underset{ij}{}\mathrm{\Delta }\varphi _i\mathrm{\Delta }\varphi _j\frac{^2L}{\varphi _i\varphi _j}\|_𝝃+𝒪(\mathrm{\Delta }\mathit{\varphi }^3)$$ $$=\frac{1}{4N}\underset{i}{}\mathrm{\Delta }\varphi _i^2\frac{1}{2N^2}\underset{ij}{}\mathrm{\Delta }\varphi _i\mathrm{\Delta }\varphi _j\mathrm{cos}^2[\xi _i\xi _j]+𝒪(N^{\frac{1}{2}},\mathrm{\Delta }\mathit{\varphi }^3)$$ $$=\frac{1}{4}\left\{\frac{1}{N}\underset{i}{}\mathrm{\Delta }\varphi _i^2[\frac{1}{N}\underset{i}{}\mathrm{\Delta }\varphi _i]^2[\frac{1}{N}\underset{i}{}\mathrm{\Delta }\varphi _i\mathrm{cos}(2\xi _i)]^2[\frac{1}{N}\underset{i}{}\mathrm{\Delta }\varphi _i\mathrm{sin}(2\xi _i)]^2\right\}+𝒪(N^{\frac{1}{2}},\mathrm{\Delta }\mathit{\varphi }^3)$$ (59) In leading order in $`N`$ the following three vectors in $`\mathrm{}^N`$ are normalised and orthogonal: $$𝐞_1=\frac{1}{\sqrt{N}}(1,1,\mathrm{},1),𝐞_2=\frac{\sqrt{2}}{\sqrt{N}}(\mathrm{cos}(2\xi _1),\mathrm{},\mathrm{cos}(2\xi _N)),𝐞_2=\frac{\sqrt{2}}{\sqrt{N}}(\mathrm{sin}(2\xi _1),\mathrm{},\mathrm{sin}(2\xi _N))$$ We may therefore use $`\mathrm{\Delta }\mathit{\varphi }^2(\mathrm{\Delta }\mathit{\varphi }𝐞_1)^2+(\mathrm{\Delta }\mathit{\varphi }𝐞_2)^2+(\mathrm{\Delta }\mathit{\varphi }𝐞_3)^2`$, insertion of which into (59) leads to $$L[𝝃+\mathrm{\Delta }\mathit{\varphi }]L[𝝃][\frac{1}{2N}\underset{i}{}\mathrm{\Delta }\varphi _i\mathrm{cos}(2\xi _i)]^2+[\frac{1}{2N}\underset{i}{}\mathrm{\Delta }\varphi _i\mathrm{sin}(2\xi _i)]^2+𝒪(N^{\frac{1}{2}},\mathrm{\Delta }\mathit{\varphi }^3)$$ Thus for large $`N`$ the second derivative of $`L`$ is non-negative at $`\mathit{\varphi }=𝝃`$, and the phase pattern $`𝝃`$ has indeed become a fixed-point attractor of the dynamics of the noise-free coupled oscillator network. The same is found to be true for the states $`\mathit{\varphi }=\pm 𝝃+\alpha (1,\mathrm{},1)`$ (for any $`\alpha `$). Storing $`p`$ Phase Patterns: Equilibrium Order Parameter Equations. We next follow the strategy of the Hopfield model and attempt to simply extend the above recipe for the synapses to the case of having a finite number $`p`$ of phase patterns $`𝝃^\mu =(\xi _1^\mu ,\mathrm{},\xi _N^\mu )[\pi ,\pi ]^N`$, giving $$J_{ij}=\frac{1}{N}\underset{\mu =1}{\overset{p}{}}\mathrm{cos}[\xi _i^\mu \xi _j^\mu ]$$ (60) (the factor $`N`$, as before, ensures a proper limit $`N\mathrm{}`$ later). In analogy with our solution of the Hopfield model we define the following averages over pattern variables: $$g[𝝃]_𝝃=\underset{N\mathrm{}}{lim}\underset{i}{}g[𝝃_i],𝝃_i=(\xi _i^1,\mathrm{},\xi _i^p)[\pi ,\pi ]^p$$ We can write the Hamiltonian $`H(\mathit{\varphi })`$ of (34) in the form $$H(\mathit{\varphi })=\frac{1}{2N}\underset{\mu =1}{\overset{p}{}}\underset{ij}{}\mathrm{cos}[\xi _i^\mu \xi _j^\mu ]\mathrm{cos}[\varphi _i\varphi _j]=\frac{N}{2}\underset{\mu =1}{\overset{p}{}}\left\{m_{cc}^\mu (\mathit{\varphi })^2+m_{cs}^\mu (\mathit{\varphi })^2+m_{sc}^\mu (\mathit{\varphi })^2+m_{ss}^\mu (\mathit{\varphi })^2\right\}$$ in which $$m_{cc}^\mu (\mathit{\varphi })=\frac{1}{N}\underset{i}{}\mathrm{cos}(\xi _i^\mu )\mathrm{cos}(\varphi _i)m_{cs}^\mu (\mathit{\varphi })=\frac{1}{N}\underset{i}{}\mathrm{cos}(\xi _i^\mu )\mathrm{sin}(\varphi _i)$$ (61) $$m_{sc}^\mu (\mathit{\varphi })=\frac{1}{N}\underset{i}{}\mathrm{sin}(\xi _i^\mu )\mathrm{cos}(\varphi _i)m_{ss}^\mu (\mathit{\varphi })=\frac{1}{N}\underset{i}{}\mathrm{sin}(\xi _i^\mu )\mathrm{sin}(\varphi _i)$$ (62) The free energy per oscillator can now be written as $$F/N=\frac{1}{\beta N}\mathrm{log}\mathrm{}𝑑\mathit{\varphi }e^{\beta H(\mathit{\varphi })}=\frac{1}{\beta N}\mathrm{log}\mathrm{}𝑑\mathit{\varphi }e^{\frac{1}{2}\beta N_\mu _{}m_{}^\mu (\mathit{\varphi })^2}$$ with $`\{cc,ss,cs,sc\}`$. Upon introducing the notation $`𝒎_{}=(m_{}^1,\mathrm{},m_{}^p)`$ we can again express the free energy in terms of the density of states $`𝒟(\{𝒎_{}\})=(2\pi )^N\mathrm{}𝑑\mathit{\varphi }_{}\delta [𝒎_{}𝒎_{}(𝝈)]`$: $$F/N=\frac{1}{\beta }\mathrm{log}(2\pi )\frac{1}{\beta N}\mathrm{log}\underset{}{}d𝒎_{}𝒟(\{𝒎_{}\})e^{\frac{1}{2}\beta N_{}𝒎_{}^2}$$ (63) Since $`p`$ is finite, the leading contribution to the density of states (as $`N\mathrm{}`$), which will give us the entropy, can be calculated by writing the $`\delta `$-functions in integral representation: $$\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}𝒟(\{𝒎_{}\})=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}\underset{}{}\left[d𝒙_{}e^{iN𝒙_{}𝒎_{}}\right]\times $$ $$\mathrm{}\frac{d\mathit{\varphi }}{(2\pi )^N}e^{i_i_\mu [x_{cc}^\mu \mathrm{cos}(\xi _i^\mu )\mathrm{cos}(\varphi _i)+x_{cs}^\mu \mathrm{cos}(\xi _i^\mu )\mathrm{sin}(\varphi _i)+x_{sc}^\mu \mathrm{sin}(\xi _i^\mu )\mathrm{cos}(\varphi _i)+x_{ss}^\mu \mathrm{sin}(\xi _i^\mu )\mathrm{sin}(\varphi _i)]}$$ $$=\mathrm{extr}_{\{𝒙_{}\}}\left\{i\underset{}{}𝒙_{}𝒎_{}+\mathrm{log}\frac{d\varphi }{2\pi }e^{i_\mu [x_{cc}^\mu \mathrm{cos}(\xi _\mu )\mathrm{cos}(\varphi )+x_{cs}^\mu \mathrm{cos}(\xi _\mu )\mathrm{sin}(\varphi )+x_{sc}^\mu \mathrm{sin}(\xi _\mu )\mathrm{cos}(\varphi )+x_{ss}^\mu \mathrm{sin}(\xi _\mu )\mathrm{sin}(\varphi )]}_𝝃\right\}$$ The relevant extremum is purely imaginary so we put $`𝒙_{}=i\beta 𝒚_{}`$ (see also our previous discussion for the Hopfield model) and, upon inserting the density of states into our original expression for the free energy per oscillator, arrive at $$\underset{N\mathrm{}}{lim}F/N=\mathrm{extr}_{\{𝒎_{},𝒚_{}\}}f(\{𝒎_{},𝒚_{}\})$$ $$f(\{𝒎_{},𝒚_{}\})=\frac{1}{\beta }\mathrm{log}(2\pi )\frac{1}{2}\underset{}{}𝒎_{}^2+\underset{}{}𝒚_{}𝒎_{}$$ $$\frac{1}{\beta }\mathrm{log}\frac{d\varphi }{2\pi }e^{\beta _\mu [y_{cc}^\mu \mathrm{cos}(\xi _\mu )\mathrm{cos}(\varphi )+y_{cs}^\mu \mathrm{cos}(\xi _\mu )\mathrm{sin}(\varphi )+y_{sc}^\mu \mathrm{sin}(\xi _\mu )\mathrm{cos}(\varphi )+y_{ss}^\mu \mathrm{sin}(\xi _\mu )\mathrm{sin}(\varphi )]}_𝝃$$ Taking derivatives with respect to the order parameters $`𝒎_{}`$ gives us $`𝒚_{}=𝒎_{}`$, with which we can eliminate the $`𝒚_{}`$. Derivation with respect to the $`𝒎_{}`$ subsequently gives the saddle-point equations $$m_{cc}^\mu =\mathrm{cos}[\xi _\mu ]\frac{𝑑\varphi \mathrm{cos}[\varphi ]e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}}{𝑑\varphi e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}}_𝝃$$ (64) $$m_{cs}^\mu =\mathrm{cos}[\xi _\mu ]\frac{𝑑\varphi \mathrm{sin}[\varphi ]e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}}{𝑑\varphi e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}}_𝝃$$ (65) $$m_{sc}^\mu =\mathrm{sin}[\xi _\mu ]\frac{𝑑\varphi \mathrm{cos}[\varphi ]e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}}{𝑑\varphi e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}}_𝝃$$ (66) $$m_{ss}^\mu =\mathrm{sin}[\xi _\mu ]\frac{𝑑\varphi \mathrm{sin}[\varphi ]e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}}{𝑑\varphi e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}}_𝝃$$ (67) The equilibrium values of the observables $`𝒎_{}`$, as defined in (61,62), are now given by the solution of the coupled equations (64-67) which minimises $$f(\{𝒎_{}\})=\frac{1}{2}\underset{}{}𝒎_{}^2\frac{1}{\beta }\mathrm{log}𝑑\varphi e^{\beta \mathrm{cos}[\varphi ]_\nu [m_{cc}^\nu \mathrm{cos}[\xi _\nu ]+m_{sc}^\nu \mathrm{sin}[\xi _\nu ]]+\beta \mathrm{sin}[\varphi ]_\nu [m_{cs}^\nu \mathrm{cos}[\xi _\nu ]+m_{ss}^\nu \mathrm{sin}[\xi _\nu ]]}_𝝃$$ (68) We can confirm that the relevant saddle-point must be a minimum by inspecting the $`\beta =0`$ limit (infinite noise levels): $`lim_{\beta 0}f(\{𝒎_{}\})=\frac{1}{2}_{}𝒎_{}^2\frac{1}{\beta }\mathrm{log}(2\pi )`$. Analysis of Order Parameter Equations: Pure States. From now on we will restrict our analysis to phase pattern components $`\xi _i^\mu `$ which have all been drawn independently at random from $`[\pi ,\pi ]`$, with uniform probability density, so that $`g[𝝃]_𝝃=(2\pi )^p_\pi ^\pi \mathrm{}_\pi ^\pi 𝑑𝝃g[𝝃]`$. At $`\beta =0`$ ($`T=\mathrm{}`$) one finds only the trivial state $`m_{}^\mu =0`$. It can be shown that there will be no discontinuous transitions to a non-trivial state as the noise level (temperature) is reduced. The continuous ones follow upon expansion of the equations (64-67) for small $`\{𝒎_{}\}`$, which is found to give (for each $`\mu `$ and each combination $``$): $$m_{}^\mu =\frac{1}{4}\beta m_{}^\mu +𝒪(\{𝒎_{}^2\})$$ Thus a continuous transition to recall states occurs at $`T=\frac{1}{4}`$. Full classification of all solutions of (64-67) is ruled out. Here we will restrict ourselves to the most relevant ones, such as the pure states, where $`m_{}^\mu =m_{}\delta _{\mu \lambda }`$ (for some pattern label $`\lambda `$). Here the oscillator phases are correlated with only one of the stored phase patterns (if at all). Insertion into the above expression for $`f(\{𝒎_{}\})`$ shows that for such solutions we have to minimise $$f(\{m_{}\})=\frac{1}{2}\underset{}{}m_{}^2\frac{1}{\beta }\frac{d\xi }{2\pi }\mathrm{log}𝑑\varphi e^{\beta \mathrm{cos}[\varphi ][m_{cc}\mathrm{cos}[\xi ]+m_{sc}\mathrm{sin}[\xi ]]+\beta \mathrm{sin}[\varphi ][m_{cs}\mathrm{cos}[\xi ]+m_{ss}\mathrm{sin}[\xi ]]}$$ (69) We anticipate solutions corresponding to the (partial) recall of the stored phase pattern $`𝝃^\lambda `$ or its mirror image (modulo overall phase shifts $`\xi _i\xi _i+\delta `$, under which the synapses are obviously invariant). Insertion into (64-67) of the state $`\varphi _i=\xi _i^\lambda +\delta `$ gives $`(m_{cc},m_{sc},m_{cs},m_{ss})=\frac{1}{2}(\mathrm{cos}\delta ,\mathrm{sin}\delta ,\mathrm{sin}\delta ,\mathrm{cos}\delta )`$. Similarly, insertion into (64-67) of $`\varphi _i=\xi _i^\lambda +\delta `$ gives $`(m_{cc},m_{sc},m_{cs},m_{ss})=\frac{1}{2}(\mathrm{cos}\delta ,\mathrm{sin}\delta ,\mathrm{sin}\delta ,\mathrm{cos}\delta )`$. Thus we can identify retrieval states as those solutions which are of the form $$\begin{array}{ccccc}(i)\hfill & & \mathrm{retrieval}\mathrm{of}𝝃^\lambda :\hfill & & (m_{cc},m_{sc},m_{cs},m_{ss})=m(\mathrm{cos}\delta ,\mathrm{sin}\delta ,\mathrm{sin}\delta ,\mathrm{cos}\delta )\hfill \\ (ii)\hfill & & \mathrm{retrieval}\mathrm{of}𝝃^\lambda :\hfill & & (m_{cc},m_{sc},m_{cs},m_{ss})=m(\mathrm{cos}\delta ,\mathrm{sin}\delta ,\mathrm{sin}\delta ,\mathrm{cos}\delta )\hfill \end{array}$$ with full recall corresponding to $`m=\frac{1}{2}`$. Insertion into the saddle-point equations and into (69), followed by an appropriate shift of the integration variable $`\varphi `$, shows that the free energy is independent of $`\delta `$ (so the above two ansätze solve the saddle-point equations for any $`\delta `$) and that $$m=\frac{1}{2}\frac{𝑑\varphi \mathrm{cos}[\varphi ]e^{\beta m\mathrm{cos}[\varphi ]}}{𝑑\varphi e^{\beta m\mathrm{cos}[\varphi ]}},f(m)=m^2\frac{1}{\beta }\mathrm{log}𝑑\varphi e^{\beta m\mathrm{cos}[\varphi ]}$$ Expansion in powers of $`m`$, using $`\mathrm{log}(1+z)=z\frac{1}{2}z^2+𝒪(z^3)`$, reveals that non-zero minima $`m`$ indeed bifurcate continuously at $`T=\beta ^1=\frac{1}{4}`$: $$f(m)+\frac{1}{\beta }\mathrm{log}[2\pi ]=(1\frac{1}{4}\beta )m^2+\frac{1}{64}\beta ^3m^4+𝒪(m^6)$$ (70) Retrieval states are obviously not the only pure states that solve the saddle-point equations. The function (69) is invariant under the following discrete (non-commuting) transformations: $$\begin{array}{ccccc}\mathrm{I}:\hfill & & (m_{cc},m_{sc},m_{cs},m_{ss})\hfill & \hfill & (m_{cc},m_{sc},m_{cs},m_{ss})\hfill \\ \mathrm{II}:\hfill & & (m_{cc},m_{sc},m_{cs},m_{ss})\hfill & \hfill & (m_{cs},m_{ss},m_{cc},m_{sc})\hfill \end{array}$$ We expect these to induce solutions with specific symmetries. In particular we anticipate the following symmetric and anti-symmetric states: $$\begin{array}{ccccc}(iii)\hfill & & \mathrm{symmetric}\mathrm{under}\mathrm{I}:\hfill & & (m_{cc},m_{sc},m_{cs},m_{ss})=\sqrt{2}m(\mathrm{cos}\delta ,\mathrm{sin}\delta ,0,0)\hfill \\ (iv)\hfill & & \mathrm{antisymmetric}\mathrm{under}\mathrm{I}:\hfill & & (m_{cc},m_{sc},m_{cs},m_{ss})=\sqrt{2}m(0,0,\mathrm{cos}\delta ,\mathrm{sin}\delta )\hfill \\ (v)\hfill & & \mathrm{symmetric}\mathrm{under}\mathrm{II}:\hfill & & (m_{cc},m_{sc},m_{cs},m_{ss})=m(\mathrm{cos}\delta ,\mathrm{sin}\delta ,\mathrm{cos}\delta ,\mathrm{sin}\delta )\hfill \\ (vi)\hfill & & \mathrm{antisymmetric}\mathrm{under}\mathrm{II}:\hfill & & (m_{cc},m_{sc},m_{cs},m_{ss})=m(\mathrm{cos}\delta ,\mathrm{sin}\delta ,\mathrm{cos}\delta ,\mathrm{sin}\delta )\hfill \end{array}$$ Insertion into the saddle-point equations and into (69) shows in all four cases the parameter $`\delta `$ is arbitrary and that always $$m=\frac{1}{\sqrt{2}}\frac{d\xi }{2\pi }\mathrm{cos}[\xi ]\frac{𝑑\varphi \mathrm{cos}[\varphi ]e^{\beta m\sqrt{2}\mathrm{cos}[\varphi ]\mathrm{cos}[\xi ]}}{𝑑\varphi e^{\beta m\sqrt{2}\mathrm{cos}[\varphi ]\mathrm{cos}[\xi ]}},f(m)=m^2\frac{1}{\beta }\frac{d\xi }{2\pi }\mathrm{log}𝑑\varphi e^{\beta m\sqrt{2}\mathrm{cos}[\varphi ]\mathrm{cos}[\xi ]}$$ Expansion in powers of $`m`$ reveals that non-zero solutions $`m`$ here again bifurcate continuously at $`T=\frac{1}{4}`$: $$f(m)+\frac{1}{\beta }\mathrm{log}[2\pi ]=(1\frac{1}{4}\beta )m^2+\frac{3}{2}.\frac{1}{64}\beta ^3m^4+𝒪(m^6)$$ (71) However, comparison with (70) shows that the free energy of the pure recall states is lower. Thus the system will prefer the recall states over the above solutions with specific symmetries. Note, finally, that the free energy and the order parameter equation for the pure recall states can be written in terms of Bessel functions as follows: $$m=\frac{1}{2}\frac{I_1(\beta m)}{I_0(\beta m)},f(m)=m^2\frac{1}{\beta }\mathrm{log}[2\pi I_0(\beta m)]$$ The behaviour of these equations and the observable $`m`$ for different noise levels is shown in figure 9. One easily proves that $`\|m\|\frac{1}{2}`$, and that $`lim_\beta \mathrm{}m=\frac{1}{2}`$. Following the transition to a state with partial recall of a stored phase pattern at $`T=\frac{1}{4}`$, further reduction of the noise level $`T`$ gives a monotonic increase of retrieval quality until retrieval is perfect at $`T=0`$. ## 5 Networks with Gaussian Distributed Synapses The type of analysis presented so far to deal with attractor networks breaks down if the number of patterns stored $`p`$ no longer remains finite for $`N\mathrm{}`$, but scales as $`p=\alpha N`$ $`(\alpha >0`$). Expressions such as (48,49) can no longer be evaluated by saddle-point methods, since the dimension of the integral diverges at the same time as the exponent of the integrand. The number of local minima (ergodic components) of Hamiltonians such as (30,32) will diverge and we will encounter phenomena reminiscent of complex disordered magnetic systems, i.e. spin-glasses. As a consequence we will need corresponding methods of analysis, in the present case: replica theory. ### 5.1 Replica Analysis Replica Calculation of the Disorder-Averaged Free Energy. As an introduction to the replica technique we will first discuss the equilibrium solution of a recurrent neural network model with binary neurons $`\sigma _i\{1,1\}`$ in which a single pattern $`𝝃=(\xi _1,\mathrm{},\xi _N)\{1,1\}^N`$ has been stored (via a Hebbian-type recipe) on a background of zero-average Gaussian synapses (equivalent to the SK model, ): $$J_{ij}=\frac{J_0}{N}\xi _i\xi _j+\frac{J}{\sqrt{N}}z_{ij},\overline{z}_{ij}=0,\overline{z^2}_{ij}=1$$ (72) in which $`J_0>0`$ measures the embedding strength of the pattern, and the $`z_{ij}`$ ($`i<j`$) are independent Gaussian random variables. We denote averaging over their distribution by $`\overline{\mathrm{}}`$ (the factors in (72) involving $`N`$ ensure appropriate scaling and statistical relevance of the two terms, and as always $`J_{ii}=0`$). Here the Hamiltonian $`H`$ (30), corresponding to sequential dynamics (2), becomes $$H(𝝈)=\frac{1}{2}NJ_0m^2(𝝈)+\frac{1}{2}J_0\frac{J}{\sqrt{N}}\underset{i<j}{}\sigma _i\sigma _jz_{ij}$$ (73) with the overlap $`m(𝝈)=\frac{1}{N}_k\sigma _k\xi _k`$ which measures pattern recall quality. We clearly cannot calculate the free energy for every given realization of the synapses, furthermore it is to be expected that for $`N\mathrm{}`$ macroscopic observables like the free energy per neuron and the overlap $`m`$ only depend on the statistics of the synapses, not on their specific values. We therefore average the free energy over the disorder distribution and concentrate on $$\overline{F}=\frac{1}{\beta }\underset{N\mathrm{}}{lim}\overline{\mathrm{log}Z},Z=\underset{𝝈}{}e^{\beta H(𝝈)}$$ (74) The disorder average is transformed into an average of powers of $`Z`$, with the identity $$\overline{\mathrm{log}Z}=\underset{n0}{lim}\frac{1}{n}\left[\overline{Z^n}1\right]\mathrm{or},\mathrm{equivalently},\overline{\mathrm{log}Z}=\underset{n0}{lim}\frac{1}{n}\mathrm{log}\overline{Z^n}$$ (75) The so-called ‘replica trick’ consists in evaluating the averages $`\overline{Z^n}`$ for integer values of $`n`$, and taking the limit $`n0`$ afterwards, under the assumption that the resulting expression is correct for non-integer values of $`n`$ as well. The integer powers of $`Z`$ are written as a product of terms, each of which can be interpreted as an equivalent copy, or ‘replica’ of the original system. The disorder-averaged free energy now becomes $$\overline{F}=\underset{n0}{lim}\frac{1}{\beta n}\mathrm{log}\overline{Z^n}=\underset{n0}{lim}\frac{1}{\beta n}\mathrm{log}\underset{𝝈^1\mathrm{}𝝈^n}{}\overline{e^{\beta _{\alpha =1}^nH(𝝈^\alpha )}}$$ From now Roman indices will refer to sites, i.e. $`i=1\mathrm{}N`$, whereas Greek indices will refer to replicas, i.e. $`\alpha =1\mathrm{}n`$. We introduce a short-hand for the Gaussian measure, $`Dz=(2\pi )^{\frac{1}{2}}e^{\frac{1}{2}z^2}dz`$, and we will repeatedly use the identity $`Dze^{xz}=e^{\frac{1}{2}x^2}`$. Upon insertion of the Hamiltonian (73) we obtain $$\begin{array}{cc}\overline{F}\hfill & =\frac{1}{\beta }N\mathrm{log}2lim_{n0}(\beta n)^1\mathrm{log}e^{\frac{\beta J_0}{N}_{i<j}\xi _i\xi _j_\alpha \sigma _i^\alpha \sigma _j^\alpha }_{i<j}\left[Dze^{\frac{\beta Jz}{\sqrt{N}}_\alpha \sigma _i^\alpha \sigma _j^\alpha }\right]_{\{𝝈^\alpha \}}\text{}\hfill \\ & \\ & =\frac{1}{\beta }N\mathrm{log}2lim_{n0}(\beta n)^1\mathrm{log}e^{\frac{\beta J_0}{2N}_\alpha _{ij}\xi _i\xi _j\sigma _i^\alpha \sigma _j^\alpha +\frac{\beta ^2J^2}{4N}_{\alpha \gamma }_{ij}\sigma _i^\alpha \sigma _j^\alpha \sigma _i^\gamma \sigma _j^\gamma }_{\{𝝈^\alpha \}}\text{}\hfill \end{array}$$ We now complete the sums over sites in this expression, $$\underset{ij}{}\sigma _i^\alpha \sigma _j^\alpha =[\underset{i}{}\sigma _i^\alpha ]^2N,\underset{ij}{}\sigma _i^\alpha \sigma _j^\alpha \sigma _i^\gamma \sigma _j^\gamma =[\underset{i}{}\sigma _i^\alpha \sigma _i^\gamma ]^2N$$ The averaging over the neuron states $`\{𝝈^\alpha \}`$ in our expression for $`\overline{F}`$ will now factorize nicely if we insert appropriate $`\delta `$-functions (in their integral representations) to isolate the relevant terms, using $$1=𝑑𝒒\underset{\alpha \beta }{}\delta \left[q_{\alpha \beta }\frac{1}{N}\underset{i}{}\sigma _i^\alpha \sigma _i^\beta \right]=\left[\frac{N}{2\pi }\right]^{n^2}𝑑𝒒𝑑\widehat{𝒒}e^{iN_{\alpha \beta }\widehat{q}_{\alpha \beta }\left[q_{\alpha \beta }\frac{1}{N}_i\sigma _i^\alpha \sigma _i^\beta \right]}$$ $$1=𝑑𝒎\underset{\alpha }{}\delta \left[m_\alpha \frac{1}{N}\underset{i}{}\xi _i\sigma _i^\alpha \right]=\left[\frac{N}{2\pi }\right]^n𝑑𝒎𝑑\widehat{𝒎}e^{iN_\alpha \widehat{m}_\alpha \left[m_\alpha \frac{1}{N}_i\xi _i\sigma _i^\alpha \right]}$$ The integrations are over the $`n\times n`$ matrices $`𝒒`$ and $`\widehat{𝒒}`$ and over the $`n`$-vectors $`𝒎`$ and $`\widehat{𝒎}`$. After inserting these integrals we obtain $$\underset{N\mathrm{}}{lim}\overline{F}/N=\frac{1}{\beta }\mathrm{log}2\underset{N\mathrm{}}{lim}\underset{n0}{lim}\frac{1}{\beta Nn}\mathrm{log}\{\left[\frac{N}{2\pi }\right]^{n^2+n}d𝒒d\widehat{𝒒}d𝒎d\widehat{𝒎}e^{\frac{1}{2}n\beta J_0\frac{1}{4}n^2\beta ^2J^2}$$ $$\times e^{N\left[i_{\alpha \gamma }\widehat{q}_{\alpha \gamma }q_{\alpha \gamma }+i_\alpha \widehat{m}_\alpha m_\alpha +\frac{1}{2}\beta J_0_\alpha m_\alpha ^2+\frac{1}{4}\beta ^2J^2_{\alpha \gamma }q_{\alpha \gamma }^2\right]}e^{i_i\left[_{\alpha \gamma }\widehat{q}_{\alpha \gamma }\sigma _i^\alpha \sigma _i^\gamma +_\alpha \widehat{m}_\alpha \xi _i\sigma _i^\alpha \right]}_{\{𝝈^\alpha \}}\}$$ The neuronal averages factorise and are therefore reduced to single-site ones. A simple transformation $`\sigma _i\xi _i\sigma _i`$ for all $`i`$ eliminates the pattern components $`\xi _i`$ from our equations, and the remaining averages involve only one $`n`$-replicated neuron $`(\sigma _1,\mathrm{},\sigma _n)`$. Finally one assumes that the two limits $`n0`$ and $`N\mathrm{}`$ commute. This allows us to evaluate the integral with the steepest-descent method: $$\underset{N\mathrm{}}{lim}\underset{n0}{lim}\frac{1}{Nn}\mathrm{log}𝑑𝒙e^{N\mathrm{\Phi }(𝒙)}=\underset{n0}{lim}\underset{N\mathrm{}}{lim}\frac{1}{Nn}\mathrm{log}e^{N\mathrm{extr}\mathrm{\Phi }+\mathrm{}}\text{}=\underset{n0}{lim}\frac{1}{n}\mathrm{extr}\mathrm{\Phi }\text{}$$ (76) The result of these manipulations is $$\underset{N\mathrm{}}{lim}\overline{F}/N=\underset{n0}{lim}\mathrm{extr}f(𝒒,𝒎;\widehat{𝒒},\widehat{𝒎})\text{}$$ (77) $$f(𝒒,𝒎;\widehat{𝒒},\widehat{𝒎})=\frac{1}{\beta }\mathrm{log}2\frac{1}{\beta n}[\mathrm{log}e^{i_{\alpha \gamma }\widehat{q}_{\alpha \gamma }\sigma _\alpha \sigma _\gamma i_\alpha \widehat{m}_\alpha \sigma _\alpha }_𝝈$$ $$+i\underset{\alpha \gamma }{}\widehat{q}_{\alpha \gamma }q_{\alpha \gamma }+i\underset{\alpha }{}\widehat{m}_\alpha m_\alpha +\frac{1}{2}\beta J_0\underset{\alpha }{}m_\alpha ^2+\frac{1}{4}\beta ^2J^2\underset{\alpha \gamma }{}q_{\alpha \gamma }^2]$$ (78) Variation of the parameters $`\{q_{\alpha \beta }\}`$ and $`\{m_\alpha \}`$ allows us to eliminate immediately the conjugate parameters $`\{\widehat{q}_{\alpha \beta }\}`$ and $`\{\widehat{m}_\alpha \}`$, since it leads to the saddle-point requirements $$\widehat{q}_{\alpha \beta }=\frac{1}{2}i\beta ^2J^2q_{\alpha \beta }\widehat{m}_\alpha =i\beta J_0m_\alpha $$ (79) Upon elimination of $`\{\widehat{q}_{\alpha \beta },\widehat{m}_\alpha \}`$ according to (79) the result (77,78) is simplified to $$\underset{N\mathrm{}}{lim}\overline{F}/N=\underset{n0}{lim}\mathrm{extr}f(𝒒,𝒎)$$ (80) $$f(𝒒,𝒎)=\frac{1}{\beta }\mathrm{log}2+\frac{\beta J^2}{4n}\underset{\alpha \gamma }{}q_{\alpha \gamma }^2+\frac{J_0}{2n}\underset{\alpha }{}m_\alpha ^2\frac{1}{\beta n}\mathrm{log}e^{\frac{1}{2}\beta ^2J^2_{\alpha \gamma }q_{\alpha \gamma }\sigma _\alpha \sigma _\gamma +\beta J_0_\alpha m_\alpha \sigma _\alpha }_𝝈$$ (81) Variation of the remaining parameters $`\{q_{\alpha \beta }\}`$ and $`\{m_\alpha \}`$ gives the final saddle-point equations $$q_{\lambda \rho }=\frac{\sigma _\lambda \sigma _\rho e^{\frac{1}{2}\beta ^2J^2_{\alpha \gamma }q_{\alpha \gamma }\sigma _\alpha \sigma _\gamma +\beta J_0_\alpha m_\alpha \sigma _\alpha }_𝝈}{e^{\frac{1}{2}\beta ^2J^2_{\alpha \gamma }q_{\alpha \gamma }\sigma _\alpha \sigma _\gamma +\beta J_0_\alpha m_\alpha \sigma _\alpha }_𝝈}$$ (82) $$m_\lambda =\frac{\sigma _\lambda e^{\frac{1}{2}\beta ^2J^2_{\alpha \gamma }q_{\alpha \gamma }\sigma _\alpha \sigma _\gamma +\beta J_0_\alpha m_\alpha \sigma _\alpha }_𝝈}{e^{\frac{1}{2}\beta ^2J^2_{\alpha \gamma }q_{\alpha \gamma }\sigma _\alpha \sigma _\gamma +\beta J_0_\alpha m_\alpha \sigma _\alpha }_𝝈}$$ (83) The diagonal elements are always $`q_{\alpha \alpha }=1`$. For high noise levels, $`\beta 0`$, we obtain the trivial result $$q_{\alpha \gamma }=\delta _{\alpha \gamma },m_\alpha =0$$ Assuming a continuous transition to a non-trivial state as the noise level is lowered, we can expand the saddle-point equations (82,83) in powers of $`𝒒`$ and $`𝒎`$ and look for bifurcations, which gives ($`\lambda \rho `$): $$q_{\lambda \rho }=\beta ^2J^2q_{\lambda \rho }+𝒪(𝒒,𝒎)^2m_\lambda =\beta J_0m_\lambda +𝒪(𝒒,𝒎)^2\text{}$$ Therefore we expect transitions either at $`T=J_0`$ (if $`J_0>J`$) or at $`T=J`$ (if $`J>J_0`$). The remaining program is: find the saddle-point $`(𝒒,𝒎)`$ for $`T<\mathrm{max}\{J_0,J\}`$ which for integer $`n`$ minimises $`f`$, determine the corresponding minimum as a function of $`n`$, and finally take the limit $`n0`$. This is in fact the most complicated part of the procedure. ### 5.2 Replica-Symmetric Solution and AT-Instability Physical Interpretation of Saddle Points. To obtain a guide in how to select saddle-points we now turn to a different (but equivalent) version of the replica trick (75), which allows us to attach a physical meaning to the saddle-points $`(𝒎,𝒒)`$. This version transforms averages over a given measure $`W`$: $$\frac{_𝝈\mathrm{\Phi }(𝝈)W(𝝈)}{_𝝈W(𝝈)}=\underset{n0}{lim}\underset{𝝈}{}\mathrm{\Phi }(𝝈)W(𝝈)[\underset{𝝈}{}W(𝝈)]^{n1}=\underset{n0}{lim}\underset{𝝈^1\mathrm{}𝝈^n}{}\mathrm{\Phi }(𝝈^1)\underset{\alpha =1}{\overset{n}{}}W(𝝈^\alpha )$$ $$=\underset{n0}{lim}\frac{1}{n}\underset{\gamma =1}{\overset{n}{}}\underset{𝝈^1\mathrm{}𝝈^n}{}\mathrm{\Phi }(𝝈^\gamma )\underset{\alpha =1}{\overset{n}{}}W(𝝈^\alpha )$$ (84) The trick again consists in evaluating this quantity for integer $`n`$, whereas the limit refers to non-integer $`n`$. We use (84) to write the distribution $`P(m)`$ of overlaps in equilibrium as $$P(m)=\frac{_𝝈\delta [m\frac{1}{N}_i\xi _i\sigma _i]e^{\beta H(𝝈)}}{_𝝈e^{\beta H(𝝈)}}=\underset{n0}{lim}\frac{1}{n}\underset{\gamma }{}\underset{𝝈^1\mathrm{}𝝈^n}{}\delta [m\frac{1}{N}\underset{i}{}\xi _i\sigma _i^\gamma ]\underset{\alpha }{}e^{\beta H(𝝈^\alpha )}$$ If we average this distribution over the disorder, we find identical expressions to those encountered in evaluating the disorder averaged free energy. By inserting the same delta-functions we arrive at the steepest descend integration (77) and find $$\overline{P(m)}=\underset{n0}{lim}\frac{1}{n}\underset{\gamma }{}\delta \left[mm_\gamma \right]$$ (85) where $`\{m_\gamma \}`$ refers to the relevant solution of (82,83). Similarly we can imagine two systems $`𝝈`$ and $`𝝈^{}`$ with identical synapses $`\{J_{ij}\}`$, both in thermal equilibrium. We now use (84) to rewrite the distribution $`P(q)`$ for the mutual overlap between the microstates of the two systems $$P(q)=\frac{_{𝝈,𝝈^{}}\delta [q\frac{1}{N}_i\sigma _i\sigma _i^{}]e^{\beta H(𝝈)\beta H(𝝈^{})}}{_{𝝈,𝝈^{}}e^{\beta H(𝝈)\beta H(𝝈^{})}}$$ $$=\underset{n0}{lim}\frac{1}{n(n1)}\underset{\lambda \gamma }{}\underset{𝝈^1\mathrm{}𝝈^n}{}\delta [q\frac{1}{N}\underset{i}{}\sigma _i^\lambda \sigma _i^\gamma ]\underset{\alpha }{}e^{\beta H(𝝈^\alpha )}$$ Averaging over the disorder again leads to the steepest descend integration (77) and we find $$\overline{P(q)}=\underset{n0}{lim}\frac{1}{n(n1)}\underset{\lambda \gamma }{}\delta \left[qq_{\lambda \gamma }\right]$$ (86) where $`\{q_{\lambda \gamma }\}`$ refers to the relevant solution of (82,83). We can now partly interpret the saddle-points $`(𝒎,𝒒)`$, since the shape of $`\overline{P(q)}`$ and $`\overline{P(m)}`$ gives direct information on the structure of phase space with respect to ergodicity. The crucial observation is that for an ergodic system one always has $$P(m)=\delta [m\frac{1}{N}\underset{i}{}\xi _i\sigma _i_{\mathrm{eq}}]P(q)=\delta [q\frac{1}{N}\underset{i}{}\sigma _i_{\mathrm{eq}}^2]$$ (87) If, on the other hand, there are $`L`$ ergodic components in our system, each of which corresponding to a pure Gibbs state with microstate probabilities proportional to $`\mathrm{exp}(\beta H)`$ and thermal averages $`\mathrm{}_{\mathrm{}}`$, and if we denote the probability of finding the system in component $`\mathrm{}`$ by $`W_{\mathrm{}}`$, we find $$P(m)=\underset{\mathrm{}=1}{\overset{L}{}}W_{\mathrm{}}\delta [m\frac{1}{N}\underset{i}{}\xi _i\sigma _i_{\mathrm{}}]P(q)=\underset{\mathrm{},\mathrm{}^{}=1}{\overset{L}{}}W_{\mathrm{}}W_{\mathrm{}^{}}\delta [q\frac{1}{N}\underset{i}{}\sigma _i_{\mathrm{}}\sigma _i_{\mathrm{}^{}}]$$ For ergodic systems both $`P(m)`$ and $`P(q)`$ are $`\delta `$-functions, for systems with a finite number of ergodic components they are finite sums of $`\delta `$-functions. A diverging number of ergodic components generally leads to distributions with continuous pieces. If we combine this interpretation with our results (85,86) we find that ergodicity is equivalent to the relevant saddle-point being of the form: $$q_{\alpha \beta }=\delta _{\alpha \beta }+q\left[1\delta _{\alpha \beta }\right]m_\alpha =m$$ (88) which is called the ‘replica symmetry’ (RS) ansatz. The meaning of $`m`$ and $`q`$ is deduced from (87) (taking into account the transformation $`\sigma _i\xi _i\sigma _i`$ we performed along the way): $$m=\frac{1}{N}\underset{i}{}\overline{\xi _i\sigma _i_{\mathrm{eq}}}q=\frac{1}{N}\underset{i}{}\overline{\sigma _i_{\mathrm{eq}}^2}$$ Replica Symmetric Solution. Having saddle-points of the simple form (88) leads to an enormous simplification in our calculations. Insertion of (88) as an ansatz into the equations (81,82,83) gives $$f(𝒒,𝒎)=\frac{1}{\beta }\mathrm{log}2\frac{1}{4}\beta J^2(1q)^2+\frac{1}{2}J_0m^2\frac{1}{\beta n}\mathrm{log}e^{\frac{1}{2}q\beta ^2J^2\left[_\alpha \sigma _\alpha \right]^2+\beta J_0m_\alpha \sigma _\alpha }_𝝈+𝒪(n)$$ $$q=\frac{\sigma _1\sigma _2e^{\frac{1}{2}q\beta ^2J^2\left[_\alpha \sigma _\alpha \right]^2+\beta J_0m_\alpha \sigma _\alpha }_𝝈}{e^{\frac{1}{2}q\beta ^2J^2\left[_\alpha \sigma _\alpha \right]^2+\beta J_0m_\alpha \sigma _\alpha }_𝝈}m=\frac{\sigma _1e^{\frac{1}{2}q\beta ^2J^2\left[_\alpha \sigma _\alpha \right]^2+\beta J_0m_\alpha \sigma _\alpha }_𝝈}{e^{\frac{1}{2}q\beta ^2J^2\left[_\alpha \sigma _\alpha \right]^2+\beta J_0m_\alpha \sigma _\alpha }_𝝈}$$ We linearise the terms $`\left[_\alpha \sigma _\alpha \right]^2`$ by introducing a Gaussian integral, and perform the average over the remaining neurons. The solutions $`m`$ and $`q`$ turn out to be well defined for $`n0`$ so we can take the limit: $$\underset{n0}{lim}f(𝒒,𝒎)=\frac{1}{\beta }\mathrm{log}2\frac{1}{4}\beta J^2(1q)^2+\frac{1}{2}J_0m^2\frac{1}{\beta }Dz\mathrm{log}\mathrm{cosh}\left[\beta J_0m+\beta Jz\sqrt{q}\right]$$ (89) $$q=Dz\mathrm{tanh}^2\left[\beta J_0m+\beta Jz\sqrt{q}\right]m=Dz\mathrm{tanh}\left[\beta J_0m+\beta Jz\sqrt{q}\right]$$ (90) Writing the equation for $`m`$ in integral form gives $$m=\beta J_0m_0^1𝑑\lambda \left[1Dz\mathrm{tanh}^2\left[\lambda \beta J_0m+\beta Jz\sqrt{q}\right]\right]$$ From this expression, in combination with (90), we conclude: $$T>J_0:m=0T>J_0\mathrm{and}T>J:m=q=0$$ Linearisation of (90) for small $`q`$ and $`m`$ shows the following continuous bifurcations: $$\begin{array}{cccc}& \mathrm{at}\hfill & \mathrm{from}\hfill & \mathrm{to}\text{}\hfill \\ J_0>J:\hfill & T=J_0\hfill & m=q=0\hfill & m0,q>0\hfill \\ J_0<J:\hfill & T=J\hfill & m=q=0\hfill & m=0,q>0\hfill \\ T<\mathrm{max}\{J_0,J\}:\hfill & T=J_0[1q]\hfill & m=0,q>0\hfill & m0,q>0\hfill \end{array}$$ Solving numerically equations $`T=J_0[1q]`$ and (90) leads to the phase diagram shown in figure 10. Breaking of Replica Symmetry: the AT Instability. If for the replica symmetric solution we calculate the entropy $`S=\beta ^2F/\beta `$ numerically, we find that for small temperatures it becomes negative. This is not possible. Firstly, straightforward differentiation shows $`S/\beta =\beta [H_{\mathrm{eq}}^2H^2_{\mathrm{eq}}]0`$, so $`S`$ increases with the noise level $`T`$. Let us now write $`H(𝝈)=H_0+\widehat{H}(𝝈)`$, where $`H_0`$ is the ground state energy and $`\widehat{H}(𝝈)0`$ (zero only for ground state configurations, the number of which we denote by $`N_01`$). We now find $$\underset{T0}{lim}S=\underset{\beta \mathrm{}}{lim}\left\{\mathrm{log}\underset{𝝈}{}e^{\beta H(𝝈)}+\beta H_{\mathrm{eq}}\right\}=\underset{\beta \mathrm{}}{lim}[\mathrm{log}\underset{𝝈}{}e^{\beta \widehat{H}(𝝈)}+\beta \widehat{H}_{\mathrm{eq}}]\mathrm{log}N_0$$ We conclude that $`S0`$ for all $`T`$. At small temperatures the RS ansatz (88) is apparently incorrect in that it no longer corresponds to the minimum of $`f(𝒒,𝒎)`$ (81). If saddle-points without replica symmetry bifurcate continuously from the RS one, we can locate the occurrence of this ‘replica symmetry breaking’ (RSB) by studying the effect on $`f(𝒒,𝒎)`$ of small fluctuations around the RS solution. It was shown that the ‘dangerous’ fluctuations are of the form $$q_{\alpha \beta }\delta _{\alpha \beta }+q\left[1\delta _{\alpha \beta }\right]+\eta _{\alpha \beta },\underset{\beta }{}\eta _{\alpha \beta }=0\alpha $$ (91) in which $`q`$ is the solution of (90) and $`\eta _{\alpha \beta }=\eta _{\beta \alpha }`$. We now calculate the resulting change in $`f(𝒒,𝒎)`$, away from the RS value $`f(𝒒_{\mathrm{RS}},𝒎_{\mathrm{RS}})`$, the leading order of which is quadratic in the fluctuations $`\{\eta _{\alpha \beta }\}`$ since the RS solution of (90) is a saddle-point: $$f(𝒒,𝒎)f(𝒒_{\mathrm{RS}},𝒎_{\mathrm{RS}})=\frac{\beta J^2}{4n}\underset{\alpha \gamma }{}\eta _{\alpha \gamma }^2\frac{\beta ^3J^4}{8n}\underset{\alpha \gamma }{}\underset{\rho \lambda }{}\eta _{\alpha \gamma }\eta _{\rho \lambda }G_{\alpha \gamma \rho \lambda }$$ with $$G_{\alpha \gamma \rho \lambda }=\frac{\sigma _\alpha \sigma _\gamma \sigma _\rho \sigma _\lambda e^{\frac{1}{2}q\beta ^2J^2\left[_\alpha \sigma _\alpha \right]^2+\beta mJ_0_\alpha \sigma _\alpha }_𝝈}{e^{\frac{1}{2}q\beta ^2J^2\left[_\alpha \sigma _\alpha \right]^2+\beta mJ_0_\alpha \sigma _\alpha }_𝝈}$$ Because of the index permutation symmetry in the above average we can write for $`\alpha \gamma `$ and $`\rho \lambda `$: $$G_{\alpha \gamma \rho \lambda }=\delta _{\alpha \rho }\delta _{\gamma \lambda }+\delta _{\alpha \lambda }\delta _{\gamma \rho }+G_4\left[1\delta _{\alpha \rho }\right]\left[1\delta _{\gamma \lambda }\right]\left[1\delta _{\alpha \lambda }\right]\left[1\delta _{\gamma \rho }\right]\text{}$$ $$+G_2\left\{\delta _{\alpha \rho }\left[1\delta _{\gamma \lambda }\right]+\delta _{\gamma \lambda }\left[1\delta _{\alpha \rho }\right]+\delta _{\alpha \lambda }\left[1\delta _{\gamma \rho }\right]+\delta _{\gamma \rho }\left[1\delta _{\alpha \lambda }\right]\right\}$$ with $$G_{\mathrm{}}=\frac{Dz\mathrm{tanh}^{\mathrm{}}\left[\beta J_0m+\beta Jz\sqrt{q}\right]\mathrm{cosh}^n\left[\beta J_0m+\beta Jz\sqrt{q}\right]}{Dz\mathrm{cosh}^n\left[\beta J_0m+\beta Jz\sqrt{q}\right]}$$ Only terms which involve precisely two $`\delta `$-functions can contribute, because of the requirements $`\alpha \gamma `$, $`\rho \lambda `$ and $`_\beta \eta _{\alpha \beta }=0`$. As a result: $$f(𝒒,𝒎)f(𝒒_{\mathrm{RS}},𝒎_{\mathrm{RS}})=\frac{\beta J^2}{4n}\left[1\beta ^2J^2\left(12G_2+G_4\right)\right]\underset{\alpha \gamma }{}\eta _{\alpha \gamma }^2$$ The condition for the RS solution to minimise $`f(𝒒,𝒎)`$, if compared to the so called ‘replicon’ fluctuations (91), is therefore $$1>\beta ^2J^2\underset{n0}{lim}\left(12G_2+G_4\right)\text{}$$ After taking the limit in the expressions $`G_{\mathrm{}}`$ this condition can be written as $$1>\beta ^2J^2Dz\mathrm{cosh}^4\left[\beta J_0m+\beta Jz\sqrt{q}\right]$$ (92) The so-called AT line in the phase diagram where this condition ceases to be met, indicates a continuous transition to a complex ‘spin-glass’ state where ergodicity is broken (i.e. the distribution $`\overline{P(q)}`$ (86) is no longer a $`\delta `$-function). It is shown in figure 10 as a dashed line for $`J_0/J>1`$, and coincides with the line $`T/J=1`$ for $`J_0<1`$. ## 6 The Hopfield Model Near Saturation ### 6.1 Replica Analysis We now turn to the Hopfield model with an extensive number of stored patterns, i.e. $`p=\alpha N`$ in (40). We can still write the free energy in the form (48), but this will not be of help since here it involves integrals over an extensive number of variables, so that steepest descent integration does not apply. Instead, following the approach of the previous model (72), we assume that we can average the free energy over the distribution of the patterns, with help of the replica-trick (75): $$\overline{F}=\underset{n0}{lim}\frac{1}{\beta n}\mathrm{log}\underset{𝝈^1\mathrm{}𝝈^n}{}\overline{e^{\beta _{\alpha =1}^nH(𝝈^\alpha )}}$$ Greek indices will denote either replica labels or pattern labels (it will be clear from the context), i.e. $`\alpha ,\beta =1,\mathrm{},n`$ and $`\mu ,\nu =1,\mathrm{},p`$. The $`p\times N`$ pattern components $`\{\xi _i^\mu \}`$ are assumed to be drawn independently at random from $`\{1,1\}`$. Replica Calculation of the Disorder-Averaged Free Energy. We first add to the Hamiltonian of (30) a finite number $`\mathrm{}`$ of generating terms, that will allow us to obtain expectation values of the overlap order parameters $`m_\mu `$ (41) by differentiation of the free energy (since all patterns are equivalent in the calculation we may choose these $`\mathrm{}`$ nominated patterns arbitrarily): $$HH+\underset{\mu =1}{\overset{\mathrm{}}{}}\lambda _\mu \underset{i}{}\sigma _i\xi _i^\mu m_\mu (𝝈)_{\mathrm{eq}}=\underset{𝝀\mathrm{𝟎}}{lim}\frac{}{\lambda _\mu }F/N$$ (93) We know how to deal with a finite number of overlaps and corresponding patterns, therefore we average only over the disorder that is responsible for the complications: the patterns $`\{𝝃^{\mathrm{}+1},\mathrm{},𝝃^p\}`$ (as in the previous section we denote this disorder-averaging by $`\overline{\mathrm{}}`$). Upon inserting the extended Hamiltonian into the replica-expression for the free energy, and assuming that the order of the limits $`N\mathrm{}`$ and $`n0`$ can be interchanged, we obtain for large $`N`$: $$\overline{F}/N=\frac{1}{2}\alpha \frac{1}{\beta }\mathrm{log}2\underset{n0}{lim}\frac{1}{\beta Nn}\mathrm{log}e^{\beta _\mu \mathrm{}_\alpha \left[\lambda _\mu _i\sigma _i^\alpha \xi _i^\mu \frac{1}{2N}\left[_i\sigma _i^\alpha \xi _i^\mu \right]^2\right]}\overline{e^{\frac{\beta }{2N}_\alpha _{\mu >\mathrm{}}\left[_i\sigma _i^\alpha \xi _i^\mu \right]^2}}_{\{𝝈^\alpha \}}$$ We linearise the $`\mu \mathrm{}`$ quadratic term using the identity (56), leading to $`n\times \mathrm{}`$ Gaussian integrals with $`D𝒎=(Dm_1^1,\mathrm{},Dm_n^{\mathrm{}})`$: $$\overline{F}/N=\frac{1}{2}\alpha \frac{1}{\beta }\mathrm{log}2\underset{n0}{lim}\frac{1}{\beta Nn}\mathrm{log}D𝒎e^{_\mu \mathrm{}_\alpha _i\sigma _i^\alpha \xi _i^\mu \left[\sqrt{\frac{\beta }{N}}m_\alpha ^\mu \beta \lambda _\mu \right]}\overline{e^{\frac{\beta }{2N}_\alpha _{\mu >\mathrm{}}\left[_i\sigma _i^\alpha \xi _i^\mu \right]^2}}_{\{𝝈^\alpha \}}$$ Anticipating that only terms exponential in the system size $`N`$ will retain statistical relevance in the limit $`N\mathrm{}`$, we rescale the $`n\times \mathrm{}`$ integration variables $`𝒎`$ according to $`𝒎𝒎\sqrt{\beta N}`$: $$\overline{F}/N=\frac{1}{2}\alpha \frac{1}{\beta }\mathrm{log}2\underset{n0}{lim}\frac{1}{\beta Nn}\mathrm{log}\{\left[\frac{\beta N}{2\pi }\right]^{\frac{n\mathrm{}}{2}}d𝒎e^{\frac{1}{2}\beta N𝒎^2}\times $$ $$e^{\beta _\mu \mathrm{}_\alpha _i\sigma _i^\alpha \xi _i^\mu \left[m_\alpha ^\mu \lambda _\mu \right]}\overline{e^{\frac{\beta }{2N}_\alpha _{\mu >\mathrm{}}\left[_i\sigma _i^\alpha \xi _i^\mu \right]^2}}_{\{𝝈^\alpha \}}\}$$ (94) Next we turn to the disorder average, where we again linearise the exponent containing the pattern components using the identity (56), with $`D𝒛=(Dz_1,\mathrm{},Dz_n)`$: $$\overline{e^{\frac{\beta }{2N}_\alpha _{\mu >\mathrm{}}\left[_i\sigma _i^\alpha \xi _i^\mu \right]^2}}=\left\{\overline{e^{\frac{1}{2}_\alpha \left[\left(\frac{\beta }{N}\right)^{\frac{1}{2}}_i\sigma _i^\alpha \xi _i\right]^2}}\right\}^p\mathrm{}=\left\{D𝒛\overline{e^{\left(\frac{\beta }{N}\right)^{\frac{1}{2}}_\alpha z_\alpha _i\sigma _i^\alpha \xi _i}}\right\}^p\mathrm{}$$ $$=\left\{D𝒛\underset{i}{}\mathrm{cosh}\left[\left(\frac{\beta }{N}\right)^{\frac{1}{2}}\underset{\alpha }{}z_\alpha \sigma _i^\alpha \right]\right\}^p\mathrm{}=\left\{D𝒛e^{\frac{\beta }{2N}_{\alpha \beta }z_\alpha z_\beta _i\sigma _i^\alpha \sigma _i^\beta +𝒪(\frac{1}{N})}\right\}^p$$ (95) We are now as in the previous case led to introducing the replica order parameters $`q_{\alpha \beta }`$: $$1=𝑑𝒒\underset{\alpha \beta }{}\delta \left[q_{\alpha \beta }\frac{1}{N}\underset{i}{}\sigma _i^\alpha \sigma _i^\beta \right]\text{}=\left[\frac{N}{2\pi }\right]^{n^2}𝑑𝒒𝑑\widehat{𝒒}e^{iN_{\alpha \beta }\widehat{q}_{\alpha \beta }\left[q_{\alpha \beta }\frac{1}{N}_i\sigma _i^\alpha \sigma _i^\beta \right]}\text{}$$ Inserting (95) and the above identities into (94) and assuming that the limits $`N\mathrm{}`$ and $`n0`$ commute gives: $$\underset{N\mathrm{}}{lim}\overline{F}/N=\frac{1}{2}\alpha \frac{1}{\beta }\mathrm{log}2\underset{N\mathrm{}}{lim}\underset{n0}{lim}\frac{1}{\beta Nn}\mathrm{log}𝑑𝒎𝑑𝒒𝑑\widehat{𝒒}e^{N\left[i_{\alpha \beta }\widehat{q}_{\alpha \beta }q_{\alpha \beta }\frac{1}{2}\beta 𝒎^2+\alpha \mathrm{log}{\scriptscriptstyle D𝒛e^{{\scriptscriptstyle \frac{\beta }{2}}_{\alpha \beta }z_\alpha z_\beta q_{\alpha \beta }}}\right]}$$ $$\times e^{\beta _\mu \mathrm{}_\alpha _i\sigma _i^\alpha \xi _i^\mu \left[m_\alpha ^\mu \lambda _\mu \right]i_{\alpha \beta }\widehat{q}_{\alpha \beta }_i\sigma _i^\alpha \sigma _i^\beta }_{\{𝝈^\alpha \}}$$ The $`n`$-dimensional Gaussian integral over $`𝒛`$ factorises in the standard way after appropriate rotation of the integration variables $`𝒛`$, with the result: $$\mathrm{log}D𝒛e^{\frac{\beta }{2}_{\alpha \beta }z_\alpha z_\beta q_{\alpha \beta }}=\frac{1}{2}\mathrm{log}det\left[1\mathrm{I}\beta 𝒒\right]$$ in which $`1\mathrm{I}`$ denotes the $`n\times n`$ identity matrix. The neuron averages factorise and are reduced to single-site ones over the $`n`$-replicated neuron $`𝝈=(\sigma _1,\mathrm{},\sigma _n)`$: $$\underset{N\mathrm{}}{lim}\overline{F}/N=\frac{1}{2}\alpha \frac{1}{\beta }\mathrm{log}2\underset{N\mathrm{}}{lim}\underset{n0}{lim}\frac{1}{\beta Nn}\mathrm{log}𝑑𝒎𝑑𝒒𝑑\widehat{𝒒}e^{N\left[i_{\alpha \beta }\widehat{q}_{\alpha \beta }q_{\alpha \beta }\frac{1}{2}\beta 𝒎^2\frac{1}{2}\alpha \mathrm{log}det\left[1\mathrm{I}\beta 𝒒\right]\right]}$$ $$\times \underset{i}{}e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _i^\mu \left[m_\alpha ^\mu \lambda _\mu \right]i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈$$ and we arrive at integrals that can be evaluated by steepest descent, following the manipulations (76). If we denote averages over the remaining $`\mathrm{}`$ patterns in the familiar way $$𝝃=(\xi _1,\mathrm{},\xi _{\mathrm{}})\text{}\mathrm{\Phi }(𝝃)_𝝃=2^{\mathrm{}}\underset{𝝃\{1,1\}^{\mathrm{}}}{}\mathrm{\Phi }(𝝃)$$ we can write the final result in the form $$\underset{N\mathrm{}}{lim}\overline{F}/N=\underset{n0}{lim}\mathrm{extr}f(𝒎,𝒒,\widehat{𝒒})\text{}$$ (96) $$f(𝒎,𝒒,\widehat{𝒒})=\frac{1}{2}\alpha \frac{1}{\beta }\mathrm{log}2\frac{1}{\beta n}[\mathrm{log}e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _\mu \left[m_\alpha ^\mu \lambda _\mu \right]i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈_𝝃$$ $$+i\underset{\alpha \beta }{}\widehat{q}_{\alpha \beta }q_{\alpha \beta }\frac{1}{2}\beta 𝒎^2\frac{1}{2}\alpha \mathrm{log}det[1\mathrm{I}\beta 𝒒]]$$ Having arrived at a saddle-point problem we now first identify the expectation values of the overlaps with (93) (note: extremisation with respect to the saddle-point variables and differentiation with respect to $`𝝀`$ commute): $$\overline{m_\mu (𝝈)_{\mathrm{eq}}}=\underset{n0}{lim}\underset{𝝀\mathrm{𝟎}}{lim}\frac{}{\lambda _\mu }\mathrm{extr}f(𝒎,𝒒,\widehat{𝒒})$$ $$=\underset{n0}{lim}\xi _\mu \frac{\frac{1}{n}_\alpha \sigma _\alpha e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _\mu m_\alpha ^\mu i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈}{e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _\mu m_\alpha ^\mu i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈}_𝝃$$ (97) which is to be evaluated in the $`𝝀=\mathrm{𝟎}`$ saddle-point. Having served their purpose, the generating fields $`\lambda _\mu `$ can be set to zero and we can restrict ourselves to the $`𝝀=\mathrm{𝟎}`$ saddle-point problem: $$f(𝒎,𝒒,\widehat{𝒒})=\frac{1}{2}\alpha \frac{1}{\beta }\mathrm{log}2\frac{1}{\beta n}[i\underset{\alpha \beta }{}\widehat{q}_{\alpha \beta }q_{\alpha \beta }\frac{1}{2}\beta 𝒎^2\frac{1}{2}\alpha \mathrm{log}det[1\mathrm{I}\beta 𝒒]$$ $$+\mathrm{log}e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _\mu m_\alpha ^\mu i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈_𝝃]$$ (98) Variation of the parameters $`\{m_\alpha ^\mu ,\widehat{q}_{\alpha \beta },q_{\alpha \beta }\}`$ gives the saddle-point equations: $$m_\alpha ^\mu =\xi _\mu \frac{\sigma _\alpha e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _\mu m_\alpha ^\mu i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈}{e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _\mu m_\alpha ^\mu i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈}_𝝃$$ (99) $$q_{\lambda \rho }=\frac{\sigma _\lambda \sigma _\rho e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _\mu m_\alpha ^\mu i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈}{e^{\beta _\mu \mathrm{}_\alpha \sigma _\alpha \xi _\mu m_\alpha ^\mu i_{\alpha \beta }\widehat{q}_{\alpha \beta }\sigma _\alpha \sigma _\beta }_𝝈}_𝝃$$ (100) $$\widehat{q}_{\lambda \rho }=\frac{1}{2}i\alpha \beta \frac{𝑑𝒛z_\lambda z_\rho e^{\frac{1}{2}𝒛\left[1\mathrm{I}\beta 𝒒\right]𝒛}}{𝑑𝒛e^{\frac{1}{2}𝒛\left[1\mathrm{I}\beta 𝒒\right]𝒛}}$$ (101) furthermore, $$\overline{m_\mu (𝝈)_{\mathrm{eq}}}=\underset{n0}{lim}\frac{1}{n}\underset{\alpha }{}m_\alpha ^\mu $$ (102) replaces the identification (97). As expected, one always has $`q_{\alpha \alpha }=1`$. The diagonal elements $`\widehat{q}_{\alpha \alpha }`$ drop out of (99,100), their values are simply given as functions of the remaining parameters by (101). Physical Interpretation of Saddle Points. We proceed along the lines of the Gaussian model (72). If we apply the alternative version (84) of the replica trick to the Hopfield model, we can write the distribution of the $`\mathrm{}`$ overlaps $`𝒎=(m_1,\mathrm{},m_{\mathrm{}})`$ in equilibrium as $$P(𝒎)=\underset{n0}{lim}\frac{1}{n}\underset{\gamma }{}\underset{𝝈^1\mathrm{}𝝈^n}{}\delta [𝒎\frac{1}{N}\underset{i}{}\sigma _i^\gamma 𝝃_i]\underset{\alpha }{}e^{\beta H(𝝈^\alpha )}$$ with $`𝝃_i=(\xi _i^1,\mathrm{},\xi _i^{\mathrm{}})`$. Averaging this distribution over the disorder leads to expressions identical to those encountered in evaluating the disorder averaged free energy. By inserting the same delta-functions we arrive at the saddle-point integration (96,98) and find $$\overline{P(𝒎)}=\underset{n0}{lim}\frac{1}{n}\underset{\gamma }{}\delta \left[𝒎𝒎_\gamma \right]$$ (103) where $`𝒎_\gamma =(m_\gamma ^1,\mathrm{},m_\gamma ^{\mathrm{}})`$ refers to the relevant solution of (99,100,101). Similarly we imagine two systems $`𝝈`$ and $`𝝈^{}`$ with identical realisation of the interactions $`\{J_{ij}\}`$, both in thermal equilibrium, and use (84) to rewrite the distribution $`P(q)`$ for the mutual overlap between the microstates of the two systems $$P(q)=\underset{n0}{lim}\frac{1}{n(n1)}\underset{\lambda \gamma }{}\underset{𝝈^1\mathrm{}𝝈^n}{}\delta [q\frac{1}{N}\underset{i}{}\sigma _i^\lambda \sigma _i^\gamma ]\underset{\alpha }{}e^{\beta H(𝝈^\alpha )}$$ Averaging over the disorder again leads to the steepest descend integration (96,98) and we find $$\overline{P(q)}=\underset{n0}{lim}\frac{1}{n(n1)}\underset{\lambda \gamma }{}\delta \left[qq_{\lambda \gamma }\right]$$ (104) where $`\{q_{\lambda \gamma }\}`$ refers to the relevant solution of (99,100,101). Finally we analyse the physical meaning of the conjugate parameters $`\{\widehat{q}_{\alpha \beta }\}`$ for $`\alpha \beta `$. We will do this in more detail, the analysis being rather specific for the Hopfield model and slightly different from the derivations above. Again we imagine two systems $`𝝈`$ and $`𝝈^{}`$ with identical interactions $`\{J_{ij}\}`$, both in thermal equilibrium. We now use (84) to evaluate the covariance of the overlaps corresponding to non-nominated patterns: $$r=\frac{1}{\alpha }\underset{\mu =\mathrm{}+1}{\overset{p}{}}\overline{\frac{1}{N}\underset{i}{}\sigma _i\xi _i^\mu _{\mathrm{eq}}\frac{1}{N}\underset{i}{}\sigma _i^{}\xi _i^\mu _{\mathrm{eq}}}$$ (105) $$=\underset{n0}{lim}\frac{N\mathrm{}/\alpha }{n(n1)}\underset{\lambda \gamma }{}\underset{𝝈^1\mathrm{}𝝈^n}{}\overline{\left[\frac{1}{N}\underset{i}{}\sigma _i^\lambda \xi _i^p\right]\left[\frac{1}{N}\underset{i}{}\sigma _i^\gamma \xi _i^p\right]\underset{\alpha }{}e^{\beta H(𝝈^\alpha )}}$$ (using the equivalence of all such patterns). We next perform the same manipulations as in calculating the free energy. Here the disorder average involves $$\overline{\left[\frac{1}{\sqrt{N}}\underset{i}{}\sigma _i^\lambda \xi _i^p\right]\left[\frac{1}{\sqrt{N}}\underset{i}{}\sigma _i^\gamma \xi _i^p\right]e^{\frac{\beta }{2N}_\alpha _{\mu >\mathrm{}}\left[_i\sigma _i^\alpha \xi _i^\mu \right]^2}}$$ $$=\left\{D𝒛\overline{e^{\left(\frac{\beta }{N}\right)^{\frac{1}{2}}_\alpha z_\alpha _i\sigma _i^\alpha \xi _i}}\right\}^{p\mathrm{}1}\frac{D𝒛}{\beta }\frac{^2}{z_\lambda z_\gamma }\overline{e^{\left(\frac{\beta }{N}\right)^{\frac{1}{2}}_\alpha z_\alpha _i\sigma _i^\alpha \xi _i}}$$ $$=\left\{D𝒛\overline{e^{\left(\frac{\beta }{N}\right)^{\frac{1}{2}}_\alpha z_\alpha _i\sigma _i^\alpha \xi _i}}\right\}^{p\mathrm{}1}D𝒛\frac{z_\lambda z_\gamma }{\beta }\overline{e^{\left(\frac{\beta }{N}\right)^{\frac{1}{2}}_\alpha z_\alpha _i\sigma _i^\alpha \xi _i}}$$ (after partial integration). We finally obtain an expression which involves the surface (98): $$r=\frac{1}{\beta }\underset{n0}{lim}\frac{1}{n(n1)}\underset{\lambda \rho }{}\underset{N\mathrm{}}{lim}\frac{𝑑𝒎𝑑𝒒𝑑\widehat{𝒒}\left[\frac{{\scriptscriptstyle 𝑑𝒛z_\lambda z_\rho e^{{\scriptscriptstyle \frac{1}{2}}𝒛\left[1\mathrm{I}\beta 𝒒\right]𝒛}}}{{\scriptscriptstyle 𝑑𝒛e^{{\scriptscriptstyle \frac{1}{2}}𝒛\left[1\mathrm{I}\beta 𝒒\right]𝒛}}}\right]e^{\beta nNf(𝒎,𝒒,\widehat{𝒒})}}{𝑑𝒎𝑑𝒒𝑑\widehat{𝒒}e^{\beta nNf(𝒎,𝒒,\widehat{𝒒})}}$$ The normalisation of the above integral over $`\{𝒎,𝒒,\widehat{𝒒}\}`$ follows from using the replica procedure to rewrite unity. The integration being dominated by the minima of $`f`$, we can use the saddle-point equations (101) to arrive at $$\underset{n0}{lim}\frac{1}{n(n1)}\underset{\lambda \rho }{}\widehat{q}_{\lambda \rho }=\frac{1}{2}i\alpha \beta ^2r$$ (106) The result (105,106) provides a physical interpretation of the order parameters $`\{\widehat{q}_{\alpha \beta }\}`$. Ergodicity implies that the distributions $`\overline{P(q)}`$ and $`\overline{P(𝒎)}`$ are $`\delta `$-functions, this is equivalent to the relevant saddle-point being of the form: $$m_\gamma ^\mu =m_\mu q_{\gamma \rho }=\delta _{\gamma \rho }+q\left[1\delta _{\gamma \rho }\right]\widehat{q}_{\gamma \rho }=\frac{1}{2}i\alpha \beta ^2\left[R\delta _{\gamma \rho }+r\left[1\delta _{\gamma \rho }\right]\right]$$ (107) which is the ‘replica symmetry’ (RS) ansatz for the Hopfield model. The RS form for $`\{q_{\alpha \beta }\}`$ and $`\{m_\alpha ^\mu \}`$ is a direct consequence of the corresponding distributions being $`\delta `$-functions, whereas the RS form for $`\{\widehat{q}_{\alpha \beta }\}`$ subsequently follows from (101). The physical meaning of $`m_\mu `$ and $`q`$ is $$m_\mu =\overline{m_\mu (𝝈)_{\mathrm{eq}}}q=\frac{1}{N}\underset{i}{}\overline{\sigma _i_{\mathrm{eq}}^2}$$ Before proceeding with a full analysis of the RS saddle-point equations, we finally make a few tentative statements on the phase diagram. For $`\beta =0`$ we obtain the trivial result $`q_{\lambda \rho }=\delta _{\lambda \rho }`$, $`\widehat{q}_{\lambda \rho }=0`$, $`m_\alpha ^\mu =0`$. We can identify continuous bifurcations to a non-trivial state by expanding the saddle-point equations in first order in the relevant parameters: $$m_\alpha ^\mu =\beta m_\alpha ^\mu +\mathrm{},q_{\lambda \rho }=2i\widehat{q}_{\lambda \rho }+\mathrm{}(\lambda \rho ),\widehat{q}_{\lambda \rho }=\frac{1}{2}\frac{i\alpha \beta }{1\beta }\left[\delta _{\lambda \rho }+\frac{\beta }{1\beta }q_{\lambda \rho }\left[1\delta _{\lambda \rho }\right]\right]+\mathrm{}$$ Combining the equations for $`𝒒`$ and $`\widehat{𝒒}`$ gives $`q_{\lambda \rho }=\alpha \left[\frac{\beta }{1\beta }\right]^2q_{\lambda \rho }+\mathrm{}`$. Thus we expect a continuous transition at $`T=1+\sqrt{\alpha }`$ from the trivial state to an ordered state where $`q_{\lambda \rho }0`$, but still $`m_\mu _{\mathrm{eq}}=0`$ (a spin-glass state). ### 6.2 Replica Symmetric Solution and AT-Instability The symmetry of the ansatz (107) for the saddle-point allows us to diagonalise the matrix $`𝚲=1\mathrm{I}\beta 𝒒`$ which we encountered in the saddle-point problem, $`\mathrm{\Lambda }_{\alpha \beta }=[1\beta (1q)]\delta _{\alpha \beta }\beta q`$: $$\begin{array}{ccc}\mathrm{eigenspace}:\hfill & \mathrm{eigenvalue}:\hfill & \mathrm{multiplicity}:\\ 𝒙=(1,\mathrm{},1)\hfill & 1\beta (1q)\beta qn\hfill & 1\\ _\alpha x_\alpha =0\hfill & 1\beta (1q)\hfill & n1\end{array}$$ so that $$\mathrm{log}det𝚲=\mathrm{log}\left[1\beta (1q)\beta qn\right]+(n1)\mathrm{log}\left[1\beta (1q)\right]$$ $$=n\left[\mathrm{log}\left[1\beta (1q)\right]\frac{\beta q}{1\beta (1q)}\right]+𝒪(n^2)$$ Inserting the RS ansatz (107) for the saddle-point into (98), utilising the above expression for the determinant and the short-hand $`𝒎=(m_1,\mathrm{},m_{\mathrm{}})`$, gives $$f(𝒎_{\mathrm{RS}},𝒒_{\mathrm{RS}},\widehat{𝒒}_{\mathrm{RS}})=\frac{1}{\beta }\mathrm{log}2+\frac{1}{2}\alpha \left[1+\beta r(1q)\right]+\frac{1}{2}𝒎^2+\frac{\alpha }{2\beta }\left[\mathrm{log}\left[1\beta (1q)\right]\frac{\beta q}{1\beta (1q)}\right]$$ $$\frac{1}{\beta n}\mathrm{log}e^{\beta 𝒎𝝃_\alpha \sigma _\alpha +\frac{1}{2}\alpha r\beta ^2\left[_\alpha \sigma _\alpha \right]^2}_𝝈_𝝃+𝒪(n)$$ We now linearise the squares in the neuron averages with (56), subsequently average over the replicated neuron $`𝝈`$, use $`\mathrm{cosh}^n[x]=1+n\mathrm{log}\mathrm{cosh}[x]+𝒪(n^2)`$, and take the limit $`n0`$: $$\underset{N\mathrm{}}{lim}\overline{F}_{\mathrm{RS}}/N=\underset{n0}{lim}f(𝒎_{\mathrm{RS}},𝒒_{\mathrm{RS}},\widehat{𝒒}_{\mathrm{RS}})$$ $$=\frac{1}{2}𝒎^2+\frac{1}{2}\alpha \left[1+\beta r(1q)+\frac{1}{\beta }\mathrm{log}\left[1\beta (1q)\right]\frac{q}{1\beta (1q)}\right]\frac{1}{\beta }Dz\mathrm{log}2\mathrm{cosh}\beta \left[𝒎𝝃+z\sqrt{\alpha r}\right]_𝝃$$ (108) The saddle-point equations for $`𝒎`$, $`q`$ and $`r`$ can be obtained either by insertion of the RS ansatz (107) into (99,100,101) and subsequently taking the $`n0`$ limit, or by variation of the RS expression (108). The latter route is the fastest one. After performing partial integrations where appropriate we obtain the final result: $$𝒎=𝝃Dz\mathrm{tanh}\beta \left[𝒎𝝃+z\sqrt{\alpha r}\right]_𝝃$$ (109) $$q=Dz\mathrm{tanh}^2\beta \left[𝒎𝝃+z\sqrt{\alpha r}\right]_𝝃r=q\left[1\beta (1q)\right]^2$$ (110) By substitution of the equation for $`r`$ into the remaining equations this set can easily be further reduced, should the need arise. In case of multiple solutions of (109,110) the relevant saddle-point is the one that minimises (108). Clearly for $`\alpha =0`$ we recover our previous results (50,51). Analysis of RS Order Parameter Equations and Phase Diagram. We first establish an upper bound for the temperature $`T=1/\beta `$ for non-trivial solutions of the set (109,110) to exist, by writing (109) in integral form: $$m_\mu =\beta \xi _\mu \left(𝝃𝒎\right)_0^1𝑑\lambda Dz\left[1\mathrm{tanh}^2\beta \left(\lambda 𝝃𝒎+z\sqrt{\alpha r}\right)\right]_𝝃$$ from which we deduce $$0=𝒎^2\beta \left(𝝃𝒎\right)^2_0^1𝑑\lambda Dz\left[1\mathrm{tanh}^2\beta \left(\lambda 𝝃𝒎+z\sqrt{\alpha r}\right)\right]_𝝃𝒎^2\beta \left(𝝃𝒎\right)^2_𝝃=𝒎^2\left[1\beta \right]$$ Therefore $`𝒎=0`$ for $`T>1`$. If $`T>1`$ we obtain in turn from (110), using $`\mathrm{tanh}^2(x)x^2`$ and $`0q1`$: $`q=0`$ or $`q1+\sqrt{\alpha }T`$. We conclude that $`q=0`$ for $`T>1+\sqrt{\alpha }`$. Secondly, for the free energy (108) to be well defined we must require $`q>1T`$. Linearisation of (109,110) for small $`q`$ and $`𝒎`$ shows the continuous bifurcations: $$\begin{array}{ccccccc}& & \mathrm{at}& & \mathrm{from}& & \mathrm{to}\\ \alpha >0:\hfill & & T=1+\sqrt{\alpha }& & 𝒎=\mathrm{𝟎},q=0& & 𝒎=\mathrm{𝟎},q>0\\ \alpha =0:\hfill & & T=1& & 𝒎=\mathrm{𝟎},q=0& & 𝒎\mathrm{𝟎},q>0\end{array}$$ The upper bound $`T=1+\sqrt{\alpha }`$ turns out to be the critical noise level indicating (for $`\alpha >0`$) a continuous transition to a spin-glass state, where there is no significant alignment of the neurons in the direction of one particular pattern, but still a certain degree of local freezing. Since $`𝒎=\mathrm{𝟎}`$ for $`T>1`$ this spin-glass state persists at least down to $`T=1`$. The quantitative details of the spin-glass state are obtained by inserting $`𝒎=\mathrm{𝟎}`$ into (110) (since (109) is fulfilled automatically). The impact on the saddle-point equations (109,110) of having $`\alpha >0`$, a smoothening of the hyperbolic tangent by convolution with a Gaussian kernel, can be viewed as noise caused by interference between the attractors. The natural strategy for solving (109,110) is therefore to make an ansatz for the nominated overlaps $`𝒎`$ of the type (52) (the mixture states). Insertion of this ansatz into the saddle-point equations indeed leads to self-consistent solutions. One can solve numerically the remaining equations for the amplitudes of the mixture states and evaluate their stability by calculating the eigenvalues of the second derivative of $`f(𝒎,𝒒,\widehat{𝒒})`$, in the same way as for $`\alpha =0`$. The calculations are just more involved. It then turns out that even mixtures are again unstable for any $`T`$ and $`\alpha `$, whereas odd mixtures can become locally stable for sufficiently small $`T`$ and $`\alpha `$. Among the mixture states, the pure states, where the vector $`𝒎`$ has only one nonzero component, are the first to stabilise as the temperature is lowered. These pure states, together with the spin-glass state ($`𝒎=0,q>0)`$, we will study in more detail. Let us first calculate the second derivatives of (108) and evaluate them in the spin-glass saddle-point. One finds, after elimination of $`r`$ with (110): $$^2f/m_\mu m_\nu =\delta _{\mu \nu }\left[1\beta (1q)\right]^2f/m_\mu q=0$$ The $`(\mathrm{}+1)\times (\mathrm{}+1)`$ matrix of second derivatives with respect to variation of $`(𝒎,q)`$, evaluated in the spin-glass saddle-point, thereby acquires a diagonal form $$^2f=(\begin{array}{cccc}1\beta (1q)& & & \\ & \mathrm{}& & \\ & & 1\beta (1q)& \\ & & & ^2f/q^2\end{array})$$ and the eigenvalues can simply be read off. The $`\mathrm{}`$-fold degenerate eigenvalue $`1\beta (1q)`$ is always positive (otherwise (108) would not even exist), implying stability of the spin-glass state in the direction of the nominated patterns. The remaining eigenvalue measures the stability of the spin-glass state with respect to variation in the amplitude $`q`$. Below the critical noise level $`T=1+\sqrt{\alpha }`$ it turns out to be positive for the spin-glass solution of (110) with nonzero $`q`$. One important difference between the previously studied case $`\alpha =0`$ and the present case $`\alpha >0`$ is that there is now a $`𝒎=\mathrm{𝟎}`$ spin-glass solution which is stable for all $`T<1+\sqrt{\alpha }`$. In terms of information processing this implies that for $`\alpha >0`$ an initial state must have a certain non-zero overlap with a pattern to evoke a final state with $`𝒎\mathrm{𝟎}`$, in order to avoid ending up in the $`𝒎=\mathrm{𝟎}`$ spin-glass state. This is clearly consistent with the observations in figure 5. In contrast, for $`\alpha =0`$, the state with $`𝒎=\mathrm{𝟎}`$ is unstable, so any initial state will eventually lead to a final state with $`𝒎\mathrm{𝟎}`$. Inserting the pure state ansatz $`𝒎=m(1,0,\mathrm{},0)`$ into our RS equations gives $$m=Dz\mathrm{tanh}\left[\beta m+\frac{z\beta \sqrt{\alpha q}}{1\beta (1q)}\right]q=Dz\mathrm{tanh}^2\left[\beta m+\frac{z\beta \sqrt{\alpha q}}{1\beta (1q)}\right]$$ (111) $$f=\frac{1}{2}m^2+\frac{1}{2}\alpha \left[(1q)\frac{1+\beta (1q)(\beta 2)}{\left[1\beta (1q)\right]^2}+\frac{1}{\beta }\mathrm{log}\left[1\beta (1q)\right]\right]\frac{1}{\beta }Dz\mathrm{log}2\mathrm{cosh}\left[\beta m+\frac{z\beta \sqrt{\alpha q}}{1\beta (1q)}\right]$$ (112) If we solve the equations (111) numerically for different values of $`\alpha `$, and calculate the corresponding ‘free energies’ $`f`$ (112) for the pure states and the spin-glass state $`𝒎=0`$, we obtain figure 11. For $`\alpha >0`$ the nontrivial solution $`m`$ for the amplitude of the pure state appears discontinously as the temperature is lowered, defining a critical temperature $`T_M(\alpha )`$. Once the pure state appears, it turns out to be locally stable (within the RS ansatz). Its ‘free energy’ $`f`$, however, remains larger than the one corresponding to the spin-glass state, until the temperature is further reduced to below a second critical temperature $`T_c(\alpha )`$. For $`T<T_c(\alpha )`$ the pure states are therefore the equilibrium states in the thermodynamics sense. By drawing these critical lines in the $`(\alpha ,T)`$ plane, together with the line $`T_g(\alpha )=1+\sqrt{\alpha }`$ which signals the second order transition from the paramagnetic to the spin-glass state, we obtain the RS phase diagram of the Hopfield model, depicted in figure 12. Strictly speaking the line $`T_M`$ would appear meaningless in the thermodynamic picture, only the saddle-point that minimises $`f`$ being relevant. However, we have to keep in mind the physics behind the formalism. The occurrence of multiple locally stable saddle-points is the manifestation of ergodicity breaking in the limit $`N\mathrm{}`$. The thermodynamic analysis, based on ergodicity, therefore applies only within a single ergodic component. Each locally stable saddle-point is indeed relevant for appropriate initial conditions and time-scales. Zero Temperature, Storage Capacity. The storage capacity $`\alpha _c`$ of the Hopfield model is defined as the largest $`\alpha `$ for which locally stable pure states exist. If for the moment we neglect the low temperature re-entrance peculiarities in the phase diagram (12) to which we will come back later, the critical temperature $`T_M(\alpha )`$, where the pure states appear decreases monotonically with $`\alpha `$, and the storage capacity is reached for $`T=0`$. Before we can put $`T0`$ in $`(\text{111})`$, however, we will have to rewrite these equations in terms of quantities with well defined $`T0`$ limits, since $`q1`$. A suitable quantity is $`C=\beta (1q)`$, which obeys $`0C1`$ for the free energy (108) to exist. The saddle-point equations can now be written in the form $$m=Dz\mathrm{tanh}\left[\beta m+\frac{z\beta \sqrt{\alpha q}}{1C}\right]C=\frac{}{m}Dz\mathrm{tanh}\left[\beta m+\frac{z\beta \sqrt{\alpha q}}{1C}\right]$$ in which the limit $`T0`$ simply corresponds to $`\mathrm{tanh}(\beta x)\mathrm{sgn}(x)`$ and $`q1`$. After having taken the limit we perform the Gaussian integral: $$m=\mathrm{erf}\left[\frac{m(1C)}{\sqrt{2\alpha }}\right]C=(1C)\sqrt{\frac{2}{\alpha \pi }}e^{m^2(1C)^2/2\alpha }$$ This set can be reduced to a single transcendental equation by introducing $`x=m(1C)/\sqrt{2\alpha }`$: $$x\sqrt{2\alpha }=F(x)F(x)=\mathrm{erf}(x)\frac{2x}{\sqrt{\pi }}e^{x^2}$$ (113) Equation (113) is solved numerically (see figure 13). Since $`F(x)`$ is anti-symmetric, solutions come in pairs $`(x,x)`$ (reflecting the symmetry of the Hamiltonian of the system with respect to an overall state-flip $`𝝈𝝈`$). For $`\alpha <\alpha _c0.138`$ there indeed exist pure state solutions $`x0`$. For $`\alpha >\alpha _c`$ there is only the spin-glass solution $`x=0`$. Given a solution $`x`$ of (113), the zero temperature values for the order parameters follow from $$\underset{T0}{lim}m=\mathrm{erf}[x]\underset{T0}{lim}C=\left[1+\sqrt{\frac{\alpha \pi }{2}}e^{x^2}\right]^1$$ with which in turn we can take the zero temperature limit in our expression (112) for the free energy: $$\underset{T0}{lim}f=\frac{1}{2}\mathrm{erf}^2[x]+\frac{1}{\pi }e^{x^2}\frac{2}{\pi }\left[e^{x^2}+\sqrt{\frac{\alpha \pi }{2}}\right]\left[x\sqrt{\pi }\mathrm{erf}[x]+e^{x^2}\text{}\right]$$ Comparison of the values for $`lim_{T0}f`$ thus obtained, for the pure state $`m>0`$ and the spin-glass state $`m=0`$ leads to figure 13, which clearly shows that for sufficiently small $`\alpha `$ the pure states are the true ground states of the system. The AT-Instability. As in the case of the Gaussian model (72), the above RS solution again generates negative entropies at sufficiently low temperatures, indicating that replica-symmetry must be broken. We can locate continuous replica symmetry breaking by studying the effect on $`f(𝒎,𝒒,\widehat{𝒒})`$ (98) of small replicon fluctuations around the RS solution: $$q_{\alpha \beta }\delta _{\alpha \beta }+q\left[1\delta _{\alpha \beta }\right]+\eta _{\alpha \beta },\eta _{\alpha \beta }=\eta _{\beta \alpha }\eta _{\alpha \alpha }=0\underset{\alpha }{}\eta _{\alpha \beta }=0$$ (114) The variation of $`𝒒`$ induces a similar variation in the conjugate parameters $`\widehat{𝒒}`$ through equation (101): $$\widehat{q}_{\alpha \beta }\frac{1}{2}i\alpha \beta ^2\left[R\delta _{\alpha \beta }+r\left[1\delta _{\alpha \beta }\right]+\widehat{\eta }_{\alpha \beta }\right]\widehat{\eta }_{\alpha \beta }=\frac{1}{2}\underset{\gamma \delta }{}\eta _{\gamma \delta }\left[g_{\alpha \beta \gamma \delta }g_{\alpha \beta }g_{\gamma \delta }\right]$$ with $$g_{\alpha \beta \gamma \delta }=\frac{𝑑𝒛z_\alpha z_\beta z_\gamma z_\delta e^{\frac{1}{2}𝒛\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]𝒛}}{𝑑𝒛e^{\frac{1}{2}𝒛\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]𝒛}}g_{\alpha \beta }=\frac{𝑑𝒛z_\alpha z_\beta e^{\frac{1}{2}𝒛\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]𝒛}}{𝑑𝒛e^{\frac{1}{2}𝒛\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]𝒛}}$$ Wick’s theorem (see e.g. ) can now be used to write everything in terms of second moments of the Gaussian integrals only: $$g_{\alpha \beta \gamma \delta }=g_{\alpha \beta }g_{\gamma \delta }+g_{\alpha \gamma }g_{\beta \delta }+g_{\alpha \delta }g_{\beta \gamma }$$ with which we can express the replicon variation in $`\widehat{𝒒}`$, using the symmetry of $`\{\eta _{\alpha \beta }\}`$ and the saddle-point equation (101), as $$\widehat{\eta }_{\alpha \beta }=\underset{\gamma \delta }{}g_{\alpha \gamma }\eta _{\gamma \delta }g_{\delta \beta }=\beta ^2\underset{\gamma \delta }{}\left[R\delta _{\alpha \gamma }+r\left[1\delta _{\alpha \gamma }\right]\right]\eta _{\gamma \delta }\left[R\delta _{\delta \beta }+r\left[1\delta _{\delta \beta }\right]\right]=\beta ^2(Rr)^2\eta _{\alpha \beta }$$ (115) since only those terms can contribute which involve precisely two $`\delta `$-symbols, due to $`_\alpha \eta _{\alpha \beta }=0`$. We can now calculate the change in $`f(𝒎,𝒒,\widehat{𝒒})`$, away from the RS value $`f(𝒎_{\mathrm{RS}},𝒒_{\mathrm{RS}},\widehat{𝒒}_{\mathrm{RS}})`$, the leading order of which must be quadratic in the fluctuations $`\{\eta _{\alpha \beta }\}`$ since the RS solution is a saddle-point: $$f(𝒎_{\mathrm{RS}},𝒒,\widehat{𝒒})f(𝒎_{\mathrm{RS}},𝒒_{\mathrm{RS}},\widehat{𝒒}_{\mathrm{RS}})=\frac{1}{\beta n}[\frac{1}{2}\alpha \mathrm{log}\frac{det\left[1\mathrm{I}\beta (𝒒_{\mathrm{RS}}+𝜼)\right]}{det\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]}i\mathrm{Tr}[\widehat{𝒒}_{\mathrm{RS}}.𝜼]$$ $$+\frac{1}{2}\alpha \beta ^2\mathrm{Tr}[\widehat{𝜼}.𝜼+\widehat{𝜼}.𝒒_{\mathrm{RS}}]\mathrm{log}\frac{e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈[\widehat{𝒒}_{\mathrm{RS}}+\frac{1}{2}i\alpha \beta ^2\widehat{𝜼}]𝝈}_𝝈}{e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}_𝝈}_𝝃]$$ (116) Evaluating (116) is simplified by the fact that the matrices $`𝒒_{\mathrm{RS}}`$ and $`𝜼`$ commute, which is a direct consequence of the properties (114) of the replicon fluctuations and the form of the replica-symmetric saddle-point. If we define the $`n\times n`$ matrix $`𝑷`$ as the projection onto the vector $`(1,\mathrm{},1)`$, we have $$P_{\alpha \beta }=n^1𝑷.𝜼=𝜼.𝑷=0𝒒_{\mathrm{RS}}=(1q)1\mathrm{I}+nq𝑷𝒒_{\mathrm{RS}}.𝜼=𝜼.𝒒_{\mathrm{RS}}=(1q)𝜼$$ (117) $$\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]^1=\frac{1}{1\beta (1q)}1\mathrm{I}+\frac{\beta nq}{\left[1\beta (1q)\beta nq\right]\left[1\beta (1q)\right]}𝑷$$ We can now simply expand the relevant terms, using the identity $`\mathrm{log}detM=\mathrm{Tr}\mathrm{log}M`$: $$\mathrm{log}\frac{det\left[1\mathrm{I}\beta (𝒒_{\mathrm{RS}}+𝜼)\right]}{det\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]}=\mathrm{Tr}\mathrm{log}\left[1\mathrm{I}\beta 𝜼\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]^1\right]$$ $$=\mathrm{Tr}\left\{\beta 𝜼\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]^1\frac{1}{2}\beta ^2\left[𝜼\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]^1\right]^2\right\}+𝒪(𝜼^3)$$ $$=\frac{1}{2}\frac{\beta ^2}{\left[1\beta (1q)\right]^2}\mathrm{Tr}𝜼^2+𝒪(𝜼^3)$$ (118) Finally we address the remaining term in (116), again using the RS saddle-point equations (109,110) where appropriate: $$\mathrm{log}\frac{e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}\left[1+\frac{1}{2}\alpha \beta ^2𝝈\widehat{𝜼}𝝈+\frac{1}{8}\alpha ^2\beta ^4(𝝈\widehat{𝜼}𝝈)^2+\mathrm{}\right]_𝝈}{e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}_𝝈}_𝝃$$ $$=\frac{1}{2}\alpha \beta ^2\mathrm{Tr}[\widehat{𝜼}.𝒒_{\mathrm{RS}}]+\frac{1}{8}\alpha ^2\beta ^4\underset{\alpha \beta \gamma \delta }{}\widehat{\eta }_{\alpha \beta }\widehat{\eta }_{\gamma \delta }[G_{\alpha \beta \gamma \delta }H_{\alpha \beta \gamma \delta }]+\mathrm{}$$ (119) with $$G_{\alpha \beta \gamma \delta }=\frac{\sigma _\alpha \sigma _\beta \sigma _\gamma \sigma _\delta e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}_𝝈}{e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}_𝝈}_𝝃$$ $$H_{\alpha \beta \gamma \delta }=\frac{\sigma _\alpha \sigma _\beta e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}_𝝈}{e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}_𝝈}\frac{\sigma _\gamma \sigma _\delta e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}_𝝈}{e^{\beta 𝝃𝒎_{\mathrm{RS}}_\alpha \sigma _\alpha i𝝈\widehat{𝒒}_{\mathrm{RS}}𝝈}_𝝈}_𝝃$$ Inserting the ingredients (115,117,118,119) into expression (116) and rearranging terms shows that the linear terms indeed cancel, and that the term involving $`H_{\alpha \beta \gamma \delta }`$ does not contribute (since the elements $`H_{\alpha \beta \gamma \delta }`$ don’t depend on the indices for $`\alpha \beta `$ and $`\gamma \delta `$), and we are left with: $$f(𝒎_{\mathrm{RS}},𝒒,\widehat{𝒒})f(𝒎_{\mathrm{RS}},𝒒_{\mathrm{RS}},\widehat{𝒒}_{\mathrm{RS}})=\frac{1}{\beta n}[\frac{1}{4}\frac{\alpha \beta ^2}{\left[1\beta (1q)\right]^2}\mathrm{Tr}𝜼^2+\frac{1}{2}\alpha \beta ^4(Rr)^2\mathrm{Tr}𝜼^2$$ $$\frac{1}{8}\alpha ^2\beta ^8(Rr)^4\underset{\alpha \beta \gamma \delta }{}\eta _{\alpha \beta }\eta _{\gamma \delta }G_{\alpha \beta \gamma \delta }]+\mathrm{}$$ Because of the index permutation symmetry in the neuron average we can write for $`\alpha \gamma `$ and $`\rho \lambda `$: $$G_{\alpha \gamma \rho \lambda }=\delta _{\alpha \rho }\delta _{\gamma \lambda }+\delta _{\alpha \lambda }\delta _{\gamma \rho }+G_4\left[1\delta _{\alpha \rho }\right]\left[1\delta _{\gamma \lambda }\right]\left[1\delta _{\alpha \lambda }\right]\left[1\delta _{\gamma \rho }\right]$$ $$+G_2\left\{\delta _{\alpha \rho }\left[1\delta _{\gamma \lambda }\right]+\delta _{\gamma \lambda }\left[1\delta _{\alpha \rho }\right]+\delta _{\alpha \lambda }\left[1\delta _{\gamma \rho }\right]+\delta _{\gamma \rho }\left[1\delta _{\alpha \lambda }\right]\right\}$$ with $$G_{\mathrm{}}=\frac{Dz\mathrm{tanh}^{\mathrm{}}\beta \left[𝒎𝝃+z\sqrt{\alpha r}\right]\mathrm{cosh}^n\beta \left[𝒎𝝃+z\sqrt{\alpha r}\right]}{Dz\mathrm{cosh}^n\beta \left[𝒎𝝃+z\sqrt{\alpha r}\right]}_𝝃$$ Only terms which involve precisely two $`\delta `$-functions can contribute, because of the replicon properties (114). As a result: $$f(𝒎_{\mathrm{RS}},𝒒,\widehat{𝒒})f(𝒎_{\mathrm{RS}},𝒒_{\mathrm{RS}},\widehat{𝒒}_{\mathrm{RS}})=\frac{1}{\beta n}\mathrm{Tr}𝜼^2[\frac{1}{4}\frac{\alpha \beta ^2}{\left[1\beta (1q)\right]^2}+\frac{1}{2}\alpha \beta ^4(Rr)^2$$ $$\frac{1}{4}\alpha ^2\beta ^8(Rr)^4[12G_2+G_4]]+\mathrm{}$$ Since $`\mathrm{Tr}𝜼^2=_{\alpha \beta }\eta _{\alpha \beta }^2`$, the condition for the RS solution to minimise $`f(𝒎,𝒒,\widehat{𝒒})`$, if compared to the ‘replicon’ fluctuations, is therefore $$\frac{1}{\left[1\beta (1q)\right]^2}+2\beta ^2(Rr)^2\alpha \beta ^6(Rr)^4\left[12G_2+G_4\right]>0$$ (120) After taking the limit in the expressions $`G_{\mathrm{}}`$ and after evaluating $$\underset{n0}{lim}R=\frac{1}{\beta }\underset{n0}{lim}g_{\alpha \alpha }=\underset{n0}{lim}\frac{1}{n\beta }\frac{𝑑𝒛𝒛^2e^{\frac{1}{2}𝒛\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]𝒛}}{𝑑𝒛e^{\frac{1}{2}𝒛\left[1\mathrm{I}\beta 𝒒_{\mathrm{RS}}\right]𝒛}}$$ $$=\underset{n0}{lim}\frac{1}{n\beta }\left[\frac{n1}{1\beta (1q)}+\frac{1}{1\beta (1q+nq)}\right]=\frac{1}{\beta }\frac{1\beta +2\beta q}{\left[1\beta (1q)\right]^2}$$ and using (110), the condition (120) can be written as $$[1\beta (1q)]^2>\alpha \beta ^2Dz\mathrm{cosh}^4\beta \left[𝒎𝝃+z\sqrt{\alpha r}\right]_𝝃$$ (121) The AT line in the phase diagram, where this condition ceases to be met, indicates a second-order transition to a spin-glass state where ergodicity is broken in the sense that the distribution $`\overline{P(q)}`$ (104) is no longer a $`\delta `$-function. In the paramagnetic regime of the phase diagram, $`𝒎=\mathrm{𝟎}`$ and $`q=0`$, the AT condition reduces precisely to $`T>T_g=1+\sqrt{\alpha }`$. Therefore the paramagnetic solution is stable. The AT line coincides with the boundary between the paramagnetic and spin-glass phase. Numerical evaluation of (121) shows that the RS spin-glass solution remains unstable for all $`T<T_g`$, but that the retrieval solution $`𝒎\mathrm{𝟎}`$ is unstable only for very low temperatures $`T<T_R`$ (see figure 12). ## 7 Epilogue In this paper I have tried to give a self-contained exposé of the main issues, models and mathematical techniques relating to the equilibrium statistical mechanical analysis of recurrent neural networks. I have included networks of binary neurons and networks of coupled (neural) oscillators, with various degrees of synaptic complexity (albeit always fully connected), ranging from uniform synapses, via synapses storing a small number of patterns, to Gaussian synapses and synapses encoding an extensive number of stored patterns. The latter (complex) cases I only worked out for binary neurons; similar calculations can be done for coupled oscillators (see ). Networks of graded response neurons could not be included, because these are found never to go to (detailed balance) equilibrium, ruling out equilibrium statistical mechanical analysis. All analytical results and predictions have later also been confirmed comprehensively by numerical simulations. Over the years we have learned an impressive amount about the operation of recurrent networks by thinking in terms of free energies and phase transitions, and by having been able to derive explicit analytical solutions (since a good theory always supersedes an infinite number of simulation experiments …). I have given a number of key references along the way; many could have been added but were left out for practical reasons. Instead I will just mention a number of textbooks in which more science as well as more references to research papers can be found. Any such selection is obviously highly subjective, and I wish to apologize beforehand to the authors which I regret to have omitted. Several relevant review papers dealing with the statistical mechanics of neural networks can be found scattered over the three volumes . Textbooks which attempt to take the interested but non-expert reader towards the expert level are . Finally, a good introduction to the methods and backgrounds of replica theory, together with a good collection of reprints of original papers, can be found in . What should we expect for the next decades, in the equilibrium statistical mechanics of recurrent neural networks ? Within the confined area of large symmetric and fully connected recurrent networks with simple neuron types we can now deal with fairly complicated choices for the synapses, inducing complicated energy landscapes with many stable states, but this involves non-trivial and cutting-edge mathematical techniques. If our basic driving force next is the aim to bring our models closer to biological reality, balancing the need to retain mathematical solvability with the desire to bring in more details of the various electro-chemical processes known to occur in neurons and synapses and spatio-temporal characteristics of dendrites, the boundaries of what can be done with equilibrium statistical mechanics are, roughly speaking, set by the three key issues of (presence or absence of) detailed balance, system size, and synaptic interaction range. The first issue is vital: no detailed balance immediately implies no equilibrium statistical mechanics. This generally rules out networks with non-symmetric synapses and all networks of graded response neurons (even when the latter are equipped with symmetric synapses). The issue of system size is slightly less severe; models of networks with $`N<\mathrm{}`$ neurons can often be solved in leading order in $`N^{\frac{1}{2}}`$, but a price will have to be paid in the form of a reduction of our ambition elsewhere (e.g. we might have to restrict ourselves to simpler choices of synaptic interactions). Finally, we know how to deal with fully connected models (such as those discussed in this paper), and also with models having dendritic structures which cover a long (but not infinite) range, provided they vary smoothly with distance. We can also deal with short-range dendrites in one-dimensional (and to a lesser extent two dimensional) networks; however, since even the relatively simple Ising model (mathematically equivalent to a network of binary neurons with uniform synapses connecting only nearest-neighbour neurons) has so far not yet been solved in three dimensions, it is not realistic to assume that analytical solution will be possible soon of general recurrent neural network models with short range interactions. On balance, although there are still many interesting puzzles to keep theorists happy for years to come, and although many of the model types discussed in this text will continue to be useful building blocks in explaining at a basic and qualitative level the operation of specific recurrent brain regions (such as the CA3 region of the hippocampus), one is therefore led to the conclusion that equilibrium statistical mechanics has by now brought us as far as can be expected with regard to increasing our understanding of biological neural networks. Dale’s law already rules out synaptic symmetry, and thereby equilibrium statistical mechanics altogether, so we are forced to turn to dynamical techniques if we wish to improve biological realism. ### Acknowledgements I is my great pleasure to thank David Sherrington and Nikos Skantzos for their contributions to content and presentation of this review.
warning/0006/hep-ph0006088.html	ar5iv	text	# Random values of the cosmological constant ## I Introduction The problem of understanding a small but non-zero cosmological constant ($`\mathrm{\Lambda }`$)<sup></sup><sup></sup>* If the cosmological constant is indeed the explanation of the recent supernova observations, the value would be $`\mathrm{\Lambda }=(1.2\pm 0.4)\times 10^{123}M_P^4`$, where $`M_P=1.22\times 10^{19}`$ GeV is the Planck mass. appears even harder than it would be if the cosmological constant were identically zero. There are many contributions to $`\mathrm{\Lambda }`$, ranging from zero-point energies to Higgs and QCD vacuum condensates. Observing a non-zero value tells us that we should not seek a principle that requires these contributions to cancel exactly. However, empirically a partial cancellation must occur and must be extremely fine-tuned in order to result in such a tiny residual. The problem is so severe that it forces us to take seriously the anthropic multiple domain solution which would naturally lead to the observation of a small non-zero $`\mathrm{\Lambda }`$. Under this hypothesis, the cosmological constant is a parameter that can take on different values in different domains of the universe, with an assumed cosmological evolution such that we live entirely within a single domain. Domains with “normal” values of the cosmological constant would collapse quickly or expand exponentially rapidly and could not lead to life of any form. Only those with a small enough residual $`\mathrm{\Lambda }`$ have the conditions appropriate for life and it is only this restricted range that we should consider. Under this hypothesis, we would expect to observe a non-zero value of $`\mathrm{\Lambda }`$, since there is no mechanism forcing it to be zero, and the magnitude would be expected to be typical of the anthropically allowed range. Weinberg has phrased this constraint in a physical way by asking about the mean value of $`\mathrm{\Lambda }`$ in universes in which matter clumps into galaxies, the clumping being a needed precursor to life. Under plausible estimates, the observed value of $`\mathrm{\Lambda }`$ is reasonably typical of the mean viable value. Rees and Tegmark have pointed out that in a more general context there is an allowed two dimensional area in the values of $`Q`$ and $`\mathrm{\Lambda }`$, where $`Q`$ is the magnitude of the initial density perturbations, again such that our values are reasonably typical. While this hypothesis could provide a natural explanation of the value of $`\mathrm{\Lambda }`$, its physical foundation remains unclear and we need to look for possible physical realizations. For the mechanism to be contained in a physical theory there must be two main ingredients. The first is the generation of an appropriately large universe with domains that are presently disconnected. This is a relatively simple requirement. There are many available ideas for having quantum fluctuations or random dynamics influence physics within one causally connected region of the early universe. Inflation, or pre-big-bang evolution can then insure that we live entirely within a region which evolved from a single such domain. Disconnected regions of the universe are a common occurrence in modern theories of cosmology. The difficult aspect of this hypothesis is contained in the second ingredient - the variability of physical parameters such as the cosmological constant. Ordinarily, coupling constants and masses are constant parameters uniquely defined within a theory. However in this hypothesis, the requirement is that these parameters can take on multiple values, yet are essentially constant throughout our domain. The values of the parameters are related to the ground state of the theory. Different ground states correspond to differences of at least some parameters of the low energy theory. The usual situation envisioned in fundamental theories is that there is a unique ground state to the theory, or at most a discrete few ground states. Even in string theory, which classically has continuous families of ground states, one normally expects that non-perturbative effects will select at most a few possible true ground states per compactification. However, it is unlikely that a set of discrete ground states is sufficient for implementing the anthropic selection (see the next section). If we turn to continuously variable states there are difficulties in maintaining a stable set of parameters. This paper examines issues associated with known mechanisms for implementing the multiple-domain/anthropic scenario. ## II Discrete versus continuous? The ground state of a theory like QCD appears to be unique. The electroweak theory has a continuous family of ground states, corresponding to different directions of the Higgs field, but they are all equivalent and all have the same parameters. In more complicated theories with multiple Higgs fields there can be several minima to the Higgs potential. These multiple minima are potentially applicable to the multiple domain problem. One possibility is that some symmetry leads to the condition where several ground states have the same energy. However, if they have the same ground state energy then they have the same cosmological constant, and are therefore not useful in this context. In the more common case where the minima are all of different energy, only one will be the true ground state. However, in some situations the time to tunnel from one ground state to the true vacua can be long enough that the metastable states can be considered in cosmology. These different ground states would correspond to different cosmological constants. Therefore it is reasonable to consider multiple metastable discrete ground states as a candidate mechanism in the multiple domain problem. However, a few discrete ground states are not enough for the anthropic solution for the cosmological constant. A theory with multiple ground states must occur at energies higher that that of the Standard Model. Let us denote the scale of this future theory by $`M_{}`$, with $`M_{}1`$ TeV. The ground states would be categorized by energies of this scale. In particular, the splitting between the ground state with the smallest negative cosmological constant and the smallest positive one would be of this size. It is extremely unlikely that a ground state would fall in the very tiny window that allows a anthropically viable cosmological constant. That window corresponds to a range $$\frac{(\mathrm{\Delta }\mathrm{\Lambda })_{\mathrm{anthropic}}}{(\mathrm{\Delta }\mathrm{\Lambda })_{\mathrm{natural}}}10^{58}\left[\frac{1TeV^4}{M_{}^4}\right]$$ (1) If a theory had a very densely packed set of states around $`\mathrm{\Lambda }=0`$, it would contain an unnaturally small parameter describing the spacing of these states (as well as possibly having difficulty arranging for these states to be metastable for long periods). The great disparity between $`M_{}^4`$ and $`\mathrm{\Lambda }`$ indicates that one would require an additional mechanism to generate the possibility of fine tuning an anthropically acceptable valueIn fact, new mechanisms that may allow closely spaced values for the cosmological constant have recently been addressed in . The alternative is that the parameters can vary continuously, yet stay frozen at an arbitrary fixed value throughout our domain. Here the requirement is only that both signs of the cosmological constant be possible. In this case, random dynamics will occasionally generate an acceptable $`\mathrm{\Lambda }`$ in the neighborhood of zero. However, this option is not without problems. Let us consider the basic difficulty in a general framework. If the parameters can have continuously different values, they would be different in causally disconnected domains in the early universe. Therefore they can be described by space-time dependent fields. This means that we will always be looking at the dynamics of some fields. Since by assumption these fields are not constrained to be near a unique value, their potential, if they have any, must be small and they would normally be described as nearly-massless fields. While inflation can readily lead to these fields becoming uniform throughout the observed universe, the difficult part is to understand why the dynamics of such a field did not lead it to evolve towards a unique ground state. Therefore we are lead to consider fields whose dynamics have been frozen at continuous values in some fashion. This is the topic of the rest of this paper. ## III Hubble damping There exists a simple mechanism that demonstrates that the freezing of dynamical fields at random values is indeed possible. It is related to the “slow roll” mechanism which is important for inflationHowever, it should be clearly stated that the discussion which follows does not apply to the field responsible for inflation, but to a different scalar field that is being invoked to address the cosmological constant problem.. Consider a scalar field in an expanding FRW universe governed by a scale factor $`a(t)`$. The equation of motion for this field is $$\ddot{\varphi }+3H\dot{\varphi }\frac{1}{a^2}^2\varphi =V^{}(\varphi ).$$ (2) where a dot denoted a derivative with respect to time and the prime denotes differentiation with respect to the field $`\varphi `$. $`V(\varphi )`$ is the potential for the scalar field and the Hubble parameter is defined by $$H=\frac{\dot{a(t)}}{a(t)}.$$ (3) For a field which is sufficiently spatially uniform, one can drop the term involving spatial gradients. In this case, a sufficiently flat potential, when compared to the Hubble constant, will lead to $`\dot{\varphi }0`$. Thus the Hubble expansion can damp the time evolution of a uniform scalar field with a sufficiently flat potential. The application of this mechanism to the question of the anthropic solution to the cosmological constant problem has been studied in detail in recent work by Garriga and Vilenken, and discussed by Weinberg. Let us look in more detail at the condition for the freezing of the field as we will see that there is a conflict related to the two uses of a flat potential in this scenario. On one hand the potential must be flat with respect the present Hubble parameter, which is a very small number, in order that the field be presently frozen. This corresponds to the intuitive expectation that the Hubble expansion plays very little role on the fields that we see around us, such that it must be a very weakly varying potential if the Hubble parameter is to provide the damping to keep the field from dynamically evolving. However, on the other hand, we need the potential to have enough variation that its contribution to the vacuum energy is sufficient to influence the cosmological constant. If the potential is too flat it contributes too weakly to the cosmological constant to be able to nearly cancel the other contributions to $`\mathrm{\Lambda }`$. These dual requirements force certain unnatural conditions on the potential and also pose an important requirement on the nature of inflation. The condition that the field remain effectively frozen today means, among other things, that it is not changing so fast as to contribute significantly to the present energy density. Thus its kinetic energy is bounded $$\frac{1}{2}\dot{\varphi }^2<<\rho _0$$ (4) where $`\rho _0`$ is the present energy density of the universe. This is related to the present Hubble constant (neglecting a possible curvature contribution to the Hubble constant) by $$H_0^2=\frac{8\pi }{3M_P^2}\rho _0$$ (5) Using the slow roll approximation, we find that this kinetic constraint implies $$V^{}(\varphi )3H_0\dot{\varphi }<<3H_0\sqrt{\frac{3M_P^2H_0^2}{8\pi }}H_0^2M_P10^{122}M_P^3.$$ (6) Here we have use the kinetic energy bound, dropped constants of order unity and used the value of the Hubble parameter $`H_010^{61}M_P`$. The conclusion is that the potential must be very flat. However, this means that reasonable variations of the magnitude of $`\varphi `$ do not change the vacuum energy much. Specifically, for this mechanism to be operable we need to be able to nearly cancel the effect of other sources of vacuum energy, which we can denote by $`\mathrm{\Lambda }_{\mathrm{other}}`$. The variation of the vacuum energy as we vary $`\varphi `$ must then be of order $`\mathrm{\Lambda }_{\mathrm{other}}`$. Let us distinguish two extreme situations: one where gravity or string theory provides the scale of the vacuum energy such that $`\mathrm{\Lambda }_{\mathrm{other}}M_P^4`$ and the other with low energy supersymmetry which implies $`\mathrm{\Lambda }_{\mathrm{other}}1`$ TeV<sup>4</sup>. In terms of the potential, the requirement is $$V^{}(\varphi )\mathrm{\Delta }\varphi \mathrm{\Lambda }_{\mathrm{other}}$$ (7) Thus the very small values of $`V^{}(\varphi )`$ can only be useful if $`\varphi \mathrm{\Delta }\varphi `$ is very large. Inserting the constraint on $`V^{}(\varphi )`$ from Hubble damping reveals just how large this value must be $$\varphi >>10^{122}M_P\left(\frac{\mathrm{\Lambda }_{\mathrm{other}}}{M_P^4}\right)$$ (8) Even if we have low energy supersymmetry this leads to a strikingly large value of $`\varphi >>10^{58}M_P`$. The extreme flatness of the potential is a potential difficulty. Allowing quadratic and quartic couplings, we will have $$V^{}(\varphi )=\mu ^2\varphi +\lambda \varphi ^3.$$ (9) Given the constraints on $`\varphi `$ and $`V^{}(\varphi )`$, we must have $`\mu ^2<10^{244}M_P^2`$ and $`\lambda <10^{488}`$ in the case where $`\mathrm{\Lambda }_{\mathrm{other}}`$ is determined by the Planck scale and $`\mu ^2<10^{180}M_P^2`$ and $`\lambda <10^{296}`$ in the most favorable case of weak scale supersymmetry breaking. At first sight this appears to be a fine tuning which is even greater than that of the cosmological constant. One might not be worried about the flatness of the potential since in supersymmetry flat potentials are ubiquitous, and one might hope that this flatness could be preserved. When supersymmetry is broken, radiative corrections will generate contributions to a potential. For a potential this flat, all matter fields certainly need to be decoupled from the $`\varphi `$ field, or else they would generate a large potential after supersymmetry breaking. The decoupling of matter fields is also required in order to not violate general relativity constraints. The field $`\varphi `$ is effectively massless, given the flatness of the potential, and would lead to observable long range forces if coupled to matter at even gravitational strength. However the constraints from the lack of radiative corrections to the potential are even stronger, and one is led to assume that matter fields can be completely decoupled from $`\varphi `$. This leads to the expectation that this field will not influence any of the other parameters of the Standard Model, as noted by Weinberg. It will be a formidable problem to generate a potential that is large enough to influence the cosmological constant, yet flat enough to not be presently evolving. It would be remarkable if the existence of a viable domain is only possible due to the existence of such a extreme potential. The needed initial conditions may also present a fine tuning problem. The size of the field $`\varphi `$ is not by itself is not the problem, since we have seen that despite this large value the energy associated with the field is still below the Planck mass. However for a field of this size to also have its kinetic energy below the Planck scale requires an unnatural spatial and temporal constancy. In this case the problem is not so much in the present epoch, when inflation could have smoothed out any spatial variation in $`\varphi `$, but in the early universe before the start of inflation. At this time, in order that the kinetic energies not exceed the Planck scale, we need the variation in time and in each of the spatial directions to satisfy $$_0\varphi _i\varphi \frac{\varphi }{L}<M_P^2$$ (10) with $`L`$ being the scale factor which describes the constancy of the field. (In an infinite universe $`L`$ would be the wavelength of the field configuration.) If the scalar field evolves classically its magnitude would be the same in the early universe, and one is then constrained to have $`L`$ $`>`$ $`10^{122}M_P^110^{80}\mathrm{light}\mathrm{years}(\mathrm{\Lambda }_{\mathrm{other}}M_P)`$ (11) $`>`$ $`10^{58}M_P^110^{16}\mathrm{light}\mathrm{years}(\mathrm{\Lambda }_{\mathrm{other}}1\mathrm{TeV}^4)`$ (12) Thus the requirement is that the field initially (at the start of inflation) have an extremely large value, but have a incredibly tiny spatial and temporal variation. These initial values are quite unnatural, and tell us that the classical evolutions is not a natural solution. However, quantum fluctuations during inflation can modify the field values, and if inflation is long enough, would remove the unnaturalness issue for the initial conditions. This then can be converted into a limit on the length of the inflationary epoch. Quantum fluctuations behave differently in an exponentially expanding space time. Long wavelength modes can get redshifted such that they become almost flat, in which case the Hubble damping freezes them to constant values that add to the value of the classical field. Since the quantum fluctuations carry either sign, this leads to a random walk character for the net field values. Different regions in the inflating domain can thus develop different values of the field, with a rms deviation that grows as $`\sqrt{t}`$. The heuristic explanation is as follows, although the result is derived from more rigorous calculations. In each causally connected region of size $`H^1`$, fluctuations are independent. (In this section $`H`$ refers to the Hubble constant during the period of inflation, rather than in the present epoch.) In an expansion time of $`H^1`$ a typical fluctuation is of size $`\mathrm{\Delta }\varphi H/2\pi `$. Since the expansion can freeze this field, over many expansion times these fluctuations then add as a random walk, resulting in a spread of values of order $$\delta \varphi ^2=\left(\frac{H}{2\pi }\right)^2Ht$$ (13) Because these fluctuations take place during inflation, and the expansion smooths out the spatial variation, the kinetic energy constraint is never violated. As long as the potential energy does not grow larger than $`M_P`$, any value of the field can be reached in some domain if inflation goes on long enough. Thus if the initial value of $`\varphi `$ starts off as $`\varphi M_P`$ and $`H,\mathrm{\Lambda }_{\mathrm{other}}`$ are also of order $`M_P`$, one requires $`N=Ht10^{244}`$ e-foldings of inflation in order to have quantum fluctuations allow the field to grow to sufficient size to be relevant for the anthropic constraint. For the case where all of these quantities are as small as 1 TeV, the required number of e-foldings is $`10^{148}`$. If a theory has only 60-100 e-foldings, the quantum fluctuations cannot solve the initial value problem. However the constraint on the amount of inflation can be solved naturally in the various versions of eternal inflation, in which inflating domains continue forever, and our domain can have undergone an unlimited amount of inflation. Therefore, the anthropic mechanism is most naturally embedded in theories of eternal inflation, as in . The multiple domain hypothesis raises the possibility of naturally providing a way to solve the fine tuning problem. The Hubble damping mechanism is interesting because it demonstrates that fields can become frozen at a continuous range of values. The difficulty with the extremely flat potential can be traced back to the the reliance on the Hubble term in the equation of motion to provide the mechanism for freezing the field. At present, $`H`$ is too small to provide much influence on the behavior of fields. It is useful to search for more efficient methods of damping the dynamics. ## IV Kinetic freezing We could ask why, in the previous analysis, we could not simply redefine the scalar field by an overall scale such that its magnitude looks more normal, at the expense of also redefining parameters in the potential. The answer is that the condition that set the scale was the requirement that the kinetic terms be conventionally normalized. This suggests that by playing with the overall factor in front of the kinetic energy, one could also freeze the dynamics. In fact, this idea has been suggested in the context of hyperextended inflation, and a variant has recently been invoked to control the dilaton potential. In this situation, one imagines that nonrenormalizable interactions are present in the action, such that the lagrangian becomes $$=\frac{1}{2}f(\varphi ,\psi )_\mu \varphi ^\mu \varphi +V(\varphi ,\psi )+\mathrm{}$$ (14) The function $`f(\varphi ,\psi )`$ is a unknown function that can depend on $`\varphi `$ and on other fields, here labeled $`\psi `$. In this case, at values where $`f`$ is large, the fields are effectively frozen even if fields are not at the minimum of the potential. This can be seen from the equation of motion for $`\varphi `$ $$f(\ddot{\varphi }+3H\dot{\varphi })+\frac{1}{2}f^{}\dot{\varphi }^2=V^{}(\varphi ).$$ (15) where $$f^{}=\frac{f(\varphi ,\psi )}{\varphi }$$ (16) If $`f`$ or $`f^{}`$ is large, this can be a mechanism for slowing further dynamical evolution, yet it is problematic when applied to the cosmological constant. The goal here is to allow a a more natural size of the potential. Using notation from the Sec. III this means that a potential of size $$V^{}(\varphi )\frac{\mathrm{\Lambda }_{\mathrm{other}}}{M_P}$$ (17) will allow $`V^{}(\varphi )\mathrm{\Delta }\varphi \mathrm{\Lambda }_{\mathrm{other}}`$ for $`\varphi `$ ranging over a natural range $`\mathrm{\Delta }\varphi M_P`$. (Here we will not worry about a few extra powers of ten). Since the smallest reasonable expectation for $`\mathrm{\Lambda }_{\mathrm{other}}`$ is of order the scale of low energy supersymmetry, in the absence of other mechanisms, this means that we need a potential of rough size $$V^{}(\varphi )\frac{1\mathrm{TeV}^4}{M_P}10^{64}M_P^3.$$ (18) Combining the equation of motion with the constraint on $`\dot{\varphi }`$, this means that we need $$f^{}>10^{58}M_P^1$$ (19) While this mechanism may also be used to freeze the fields it is questionable whether it is reasonable to get non-renormalizable terms so large. In an effective field theory description, non renormalizable terms occur as small corrections to the basic theory, due to interactions with degrees of freedom which are much heavier. The expectation of effective field theories, born out in known examples, is that once the nonrenomrmalizable terms are of order unity, we excite the high energy degrees of freedom directly and the theory changes to a new effective theory in which these fields are dynamical variables. It is not natural to achieve such extremely large nonrenormalizable interactions. In fact, one can see that this is related to the mechanism of the previous section in the special situation where $`f`$ either does not depend on other fields $`\psi `$, or these fields are held fixed at the minimum of a potential, $`\psi =<\psi >`$, and an integrability constraint is satisfied. In this case a field redefinition changes the problem exactly back to the situation of the previous section. Define $`\chi `$ $`=`$ $`g(\varphi )`$ (20) $`_\mu \chi `$ $`=`$ $`g^{}(\varphi )_\mu \varphi `$ (21) If we then identify $$g^{}(\varphi )=f^{\frac{1}{2}}(\varphi ,<\psi >)$$ (22) and this can be integrated to obtain $`g`$, the Lagrangian is transformed into $$=\frac{1}{2}_\mu \chi ^\mu \chi +V(g^1(\chi ))+\mathrm{}$$ (23) This is just a conventionally normalized action with a suppressed potential. ## V Radiative damping Finally, what about other forms of damping? It is also possible to damp the motion of a field through the radiation of particles. In effect, a changing field can produce particles, which takes energy out of the field and hence slows down the rate of change. This effect has been studied in Ref and is used in the theory of “warm inflation” Let us consider the equations of motion $$\ddot{\varphi }+(\mathrm{\Gamma }+3H)\dot{\varphi }=V^{}(\varphi )$$ (24) with some unspecified damping $`\mathrm{\Gamma }`$. This structure arises from the coupling of the field to other particles, with the radiation of the other particles damping the dynamics of the field. The proportionality of the damping to $`\dot{\varphi }`$ is indicative that if the field is not changing it does not radiate. Perhaps this mechanism could lead to naturally frozen fields. The difficulty in this case comes from the fact that $`\mathrm{\Gamma }`$ must be very small at present. Can we have $`\mathrm{\Gamma }\dot{\varphi }`$ as large as $`10^{64}M_P^3`$ as required by the constraint of Eq. 18? The constraint on $`\mathrm{\Gamma }`$ comes from the generation of particles in the universe. The equation of motion is equivalent to the conservation of energy in a co-moving volume $`a^3`$ $$\frac{d}{dt}(a^3(t)\rho )=p\frac{d}{dt}a^3\mathrm{\Gamma }\dot{\varphi }^2$$ (25) with energy density and pressure $`\rho `$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2+V(\varphi )`$ (26) $`p`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\varphi }^2V(\varphi )`$ (27) such that $`\mathrm{\Gamma }\dot{\varphi }^2`$ represents the rate of energy flow out of the field $`\varphi `$. Higher power dependence on $`\dot{\varphi }`$ will not change our argument. A reasonably loose constraint on the rate of energy production is that it is smaller that the production of the full present energy in one Hubble time. $$\mathrm{\Gamma }\dot{\varphi }^2<H\rho _010^{183}M_P^5$$ (28) This then lets us put a constraint on the damping term in the equation of motion using $$\mathrm{\Gamma }\dot{\varphi }=(\mathrm{\Gamma }\dot{\varphi }^2\mathrm{\Gamma })^{\frac{1}{2}}<10^{91}M_P^3(\frac{\mathrm{\Gamma }}{M_P})^{1/2}.$$ (29) Even if this unspecified damping mechanism was able to produce $`\mathrm{\Gamma }M_P`$, this fails by 27 orders of magnitude to provide enough damping to allow a reasonably sized potential. ## VI Form fields We may also turn to other ideas for fields with frozen dynamics. Another possibility is known in the supergravity literature, as first pointed out by Consider a field like a gauge potential but with three totally antisymmetric Lorentz indices $$A_{\alpha \beta \gamma }(x)=A_{\beta \alpha \gamma }=A_{\gamma \beta \alpha }$$ (30) such that its field strength tensor is also formed antisymmetrically $$F_{\alpha \beta \gamma \delta }=_{[\alpha }A_{\beta \gamma \delta ]}$$ (31) where the square brackets denote the antisymmetrization of the indices. The Bianchi identity $$_{[\alpha }F_{\beta \gamma \delta \rho ]}=0$$ (32) is then always satisfied in 4 dimensions since there is no totally antisymmetric object with five Lorentz indices. The action $$S_F=\frac{1}{48}d^4x\sqrt{g}F_{\alpha \beta \gamma \delta }F^{\alpha \beta \gamma \delta }$$ (33) leads to the equation of motion $$^\alpha \left[\sqrt{g}F_{\alpha \beta \gamma \delta }\right]=0.$$ (34) The only solution to this is $$F_{\alpha \beta \gamma \delta }=\frac{c}{\sqrt{g}}ϵ_{\alpha \beta \gamma \delta }$$ (35) for arbitrary c. Thus this field is nondynamical, with only a constant solution. Substitution of this solution in Einstein’s equations shows that it behaves as a positive cosmological constant. In the language of differential forms, $`A`$ is a 3-form potential, and $`F`$ a 4-form field strength, with equations of motion and Bianchi identity $`dF_4`$ $`=`$ $`0`$ (36) $`dF_4`$ $`=`$ $`0.`$ (37) Form fields appear in the low energy limit of string theory and M theory. The most obvious is the type II supergravity in the low energy limit of M theory, where the 4-form field strengths occur explicitly. However, they can also be obtained by dimensional reduction from higher form fields. Consider a form field strength with more than four indices, $`F_{\alpha \beta \gamma \delta \mathrm{}.\rho }`$. Upon compactification, some of the indices can be assigned to the compact directions, becoming internal indices. The number of such 4-forms will depend on the particular number and symmetries of the compact subspaces. Four-forms may also appear from lower dimension forms. For an n-form in d dimensions, its dual is a d-n form. Likewise duality relates a 4- form in 4-d to a zero-form -i.e. a constant. Hawking and Turok have proposed the generation of a non-zero 4-form through a tunneling mechanism involving an special instanton in the case of an open universe. This calculation remains controversial, with a dispute over the meaning of the instanton solution. The mechanism has phenomenological problems, as it naturally predicts an almost empty universe. Moreover, the mechanism generates one value of $`\mathrm{\Lambda }`$ through out the entire universe, such that the naturalness of the anthropic selection is lost, and it does not correspond to the multiple domain structure under consideration here. If a multiple domain structure is to be realized in nature, it will be generated in the early universe. Therefore we should look to cosmology for possible mechanisms. Here we suggest a mechanism which exploits the dimensional reduction that may take place in string theories The non-dynamical nature of the 4-form fields is only true in four dimensions. In higher dimensions, the equations of motion allow the usual plane wave solutions. The lack of dynamics in 4-d results from the restriction that the Lorentz indices and the space-time variability lie entirely within the 4-d space. This suggests a potential cosmological mechanism for the generation of the 4-form. Consider higher dimensional theories where compactification leads to a 4-d low energy theory. If cosmology goes through a phase where fields above the compactification scale are excited at some time in the early universe, 4-form fields will be dynamical. They will have fluctuating values, with a non-zero rms field strength. As the universe expands and the average energy decreases, the Kaluza Klein modes with excitations in the compact dimensions will decouple leaving an effective four dimensional theory. As this transition occurs, the 4-form fields will become non-dynamical and will be frozen into random values in different space-time regions. As the universe evolves to lower energies, these values remain frozen. When supplemented by inflation, such that we see only the field from a very small initial patch, this can result in the multiple domain scenario. In string theory there appears to be a barrier to the use of form fields to generate random values of $`\mathrm{\Lambda }`$. In a string theory ground state, the values of the form field strengths are quantized. This occurs because there are both electric and magnetic charges coupled to the form fields. By analogy to the usual electric and magnetic charges, these charges are quantized. Construction of various Gaussian surfaces then imply that the flux, and hence the magnitudes of the constant form fields, are also quantized. The cosmological mechanism described above could also generate different values of the quantized form fields, but it might appear that unless the size of the quanta are extremely small, the likelihood of solving the cosmological constant problem is small<sup>§</sup><sup>§</sup>§Of course it is also possible to imagine form fields without invoking string theory. It has also been argued that there can be form fields which are not coupled to string theory charges.. However the quantization constraint can still allow the form fields to take on all values in a continuous range providing other fields adjust accordingly. The quantization constraint involves $`V_7`$, the volume of the compact seven dimensional manifold. There are also additive contributions from possible flat background gauge potentials and constant fermion densitiesThis can be seen from the supergravity equations of motion.. If these were all to attain their low-energy values first, then the form field condensate would be forced to certain discrete values. However, in the early universe the moduli controlling $`V_7`$, the gauge potentials and the fermions are fluctuating. The form field can take on any continuous value as long as the other fields are adjusted to values consistent with the quantization constraint. As the universe cools to lower energies, the form field will become non-dynamical and will stay at its constant value. At low energies the potentials for the moduli and other fields will become important and will approach their zero-temperature form. These fields will then seek the minimum of their potentials, with the quantization constraint being a constraint on what values are possible. On other words, the form field value will become a constraint on vacuum selection because it is no longer able to evolve. This inverts the usual reasoning, with the result that the form fields could end up at any value but the vacuum state adjusts in order to satisfy the quantization condition. The frozen fields will have two effects. First, they can contribute directly to the cosmological constant. However, there is also an indirect secondary effect through the dilaton and moduli fields. As emphasized in Ref, the form fields which carry string theory charges can influence the potentials for the moduli and dilaton fields. The moduli and dilaton potentials vanish perturbatively, yet it is expected that non-perturbative effects will generate potentials for these fields. The frozen background of form fields will give additive contributions to the potentials. This would amount to random shifts in the moduli potentials in different domains, and would influence the ground state solution and the parameters of the low energy theory. This will then provide a further shift in the ultimate cosmological constant, since every mass and coupling contributes to some extent to the vacuum energy. The influence on the moduli values may lead to the expectation that other parameters in the theory also are variable. In general, non-zero values of the form fields break supersymmetry. It is known that there are special combinations of compactification and vacuum expectation values that allow the existence of low-energy supersymmetry. Whether it is natural that such special situations occur in the early universe is an open question. However, it may even be preferable that the supersymmetry is broken at high scales (depending ultimately on the outcome of future experiments, of course.) In theories with random coupling constants, the fine tuning problem of the Higgs vev may not be the most serious issue. As with the cosmological constant, there is a plausible anthropic constraint such that we would only live in regions with a small Higgs vev. This occurs because if the vev is much larger than observed, the elements other than hydrogen do not exist and we lack the complexity needed for life. The variability of the form fields and the moduli could allow the realization of this anthropic constraint also. Moreover, low energy supersymmetry poses significant problems for cosmology. Scalar particles with TeV scale masses are ubiquitous in such theories, and the dilaton in particular is model independent. These particles dominate the energy density of the universe for so long that they spoils nucleosynthesis. This problem, a string variant of the Polonyi problem, has proven difficult to overcome. Moreover, it appears difficult to implement inflation in theories with low energy supergravity. So cosmology may welcome the situation where supersymmetry is broken at a high scale. Let us then summarize the ingredients of a cosmology that would make use of this mechanism. The first obvious requirement is that the evolution of the universe must involve an early period where energies above the compactification scale are excited. This is needed in order to excite the form fields. The required features of compactification has not yet been studied much because most analyses have been done under the assumption that supersymmetry survives to low energy. So we don’t yet know the full possibility for the field content below the scale of non-supersymmetric compactifications. However, the supersymmetric spectrum above the compactification scale has many fields, the dilaton and moduli, that have the possibility of playing the role of the inflaton. Use of these fields would likely be possible if inflation and compactification occur at the same scale. Finally we clearly need sufficient inflation to smooth out any initial gradients in the fields. ## VII summary This paper is a preliminary investigation into the field theory dynamics that could lead to continuous random contributions to the cosmological constant in theories with multiple domains with different parameters. Damping mechanisms appear to require rather extreme values for the potentials, the fields and/or the nonrenormalizable interactions. As noted by Weinberg, the need to decouple all other fields from this scalar field, in order to preserve the flatness of the potential, has the consequence that it will not influence other parameters in the theory - that the cosmological constant will be the only parameter for which an anthropic constraint is relevant However 4-form fields appear as a quite natural mechanism. For this to be applicable, we would want an energetic initial condition, to excite the form fields, and a inflationary phase to generate the uniformity of the observed universe. The fact that the form fields also influence the dilaton and moduli fields of string theory is also interesting. This would generate a chaotic component to the vacuum selection procedure and would thus influence the other parameters in the theory also. This may then also for the Higgs vev fine-tuning problem. There exists the possibility of testing the distribution of some of the parameters through the weight of the quark mass distribution. It remains to be seen whether a fully complete model along these lines may be developed. This paper has explored the situation in which the field variables influencing the cosmological constant are continuous. In this situation it is quite natural that the cosmological constant should occasionally be close enough to zero to satisfy Weinberg’s anthropic constraint. In a recent paper, Bousso and Polchinski have addressed the situation where multiple form fields can plausibly lead to discrete but closely spaced values for the cosmological constant appropriate for an anthropic selection. The spacing of the values with separation of order $`10^{122}M_P^4`$ appears to require very large internal dimensions or very many (of order 100) form fields. ## Acknowledgments I would like to thank David Kastor, Renata Kallosh, Andrei Linde, Joe Polchinski and especially Daniel Waldram, Jaume Garriga and Alexander Vilenkin for useful conversations. I also thank CERN for their kind hospitality during most of time devoted to this work. This work has been supported in part by the U.S. National Science Foundation and by the John Templeton Foundation.
warning/0006/hep-ex0006008.html	ar5iv	text	# I Measurement results. Displayed are the decay mode, event yield from the fit, total efficiency including secondary branching fraction ϵ, statistical significance (𝜎), branching fraction from the fit ℬ_{𝑓⁢𝑖⁢𝑡} (in units of 10-6), the measured branching fraction (ℬ) or 90% confidence level upper limit (in units of 10-6) and theoretical prediction [] (in units of 10-6). For the branching fraction measurement, the first error is statistical and the second systematic. We assume equal branching fractions for Υ⁢(4⁢𝑆)→"B0¯B0" and 𝐵⁺𝐵⁻. CLNS 99/1652 CLEO 99-19 Study of Charmless Hadronic $`B`$ Meson Decays to Pseudoscalar-Vector Final States CLEO Collaboration (August 4, 2000) Abstract We report results of searches for charmless hadronic $`B`$ meson decays to pseudoscalar($`\pi ^\pm `$, $`K^\pm `$, $`\pi ^0`$ or $`K_S^0`$)-vector($`\rho `$, $`K^{}`$ or $`\omega `$) final states. Using $`9.7\times 10^6`$ $`B\overline{B}`$ pairs collected with the CLEO detector, we report first observation of $`B^{}\pi ^{}\rho ^0`$, $`\overline{B}^0\pi ^\pm \rho ^{}`$ and $`B^{}\pi ^{}\omega `$, which are expected to be dominated by hadronic $`b`$ $``$ $`u`$ transitions. The measured branching fractions are (10.4$`{}_{3.4}{}^{+3.3}\pm 2.1)\times 10^6`$, (27.6$`{}_{7.4}{}^{+8.4}\pm 4.2)\times 10^6`$ and (11.3$`{}_{2.9}{}^{+3.3}\pm 1.4)\times 10^6`$, respectively. Branching fraction upper limits are set for all the other decay modes investigated. PACS numbers: 13.20.He,13.25.-k,13.25.Hw,13.30.Eg,14.40.Nd C. P. Jessop,<sup>1</sup> H. Marsiske,<sup>1</sup> M. L. Perl,<sup>1</sup> V. Savinov,<sup>1</sup> D. Ugolini,<sup>1</sup> X. Zhou,<sup>1</sup> T. E. Coan,<sup>2</sup> V. Fadeyev,<sup>2</sup> Y. Maravin,<sup>2</sup> I. Narsky,<sup>2</sup> R. Stroynowski,<sup>2</sup> J. Ye,<sup>2</sup> T. Wlodek,<sup>2</sup> M. Artuso,<sup>3</sup> R. Ayad,<sup>3</sup> C. Boulahouache,<sup>3</sup> K. Bukin,<sup>3</sup> E. Dambasuren,<sup>3</sup> S. Karamov,<sup>3</sup> G. Majumder,<sup>3</sup> G. C. Moneti,<sup>3</sup> R. Mountain,<sup>3</sup> S. Schuh,<sup>3</sup> T. Skwarnicki,<sup>3</sup> S. Stone,<sup>3</sup> G. Viehhauser,<sup>3</sup> J.C. Wang,<sup>3</sup> A. Wolf,<sup>3</sup> J. Wu,<sup>3</sup> S. Kopp,<sup>4</sup> S. E. Csorna,<sup>5</sup> I. Danko,<sup>5</sup> K. W. McLean,<sup>5</sup> Sz. Márka,<sup>5</sup> Z. Xu,<sup>5</sup> R. Godang,<sup>6</sup> K. Kinoshita,<sup>6,</sup><sup></sup><sup></sup>*Permanent address: University of Cincinnati, Cincinnati, OH 45221 I. C. Lai,<sup>6</sup> S. Schrenk,<sup>6</sup> G. Bonvicini,<sup>7</sup> D. Cinabro,<sup>7</sup> S. McGee,<sup>7</sup> L. P. Perera,<sup>7</sup> G. J. Zhou,<sup>7</sup> E. Lipeles,<sup>8</sup> M. Schmidtler,<sup>8</sup> A. Shapiro,<sup>8</sup> W. M. Sun,<sup>8</sup> A. J. Weinstein,<sup>8</sup> F. Würthwein,<sup>8,</sup>Permanent address: Massachusetts Institute of Technology, Cambridge, M A 02139. D. E. Jaffe,<sup>9</sup> G. Masek,<sup>9</sup> H. P. Paar,<sup>9</sup> E. M. Potter,<sup>9</sup> S. Prell,<sup>9</sup> V. Sharma,<sup>9</sup> D. M. Asner,<sup>10</sup> A. Eppich,<sup>10</sup> T. S. Hill,<sup>10</sup> R. J. Morrison,<sup>10</sup> H. N. Nelson,<sup>10</sup> R. A. Briere,<sup>11</sup> B. H. Behrens,<sup>12</sup> W. T. Ford,<sup>12</sup> A. Gritsan,<sup>12</sup> J. Roy,<sup>12</sup> J. G. Smith,<sup>12</sup> J. P. Alexander,<sup>13</sup> R. Baker,<sup>13</sup> C. Bebek,<sup>13</sup> B. E. Berger,<sup>13</sup> K. Berkelman,<sup>13</sup> F. Blanc,<sup>13</sup> V. Boisvert,<sup>13</sup> D. G. Cassel,<sup>13</sup> M. Dickson,<sup>13</sup> P. S. Drell,<sup>13</sup> K. M. Ecklund,<sup>13</sup> R. Ehrlich,<sup>13</sup> A. D. Foland,<sup>13</sup> P. Gaidarev,<sup>13</sup> L. Gibbons,<sup>13</sup> B. Gittelman,<sup>13</sup> S. W. Gray,<sup>13</sup> D. L. Hartill,<sup>13</sup> B. K. Heltsley,<sup>13</sup> P. I. Hopman,<sup>13</sup> C. D. Jones,<sup>13</sup> D. L. Kreinick,<sup>13</sup> M. Lohner,<sup>13</sup> A. Magerkurth,<sup>13</sup> T. O. Meyer,<sup>13</sup> N. B. Mistry,<sup>13</sup> E. Nordberg,<sup>13</sup> J. R. Patterson,<sup>13</sup> D. Peterson,<sup>13</sup> D. Riley,<sup>13</sup> J. G. Thayer,<sup>13</sup> P. G. Thies,<sup>13</sup> B. Valant-Spaight,<sup>13</sup> A. Warburton,<sup>13</sup> P. Avery,<sup>14</sup> C. Prescott,<sup>14</sup> A. I. Rubiera,<sup>14</sup> J. Yelton,<sup>14</sup> J. Zheng,<sup>14</sup> G. Brandenburg,<sup>15</sup> A. Ershov,<sup>15</sup> Y. S. Gao,<sup>15</sup> D. Y.-J. Kim,<sup>15</sup> R. Wilson,<sup>15</sup> T. E. Browder,<sup>16</sup> Y. Li,<sup>16</sup> J. L. Rodriguez,<sup>16</sup> H. Yamamoto,<sup>16</sup> T. Bergfeld,<sup>17</sup> B. I. Eisenstein,<sup>17</sup> J. Ernst,<sup>17</sup> G. E. Gladding,<sup>17</sup> G. D. Gollin,<sup>17</sup> R. M. Hans,<sup>17</sup> E. Johnson,<sup>17</sup> I. Karliner,<sup>17</sup> M. A. Marsh,<sup>17</sup> M. Palmer,<sup>17</sup> C. Plager,<sup>17</sup> C. Sedlack,<sup>17</sup> M. Selen,<sup>17</sup> J. J. Thaler,<sup>17</sup> J. Williams,<sup>17</sup> K. W. Edwards,<sup>18</sup> R. Janicek,<sup>19</sup> P. M. Patel,<sup>19</sup> A. J. Sadoff,<sup>20</sup> R. Ammar,<sup>21</sup> A. Bean,<sup>21</sup> D. Besson,<sup>21</sup> R. Davis,<sup>21</sup> N. Kwak,<sup>21</sup> X. Zhao,<sup>21</sup> S. Anderson,<sup>22</sup> V. V. Frolov,<sup>22</sup> Y. Kubota,<sup>22</sup> S. J. Lee,<sup>22</sup> R. Mahapatra,<sup>22</sup> J. J. O’Neill,<sup>22</sup> R. Poling,<sup>22</sup> T. Riehle,<sup>22</sup> A. Smith,<sup>22</sup> J. Urheim,<sup>22</sup> S. Ahmed,<sup>23</sup> M. S. Alam,<sup>23</sup> S. B. Athar,<sup>23</sup> L. Jian,<sup>23</sup> L. Ling,<sup>23</sup> A. H. Mahmood,<sup>23,</sup>Permanent address: University of Texas - Pan American, Edinburg, TX 78 539. M. Saleem,<sup>23</sup> S. Timm,<sup>23</sup> F. Wappler,<sup>23</sup> A. Anastassov,<sup>24</sup> J. E. Duboscq,<sup>24</sup> K. K. Gan,<sup>24</sup> C. Gwon,<sup>24</sup> T. Hart,<sup>24</sup> K. Honscheid,<sup>24</sup> D. Hufnagel,<sup>24</sup> H. Kagan,<sup>24</sup> R. Kass,<sup>24</sup> T. K. Pedlar,<sup>24</sup> H. Schwarthoff,<sup>24</sup> J. B. Thayer,<sup>24</sup> E. von Toerne,<sup>24</sup> M. M. Zoeller,<sup>24</sup> S. J. Richichi,<sup>25</sup> H. Severini,<sup>25</sup> P. Skubic,<sup>25</sup> A. Undrus,<sup>25</sup> S. Chen,<sup>26</sup> J. Fast,<sup>26</sup> J. W. Hinson,<sup>26</sup> J. Lee,<sup>26</sup> N. Menon,<sup>26</sup> D. H. Miller,<sup>26</sup> E. I. Shibata,<sup>26</sup> I. P. J. Shipsey,<sup>26</sup> V. Pavlunin,<sup>26</sup> D. Cronin-Hennessy,<sup>27</sup> Y. Kwon,<sup>27,</sup><sup>§</sup><sup>§</sup>§Permanent address: Yonsei University, Seoul 120-749, Korea. A.L. Lyon,<sup>27</sup> and E. H. Thorndike<sup>27</sup> <sup>1</sup>Stanford Linear Accelerator Center, Stanford University, Stanford, California 94309 <sup>2</sup>Southern Methodist University, Dallas, Texas 75275 <sup>3</sup>Syracuse University, Syracuse, New York 13244 <sup>4</sup>University of Texas, Austin, TX 78712 <sup>5</sup>Vanderbilt University, Nashville, Tennessee 37235 <sup>6</sup>Virginia Polytechnic Institute and State University, Blacksburg, Virginia 24061 <sup>7</sup>Wayne State University, Detroit, Michigan 48202 <sup>8</sup>California Institute of Technology, Pasadena, California 91125 <sup>9</sup>University of California, San Diego, La Jolla, California 92093 <sup>10</sup>University of California, Santa Barbara, California 93106 <sup>11</sup>Carnegie Mellon University, Pittsburgh, Pennsylvania 15213 <sup>12</sup>University of Colorado, Boulder, Colorado 80309-0390 <sup>13</sup>Cornell University, Ithaca, New York 14853 <sup>14</sup>University of Florida, Gainesville, Florida 32611 <sup>15</sup>Harvard University, Cambridge, Massachusetts 02138 <sup>16</sup>University of Hawaii at Manoa, Honolulu, Hawaii 96822 <sup>17</sup>University of Illinois, Urbana-Champaign, Illinois 61801 <sup>18</sup>Carleton University, Ottawa, Ontario, Canada K1S 5B6 and the Institute of Particle Physics, Canada <sup>19</sup>McGill University, Montréal, Québec, Canada H3A 2T8 and the Institute of Particle Physics, Canada <sup>20</sup>Ithaca College, Ithaca, New York 14850 <sup>21</sup>University of Kansas, Lawrence, Kansas 66045 <sup>22</sup>University of Minnesota, Minneapolis, Minnesota 55455 <sup>23</sup>State University of New York at Albany, Albany, New York 12222 <sup>24</sup>Ohio State University, Columbus, Ohio 43210 <sup>25</sup>University of Oklahoma, Norman, Oklahoma 73019 <sup>26</sup>Purdue University, West Lafayette, Indiana 47907 <sup>27</sup>University of Rochester, Rochester, New York 14627 $`CP`$ violation in the Standard Model (SM) is a consequence of the complex phase in the Cabibbo-Kobayashi-Maskawa (CKM) quark-mixing matrix . The study of charmless hadronic decays of $`B`$ mesons plays a key role in testing the SM picture of $`CP`$ violation. For example, the angle $`\alpha `$ $``$ $`arg`$ \[($`V_{td}V_{tb}^{})/(V_{ud}V_{ub}^{})`$\] of the unitarity triangle can be measured by performing a full Dalitz analysis of the decays $`B^0(\overline{B}^0)\pi ^+\rho ^{}`$, $`\pi ^{}`$$`\rho ^+`$ and $`\pi ^0`$$`\rho ^0`$ . While the CLEO data do not yet have the sensitivity for the $`CP`$ violation measurements, experimental measurements of these decay modes will be useful to test various theoretical predictions that typically make use of effective Hamiltonians, often with factorization assumptions . Recently, it has been suggested , with model dependency, that published experimental results on charmless hadronic $`B`$ decays indicate that $`\mathrm{cos}\gamma <0`$, in disagreement with current fits to the information most sensitive to CKM matrix elements . In this Letter, we present results of searches for $`B`$ meson decays to exclusive pseudoscalar-vector ($`BPV`$) final states that include a pseudoscalar meson $`\pi ^\pm `$, $`K^\pm `$, $`\pi ^0`$ or $`K_S^0`$ and a vector meson $`\rho `$, $`K^{}`$ or $`\omega `$. In particular we present first observation of the decays $`B^{}\pi ^{}\rho ^0`$, $`\overline{B}^0\pi ^\pm \rho ^{}`$ and $`B^{}\pi ^{}\omega `$ (charge-conjugate modes are implied) which are expected to be dominated by hadronic $`b`$ $``$ $`u`$ transitions. Our results supersede previous CLEO results on these decay modes . The data were collected with two configurations (CLEO II and CLEO II.V ) of the CLEO detector at the Cornell Electron Storage Ring (CESR). They consist of 9.1 fb<sup>-1</sup> taken at the $`\mathrm{{\rm Y}}`$(4S), which corresponds to $`9.7\times 10^6`$ $`B\overline{B}`$ pairs, and 4.4 fb<sup>-1</sup> taken below $`B\overline{B}`$ threshold, used for continuum background studies. The data sample contains a factor of 3 more statistics than previously published results . In addition, the CLEO II data were reanalyzed with improved calibration and track-fitting, allowing for larger geometric acceptance and more efficient track quality requirements. The final states of the decays under study are reconstructed by combining detected photons and charged pions and kaons. The detector elements most important for the results presented here are the tracking system, which consists of several concentric detectors operating inside a 1.5 T superconducting solenoid, and the high-resolution electromagnetic calorimeter, consisting of 7800 CsI(Tl) crystals. For CLEO II, the tracking system consists of a 6-layer straw tube chamber, a 10-layer precision drift chamber, and a 51-layer main drift chamber. The main drift chamber also provides a measurement of the specific ionization loss, $`dE/dx`$, used for particle identification. For CLEO II.V the straw tube chamber was replaced by a 3-layer, double-sided silicon vertex detector, and the gas in the main drift chamber was changed from an argon-ethane to a helium-propane mixture. The resonances in the final state are identified via the decay modes $`\rho \pi \pi `$, $`K^{}K\pi `$ ($`K^0K^+\pi ^{}`$, $`K^+K^+\pi ^0`$) and $`\omega \pi ^+\pi ^{}\pi ^0`$. Reconstructed charged tracks are required to pass quality cuts based on their track fit residuals and impact parameter in both the $`r`$$`\varphi `$ and $`r`$$`z`$ planes, and on the number of main drift chamber measurements. Each event must have a total of at least four such charged tracks. The $`dE/dx`$ measured by the main drift chamber is used to distinguish kaons from pions. Electrons are rejected based on $`dE/dx`$ information and the ratio of the measured track momentum and the associated shower energy in the calorimeter. Muons are rejected by requiring that charged tracks penetrate fewer than seven interaction lengths of material. Pairs of charged tracks used to reconstruct $`K_S^0`$ (via $`K_S^0\pi ^+\pi ^{}`$) are required to have a common vertex displaced from the primary interaction point. The invariant mass of the two charged pions is required to be within two standard deviations ($`\sigma `$) of the known $`K_S^0`$ mass . Furthermore, the $`K_S^0`$ momentum vector, obtained from a kinematic fit of the charged pions’ momenta, is required to point back to the beam spot. To form $`\pi ^0`$ candidates, pairs of photon candidates with an invariant mass within 2.5$`\sigma `$ of the nominal $`\pi ^0`$ mass are kinematically fitted with the mass constrained to the known $`\pi ^0`$ mass . The primary means of identification of $`B`$ meson candidates is through their measured mass and energy. The beam-constrained mass of the candidate is defined as $`M_B\sqrt{E_b^2\|𝐩\|^2}`$, where $`𝐩`$ is the measured momentum of the candidate and $`E_b`$ is the beam energy. The resolution of $`M_B`$ ranges from 2.5 to 3.5 MeV, where the larger resolutions correspond to decay modes with neutral pion(s). The second observable $`\mathrm{\Delta }E`$ is defined as $`\mathrm{\Delta }EE_1+E_2E_b`$, where $`E_1`$ and $`E_2`$ are the energies of the two final state mesons. The resolution of $`\mathrm{\Delta }E`$ is mode dependent. For final states without a neutral pion, the $`\mathrm{\Delta }E`$ resolution is about 20 MeV. For decay modes with one or two energetic neutral pions ($`\overline{B}^0\pi ^\pm \rho ^{}`$, $`\overline{B}^0\pi ^0\rho ^0`$ and $`B^{}\pi ^0\rho ^{}`$ etc), the $`\mathrm{\Delta }E`$ resolution worsens by approximately a factor of 2 or 3 and becomes slightly asymmetric because of energy loss out of the back of the CsI crystals. We accept events with $`M_B`$ $`>`$ 5.2 GeV and $`\|\mathrm{\Delta }E\|`$ $`<`$ 100 to 300 MeV depending on the decay mode. The vector meson $`\rho `$, $`K^{}`$ and $`\omega `$ candidates are required to have masses within 200, 75 and 50 MeV of their known masses , respectively. In the simultaneous analysis of $`\overline{B}^0\pi ^0\rho ^0`$ and $`\pi ^0`$$`K^0`$, the $`\rho ^0`$ or $`K^0`$ candidate is required to have mass between 0.3 GeV to 1.0 GeV under the $`\pi ^+`$$`\pi ^{}`$ decay hypothesis so that both $`\rho ^0`$ and $`K^0`$ enter into the sample. Because of the polarization of the vector meson, the soft decay product from the vector meson may have momentum as low as 150 MeV. To reduce the large combinatoric background from soft $`\pi ^0`$s, only half of the helicity ($``$) range, corresponding to a hard $`\pi ^0`$, is selected when a $`\rho ^+`$ or $`K^+`$ decays to a $`\pi ^+`$$`\pi ^0`$ or $`K^+`$$`\pi ^0`$. The helicity is defined as the cosine of the angle between one of the vector meson decay products in the vector meson rest frame and the direction of the vector meson momentum in the lab frame. The main background comes from continuum $`e^+e^{}`$ $``$ $`q\overline{q}`$, where $`q`$ $`=`$ $`u`$, $`d`$, $`s`$, $`c`$. This background typically exhibits a two-jet structure and can be reduced with event shape criteria. We calculate the angle $`\theta _S`$ ($`\theta _T`$) between the sphericity axis (thrust axis ) of the candidate and the sphericity axis (thrust axis) of the rest of the event. The distribution of $`\mathrm{cos}\theta _S(\theta _T)`$ should be flat for $`B`$ mesons and strongly peaked at $`\pm `$1.0 for continuum background. We require $`\|\mathrm{cos}\theta _S\|`$ $`<`$ 0.8 when there is a $`\rho `$ or $`K^{}`$ meson in the final state, and $`\|\mathrm{cos}\theta _T\|`$ $`<`$ 0.8 when there is a $`\omega `$ meson in the final state. We also form a Fisher discriminant ($``$) with event shape observables . We then perform unbinned maximum-likelihood fits where the likelihood of an event is parameterized by the sum of probabilities for all relevant signal and background hypotheses, with relative weights determined by maximizing the likelihood function ($``$) . The probability of a particular hypothesis is calculated as a product of the probability density functions (PDFs) for each of the input observables. The observables used in the fit are $`\mathrm{\Delta }E`$, $`M_B`$, $``$, $``$ and the invariant mass of the resonance candidate. For final states with the same vector meson but different charged light mesons (pion or kaon), we also use the $`dE/dx`$ measurement of the high-momentum track and fit for both modes simultaneously. Similarly, $`dE/dx`$ measurements of the vector meson decay daughters are used in the simultaneous fit for $`\overline{B}^0\pi ^0\rho ^0`$ and $`\pi ^0`$$`K^0`$. For each decay mode investigated, the signal PDFs are determined with fits to GEANT-based simulation samples. The parameters of the continuum background PDFs are determined with similar fits to simulated continuum samples as well as continuum data. Simulated continuum distributions are in excellent agreement with the data taken below the $`B\overline{B}`$ threshold. Correlations between observables used in the fits are investigated and their effect is found to be negligible. In all cases, the fit includes hypotheses for signal decay modes and the dominant continuum background. Using the PDFs formed by the above observables, signal and continuum background can be well separated. For a few channels where the selected sample contains contributions from other $`B`$ decays, we also include hypotheses for background from other $`B`$ decay modes. These background decay modes can also be separated efficiently from the signal decay modes. We select a sample that contains both $`B^{}\pi ^{}\rho ^0`$, $`K^{}`$$`\rho ^0`$ and some contamination from $`B^{}\pi ^{}\overline{K}^0`$. We then fit simultaneously for $`B^{}\pi ^{}\rho ^0`$, $`K^{}`$$`\rho ^0`$ with and without a $`B^{}\pi ^{}\overline{K}^0`$ contribution. Similarly, we select a sample that contains both $`B^{}\pi ^{}\overline{K}^0`$, $`K^{}`$$`\overline{K}^0`$ with some contamination from $`B^{}\pi ^{}\rho ^0`$, $`K^{}`$$`\rho ^0`$. Then we perform a simultaneous fit for $`B^{}\pi ^{}\overline{K}^0`$, $`K^{}`$$`\overline{K}^0`$ with or without the $`B^{}\pi ^{}\rho ^0`$, $`K^{}`$$`\rho ^0`$ contributions. In both cases the fits with and without the background modes are consistent with each other. For each of the combinations $`\overline{B}^0\pi ^0\rho ^0`$, $`\pi ^0`$$`K^0`$, $`\overline{B}^0\pi ^\pm \rho ^{}`$, $`K^\pm `$$`\rho ^{}`$, and $`B^{}\pi ^{}\omega `$, $`K^{}`$$`\omega `$, contributions from other $`B`$ decays are negligible and we select a common sample to fit for both modes. Finally individual samples are selected and fit for the $`B^{}\pi ^0\rho ^{}`$, $`B^{}\pi ^0K^{}`$, $`\overline{B}^0\pi ^0\omega `$ and $`\overline{B}^0K_S^0\omega `$ searches. The contributions of $`b`$ $``$ $`c`$ and other $`B`$ decays are small in the selected samples of final states containing three tracks or two tracks and a $`\pi ^0`$, and their effects on the signal yields are negligible, except in the samples of $`B^{}\pi ^{}\rho ^0`$, $`K^{}`$$`\rho ^0`$ and $`B^{}\pi ^{}\overline{K}^0`$, $`K^{}`$$`\overline{K}^0`$. Events from $`B^{}D^0\pi ^{}`$ where $`D^0K^\pm \pi ^{},\pi ^+\pi ^{}`$ can enter into these samples and mimic our signal. We therefore impose a 30 MeV ($``$ 4$`\sigma `$) wide $`D^0\pi ^+\pi ^{},K^\pm \pi ^{}`$ invariant mass veto in all the charged track pair combinations. We have also studied background from $`B^{}K^{}\eta ^{}`$, with $`\eta ^{}\rho ^0\gamma `$ . This background has exactly the same final state particles as $`B^{}K^{}\rho ^0`$ with an extra photon. Approximately 3% of this background can pass the selection for the $`B^{}\pi ^{}\rho ^0`$, $`K^{}`$$`\rho ^0`$ sample, therefore we include a component in the fit to describe this contribution. For $`B^{}\pi ^0\rho ^{}`$ and $`B^{}\pi ^0K^{}`$ modes, due to the limited $`\mathrm{\Delta }E`$ resolution for the final state with two neutral pions, the selected sample may contain background from other $`B`$ processes such as $`B\pi a_1`$, $`\rho \rho `$. Table I shows the results of these measurements. The one standard deviation statistical error on the yield is determined by finding the ranges for which the quantity $`\chi ^2=2\mathrm{ln}`$ changes by one unit. We observe significant yields for the decays $`B^{}\pi ^{}\rho ^0`$, $`\overline{B}^0\pi ^\pm \rho ^{}`$, $`B^{}\pi ^{}\omega `$ and $`B^{}\pi ^0\rho ^{}`$. To verify that the yields we observe in $`B`$ meson decays to three-pion final states are indeed due to $`\pi \rho `$ decays, we repeat the standard fit allowing for an additional three-pion “non-resonant” contribution. The PDFs for this contribution are identical to the ones used for $`B\pi \rho `$ signals except that we use PDFs that are constants in the $`\rho `$ mass and $``$. We find that this has no effect on the yield and the significance for $`B^{}\pi ^{}\rho ^0`$ and $`\overline{B}^0\pi ^\pm \rho ^{}`$ signals. Possible contributions from all other $`B`$ processes, including higher mass pseudoscalar-vector decays, were also investigated for these channels and found to be negligible. However, the signal yield for $`B^{}\pi ^0\rho ^{}`$ drops from $`23.7_{7.4}^{+8.4}`$ with a significance of 5.1$`\sigma `$ to $`8.0_{7.9}^{+9.1}`$ events with a significance of only 1$`\sigma `$. We can not rule out the possibility that a significant fraction of the observed yield in $`\pi ^0`$$`\rho ^{}`$ comes from poorly measured processes such as non-resonant $`\pi ^{}`$$`\pi ^0`$$`\pi ^0`$, $`\pi a_1`$ and $`\rho \rho `$ processes . Therefore we calculate a conservative upper limit on the branching fraction assuming that the observed yield is due to $`B^{}\pi ^0\rho ^{}`$ decays only. Fig. 1 shows the likelihood contours from fits to $`B^{}\pi ^{}\rho ^0`$, $`K^{}`$$`\rho ^0`$, $`\overline{B}^0\pi ^\pm \rho ^{}`$, $`K^\pm `$$`\rho ^{}`$ and $`B^{}\pi ^{}\omega `$, $`K^{}`$$`\omega `$. The resulting branching fractions are given in Table I. Fig. 2 shows the $`M_B`$ and $`\mathrm{\Delta }E`$ distributions after further requirements are made on event probability to reduce background. For the remaining processes in Table I we do not consider the signal yields to be significant (i.e. significance drops to less than 3 $`\sigma `$ after all the possible systematics are taken into account), and therefore set 90% C. L. upper limits for their branching fractions. Note that for the $`B^{}K^{}\omega `$ decay mode the additional CLEO II.V data and the re-analysis of CLEO II data no longer support its previously reported observation . However, the combined branching fraction $`(B^{}h^{}\omega `$) = $`(14.3_{3.2}^{+3.6}\pm 2.0)\times 10^6`$ (where $`h`$ = $`K`$ or $`\pi `$) is still consistent with the previous result. Systematic errors are separated into two categories. The first consists of systematic errors in the PDFs, which are determined by varying the PDF parameters within their uncertainty. The second consists of systematic errors associated with event selection and efficiency factors. These are determined with studies of independent data samples. For branching fraction central values, the systematic error is the quadrature sum of the two components. For upper limits, the likelihood function is integrated to find the yield value that corresponds to 90% of the total area. The PDF systematic errors are taken into account in this procedure. The selection efficiency is then reduced by one standard deviation when calculating the final upper limit. As a goodness-of-fit check we compare $`2\mathrm{ln}`$ at the minimum for our fits with expectations from fits to Monte Carlo experiments, and find them to be consistent in all cases. In summary, we have made first observation of the decays $`B^{}\pi ^{}\rho ^0`$, $`\overline{B}^0\pi ^\pm \rho ^{}`$ and $`B^{}\pi ^{}\omega `$. All of these $`\mathrm{\Delta }S=0`$ decay modes are expected to be dominated by hadronic $`b`$ $``$ $`u`$ transitions. We see no significant yields in any of the $`\mathrm{\Delta }S=1`$ transitions. This is in contrast to the corresponding charmless hadronic $`B`$ decays to two pseudo-scalar mesons ($`BPP`$) $`BK\pi `$, $`\pi \pi `$, where $`\mathrm{\Delta }S=1`$ transitions clearly dominate . It indicates that gluonic penguin decays play less of a role in $`BPV`$ decays than in $`BPP`$ decays. This is consistent with theoretical predictions that uses factorization which predicts destructive (constructive) interference between penguin operators of opposite chirality for $`BK\rho `$ ($`BK\pi `$), leading to a rather small (large) penguin contribution in these decays. We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. This work was supported by the National Science Foundation, the U.S. Department of Energy, the Research Corporation, the Natural Sciences and Engineering Research Council of Canada, the A.P. Sloan Foundation, the Swiss National Science Foundation, the Texas Advanced Research Program, and the Alexander von Humboldt Stiftung.
warning/0006/cond-mat0006408.html	ar5iv	text	# Charge-ordering phase transition and order-disorder effects in the Raman spectra of NaV2O5 ## Abstract In the ac polarized Raman spectra of NaV<sub>2</sub>O<sub>5</sub> we have found anomalous phonon broadening, and an energy shift of the low-frequency mode as a function of the temperature. These effects are related to the breaking of translational symmetry, caused by electrical disorder that originates from the fluctuating nature of the V<sup>4.5+</sup> valence state of vanadium. The structural correlation length, obtained from comparisons between the measured and calculated Raman scattering spectra, diverges at $`T<`$ 5 K, indicating the existence of the long-range charge order at very low temperatures, probably at $`T=0`$ K. PACS: 78.30.-j, 78.30.Ly, 64.60.Cn, 78.30.Hv The investigation of the physical phenomena associated to the low-dimensional magnetic structures, like a spin-Peierls (SP) transition discovered in the inorganic compound CuGeO<sub>3</sub> , are of great importance for better understanding of strong electron correlations. A very interesting interplay between spin and charge dynamics results in a phase transition discovered in NaV<sub>2</sub>O<sub>5</sub> . At higher temperatures ($`T>34K`$) the magnetic susceptibility of NaV<sub>2</sub>O<sub>5</sub> is in excellent agreement with the Bonner-Fisher curve for the one-dimensional Heisenberg antiferromagnet. At low temperatures, however, the susceptibility decreases rapidly to zero suggesting a SP transition at T=34 K. On the other hand, the temperature dependent nuclear magnetic resonance (NMR) spectra, showed a change of the vanadium valence across the phase transition, from uniform V<sup>4.5+</sup> to two different V<sup>4+</sup> and V<sup>5+</sup> states. These measurements gave direct evidence for the charge ordering (CO) phase-transition scenario in NaV<sub>2</sub>O<sub>5</sub>. A structural analysis, and the polarized Raman and infrared (IR) spectra of NaV<sub>2</sub>O<sub>5</sub> , are also consistent with the existence of uniform vanadium valence in the high-temperature phase. However, despite extensive experimental and theoretical effort, no consistent picture has yet emerged for the low-temperature phase . None of the models proposed to date, have explained fully the two intimately connected issues: the nature of the phase transition and the low-temperature ground state. Raman spectroscopy is a powerful method for the study of dynamics of solids and it can be used to address these issues. Again, despite a considerable amount of data, the Raman spectra of NaV<sub>2</sub>O<sub>5</sub> are still not completely understood . This holds for the vibrational modes as well as for possible magnetic excitations. In this report, we present the study of the vibrational modes in the Raman spectra of NaV<sub>2</sub>O<sub>5</sub> and analysis of the effects associated with non conservation of the light-scattering wave vector selection rule k=0. We argue that anomalous phonon broadening and frequency shifts, that are observed in the ac spectra as a function of the temperature, are caused by strong electrical disorder due to the fluctuating nature of the V<sup>4.5+</sup> valence state of vanadium. The temperature dependence of the structural correlation length is obtained from the comparison between measured and calculated Raman scattering spectra, and used to characterize the phase transition in NaV<sub>2</sub>O<sub>5</sub>. This analysis shows, that the phase transition at $`T_c=34K`$ corresponds to the onset of the short-range electron correlations, with a true long-range order (CO) at $`T`$0. Thus, the charge order is not static below Tc, in contrast to the conclusions of Lohmann et al . Polarized Raman scattering experiments were performed on NaV<sub>2</sub>O<sub>5</sub> single crystals (size $``$ $`1\times 3\times 1`$ mm<sup>3</sup> along a, b, and c respectively) prepared as described in Refs. 2 and 3. As an excitation source we used the 514.5 nm laser line from an Ar<sup>+</sup> ion laser. The beam, with an average power of 5 mW, was focused (spot diameter $`80\mu `$m) on the (001) and (101) surfaces of the crystals. The spectra were measured in a backscattering geometry using a dilor triple monochromator equipped with a LN<sub>2</sub> cooled CCD camera. NaV<sub>2</sub>O<sub>5</sub> crystallizes above T<sub>c</sub> in the orthorombic centrosymmetric space group Pmmn (D$`{}_{}{}^{13}{}_{2h}{}^{}`$), with two molecules in the unit cell of a size: a=1.1318 nm, b=0.3611nm and c=0.4797 nm. Each vanadium atom, in the average valence state 4.5+, is surrounded by five oxygen atoms forming VO<sub>5</sub> pyramids. These pyramids are mutually connected via common edges and corners to form layers in the ab plane, see Fig. 1. The Na atoms are situated between these layers as intercalants. The structure of NaV<sub>2</sub>O<sub>5</sub> can also be described as an array of parallel ladders (running along the b direction) coupled in a trellis lattice (each rung is made of a V-O-V bond). The vibrational properties of NaV<sub>2</sub>O<sub>5</sub> are studied in great details by measuring the Raman spectra only from (001) surface . In our previous paper we observed an interesting low-frequency structure in the ac scattering geometry which, as we will show below, exhibits an unusual temperature dependence. The polarized Raman scattering spectra from (001) and (101) planes, at temperatures above and below phase-transition temperature, are presented in Fig. 2. The b axis of the crystals is set to be parallel to the laboratory H horizontal) axis. Thus the VV (V equals vertical) polarized scattering geometry from the (101) plane gives the mixture of aa, cc, and ac contributions (upper spectra in Fig. 2). The VV polarized configuration from the (001) plane gives only aa contribution (lower spectra in Fig. 2). The corresponding HH spectra (gives the bb contribution for both planes) are found to be identical. In this way we were able to determine the ac contribution at very low frequencies, otherwise impossible to get. For example, the low quality of the (010) surface, prevent us from direct measurements of the low frequency scattering in the ac geometry. In addition to the phonon modes, and a continuum centered at 650 $`cm^1`$ which belongs to the aa channel, we found three modes in the VV spectra obtained from the (101) planes (Fig. 2), that represent the ac contribution. Their detailed temperature dependencies are shown in inset of Fig. 2. The lowest frequency ac mode is centered at about 107 $`cm^1`$. The mode is asymmetric with a frequency cutoff around 120 cm<sup>-1</sup> (Fig.2 and Fig.3b). As the temperature is decreased below $`T_c`$, the mode becomes more symmetric, and hardens by amount of 10 cm<sup>-1</sup> \[Fig. 2 and Fig. 3b\]. The integrated intensity of this feature increases rapidly by the increasing temperature below $`T_c`$, and becomes constant for $`T>T_c`$ (Fig. 4). The temperature behavior of the frequency and the integrated intensity is similar to what is usually observed for a two-magnon mode . However, the two-magnon origin of this mode can be ruled out since its energy (117 cm<sup>-1</sup>) below the phase transition temperatures is much smaller then $`2\mathrm{\Delta }_s160`$ cm<sup>-1</sup>. The additional argument for excluding the two-magnon scattering process comes from observation of the mode at temperatures as high as $`10\times T_c`$. At these temperatures the two-magnon mode should not be visible in the Raman spectra . A one-magnon scattering process, and other magnetic-related scattering mechanisms, can be also eliminated since we did not find any change in the spectra in the magnetic fields up to 12 T. Our analysis shows that this structure might be related to the low frequency phonon density of states, due to the breakdown of the k=0 conservation rule in the light scattering process. In NaV<sub>2</sub>O<sub>5</sub> the V<sup>4.5+</sup> valence state of vanadium causes the random nature of the coupling between the atomic displacements and the fluctuations of the dielectric susceptibility. Such randomness actually appears because of irregular atomic bonding. The visualization of randomness in NaV<sub>2</sub>O<sub>5</sub> is schematically presented in Fig. 1 and it can be referred as an ”electrical disorder”, . Imagine that we ”freeze” the electrical configuration of the charges in the rungs of the ladders, at certain time in the high-temperature phase. We find random electron configurations among rungs, Fig.1b. The average value for all rung configurations gives the V<sup>4.5+</sup> in which the electron is shared by the V atom at the each end of a rung. Below the phase transition temperature the charges start to order, (V<sup>4.5±α</sup>, $`\alpha 0`$) Fig. 1c, reaching complete ”zig-zag” order at T=0 ($`\alpha =0.5`$), Fig. 1d. The ”zig- zag” phase is presented as a real low-temperature geometry, even though there are some other proposed CO configurations . It will be evident from our results that we were not able to discriminate among the various possible CO patterns, and Fig.1 must be regarded only as illustrative of the electrical disorder. However, the ”zig-zag” charge order is consistent with observed phonons in the Raman and IR spectra of NaV<sub>2</sub>O<sub>5</sub> . If so, even for the perfect plane-wave phonon, for example an acoustic or optical phonon with a finite wave vector k, disorder of the atomic coupling allows inelastic scattering of light from this mode . Then, the light scattering is expected to be proportional to the phonon density of states properly weighted by the coupling constant which is in fact a function of $`\omega `$. The light-scattering cross section is proportional to the Fourier transform of the correlation function of the polarizability fluctuations: $$I_{i,j}(\omega )𝑑td(𝐫_1𝐫_2)e^{i\omega ti𝐤(𝐫_1𝐫_2)}<\delta \chi _i^{}(𝐫_1,t)\delta \chi _j(𝐫_2,0)>,$$ (1) where $`𝐤=𝐤_I𝐤_S`$ and $`\omega =\omega _I\omega _S`$ are wavevector and frequency of phonon(s) that participate in the light scattering process. The $`<\mathrm{}>`$ denotes the equilibrium expectation value. For the ideal crystals, both energy and wavevector conservation rules are fulfilled, and the first-order Raman intensity of phonons is proportional to two delta functions, $`\delta (\omega _I\omega _S\omega )\delta (𝐤q)`$. In disordered crystals the k=0 selection rule is broken due to lack of the translational invariance, and the Stokes part of the Raman scattering intensity is $`I(\omega )\omega /[1+n(\omega )]_jC_j(\omega )g_j(\omega )`$. The $`g_j(\omega )`$ is a density of states for the band j, and the $`n(\omega )`$ is Bose distribution function . $`C_j(\omega )`$ describes the coupling of the light and the vibrational mode $`\omega `$. In a continuum description the fluctuations of the susceptibility result from the elastic strain field $`e_j(\omega ,𝐫)=ke_j(\omega )exp(i\mathrm{𝐤𝐫})`$ of the phonon (plane wave), which is modulated by static fluctuations $`\delta p_{ik}(𝐫)`$ of the elasto-optical constants: $`\delta \chi _i(𝐫)=ϵ^2/(4\pi )[p_{ij}+\delta p_{ij}(𝐫)]e_j(𝐫)`$. In this way, the coupling constant $`C(\omega )k^2𝑑rexp(i\mathrm{𝐤𝐫})F(𝐫)`$ is expressed as a correlation function $`F(𝐫)`$ of the fluctuations $`\delta p`$ of the elastoptical constants, which characterize the electrical disorder. The form of the correlations is usually postulated to be either exponential damping $`exp(r/l_c)`$; or Gaussian damping $`exp(r^2/l_c^2)`$, where $`l_c`$ is correlation length. In the case of NaV<sub>2</sub>O<sub>5</sub> we assume Gaussian damping to describe the electrical correlations. The correlation length may be defined as a length over which electrons in neighboring rungs ”see” each other. In fact, this is just a positional (structural) correlation length. One simplified picture, where the intersite Coulomb interactions are ”switched on” at T=34 K has been previously suggested . Thus, the coupling constant is $`C(\omega )k^2e^{k^2l_c^2/4}`$. The $`\omega `$ dependence of $`C(\omega )`$ comes from the dispersion relation between the frequency and the wavevector \[for example, if $`\omega =ck`$ then $`C(\omega )=\omega /c)^2exp[(\omega /c)^2l_c^2/4`$\]. The same type of the correlation function and the coupling constant have been obtained by Martin and Brenig in their analysis of Brillouin scattering in the amorphous solids. Finally, the normalized Raman intensity is: $$I(\omega )\frac{\omega }{1+n(\omega )}\underset{𝐤}{}k^2e^{k^2l_c^2/4}\delta [\omega \omega (𝐤)],$$ (2) where $`\omega (𝐤)`$ is a phonon dispersion. If we confine our analysis to the energy range of the acoustic or/and low frequency optical phonons ($`\omega 150`$ cm<sup>-1</sup>) the Raman spectrum in the high-T phase is influenced by two contributions: acoustic or optic phonon density of state contribution, $`_k\delta [\omega \omega (𝐤)]`$ and the coupling function $`C(k)`$. Since dispersion curves of the phonons have not yet been measured, we are forced to make assumption about which phonons are involved in the light scattering process. There are two possibilities: (a) Broad feature corresponds to the acoustic phonon with a zone-boundary energy of 117 cm<sup>-1</sup>. Bellow $`T_c`$ this mode is introduced to k=0 by zone folding effect. The x-ray diffraction experiments showed the existence of superlattice reflections below $`T_c`$ with a lattice modulation vector q=(1/2, 1/2, 1/4) . In this case the strong anomaly of the elastic constants is expected at $`T_c`$ and indeed observed by Schwenk et al. . (b) The mode corresponds to low-$`\omega `$ optical phonon, also allowed by the symmetry of the low-T phase, with energy that decreases as a function of the wavevector. According to the lattice dynamical calculations a good candidate for that phonon could be the low-$`\omega `$ $`B_{2g}`$ phonon (active in ac polarized geometry), with Na vibrations mainly along the a axis. By examining the Raman spectra it is difficult to conclude which one of these two assumptions is more appropriate, because of the strong quasi elastic background at low frequencies. However, this choice does not critically influence our conclusions, and for the sake of simplicity we assume the cos k form of the phonon dispersion, $`\omega =\omega _0cosk/2`$, $`\omega _0=117`$ cm<sup>-1</sup>, $`k[0,\pi /a]`$. The calculated Raman spectra are presented in Fig. 3a and compared with measured ones, Fig 3b. Both the shift and the broadening of the mode are in good agreement with the experiment. The spectra are obtained by varying just one parameter $`l_c`$, evaluating equation 2 in one dimension, and by taking the values of k and $`l_c`$ in appropriate units of a-lattice constant. By increasing temperature (from T=0), disorder is introduced and the contribution of the $`C(k)`$ in the spectra becomes important. The increase of the degree of disorder is produced by the increase of the $`k0`$ contributions, directed with $`C(k)`$. Therefore, the broadening of the mode and its energy shift towards lower energies are produced by decreasing $`l_c`$ which is a measure of the degree of disorder; complete disorder is characterized with $`l_c`$ a-interatomic distance and long range-order with $`l_c=\mathrm{}`$. The long-range order solution, $`l_c=\mathrm{}`$ and $`C(k)=0`$ of equation 2, gives the vanishing of the Raman intensity. This is not unphysical. It is just telling us that one has to analyze the Raman spectra at T=0 using equation 1 instead of equation 2. The peak position of the mode as a function of the temperature is shown in Fig. 4a. The circles (full lines) represent the experiment (theory). Since in our calculation of the Raman spectra the temperature does not enter as a parameter, the temperature dependence of the peak maximum is included only through the temperature dependence of the correlation length. It is generally expected that correlation length changes with temperature, and this dependence can be obtained by examining the phonon frequency change close to the phase transition-temperature. The best agreement between the measured and the calculated spectra is obtained by assuming a quadratic relation between the correlation length and the temperature, $`1/l_cT^2`$. Therefore, the peak position can be used as a measure of disorder. Accordingly, a similar temperature dependence is expected for the order parameter $`\alpha `$. Furthermore, we present the comparison between the temperature dependence of the measured and calculated integrated intensities, Fig. 4.b. Here, we also find very good agreement between theory and experiment using the same quadratic relation between $`1/l_c`$ and T, and by introducing a non zero intensity offset at T=0 K ($`I(T=0)/I(T=100)=0.3`$). The temperature dependence of the correlation length is presented in inset of Fig. 4. The particular value of lc that is used to calculate the Raman spectra, does not have deeper physical meaning because of an arbitrary factor in the exponent of equation 2. But its temperature dependence does. The correlation length has a approximately constant value above $`T_c`$ and increases below $`T_c`$. This signals that the critical temperature $`T_c=34`$ K represents the onset of short-range electron correlations. As the temperature is lowered below $`T_c`$, the correlation length rapidly increases indicating the existence of a possible singularity at temperatures below T=5 K. In this case the divergence of the correlation length corresponds to the appearance of the true long-range charge order at very low temperatures (probably at T=0). The short range electron correlations could correspond to intersite Coulomb interaction effects , which become important below 34 K, but the electron correlations should also persist in some form at temperatures above 34 K. Following the same arguments, the change of the spectral shape above 100 K (it is also found that some IR spectral changes occurs around 100 K, ) could be a consequence of an additional phase transition, magnetic in origin for example. In conclusion, we have presented the evidence for the existence of translational symmetry breaking effects in the Raman spectra of NaV<sub>2</sub>O<sub>5</sub>. Non conservation of the light-scattering wave vector selection rule, $`𝐤0`$, is caused by strong electrical disorder due to the fluctuating nature of the valence of vanadium. The temperature dependence of the structural correlation length has been obtained from the comparison between the measured and the calculated Raman scattering spectra. This suggests that the phase transition at $`T_c=34K`$, represents the onset of the short range electron correlations, with true long-range order developing at $`T`$0. Acknowledgments MJK thanks to J.C.Irwin, I.Herbut and P.H.M. van Loosdrecht for helpful discussions and comments. This work is supported by Natural Sciences and Engineering Research Council of Canada. MJK also thanks to MPI-FKF Stuttgart, Germany for partial financial support. $`{}_{}{}^{}mkonstan\mathrm{@}sfu.ca`$
warning/0006/cond-mat0006251.html	ar5iv	text	# Mean-field theory of learning: from dynamics to statics ## I Introduction The mean-field theory was first developed as an approximation to many physical systems in magnetic or disordered materials . However, it is interesting that they become exact in many systems in information processing. The major reason of its success is that when compared with physical systems, these artificial systems have extensive interactions among their components. Hence when one component is considered, the influence of the rest of the system can be regarded as a background satisfying some averaged properties. Learning in large neural networks is a mean-field process since the examples and weights strongly interact with each other during the learning process. Learning is often achieved by defining an energy function which involves a training set of examples. The energy function is then minimized by a gradient descent process with respect to the weights until a steady state is reached. Each of the many weights is thus dependent on each of the many examples and vice versa. This makes it an ideal area for applying mean-field theories. There have been attempts using mean-field theories to describe the dynamics of learning. In batch learning, the same restricted set of examples is provided for each learning step. Using the dynamical mean field theory, early work has been done on the steady-state behavior and asymptotic time scales in perceptrons with binary weights, rather than the continuous weights of more common interest . Much benchmarking of batch learning has been done for linear learning rules such as Hebbian learning or Adaline learning . The work on Adaline learning was further extended to the study of linear perceptrons learning nonlinear rules . However, not much work has been done on the learning of nonlinear rules with continuous weights. In this respect, it is interesting to note the recent attempts using the dynamical replica theory . It approximates the temporal correlations during learning by instantaneous effective macroscopic variables. Further approximations facilitate results for nonlinear learning. However, the rigor of these approximations remain to be confirmed in the general case. Batch learning is different from idealized models of on-line learning of infinite training sets, which has gained much progress . In this model, an independent example is generated for each learning step. Since statistical correlations among the examples can be ignored, the many-body interactions among the examples, and hence among the weights, are absent. Hence they do not address the many-body aspects of the dynamics, which will be discussed here. Nevertheless, this simplification enables the dynamics to be simply described by instantaneous dynamical variables, resulting in a significant reduction in the complexity of analysis, thereby leading to great advances in our understanding of on-line learning. In multilayer perceptrons, for instance, the persistence of a permutation symmetric stage which retards the learning process was well studied. Subsequent proposals to speed up learning were made, illustrating the usefulness of the on-line approach . Here we review models of batch learning where, however, such simplifications are not available. Since the same restricted set of examples is recycled during the learning process, there now exist temporal correlations of the parameters in the learning history. Nevertheless, we manage to consider the learning model as a many-body system. Each example makes a small contribution to the learning process, which can be described by linear response terms in a sea of background examples. Two ingredients are important to our theory: (a) The cavity method – Originally developed as the Thouless-Anderson-Palmer approach to magnetic systems and spin glasses , the method was adopted to learning in perceptrons , and subsequently extended to the teacher-student perceptron , the AND machine , the multiclass perceptron , the committee tree , Bayesian learning and pruned perceptrons . These studies only considered the equilibrium properties of learning, whereas here we are generalizing the method to study the dynamics . It uses a self-consistency argument to compare the evolution of the activation of an example when it is absent or present in the training set. When absent, the activation of the example is called the cavity activation, in contrast to its generic counterpart when it is included in the training set. The cavity method yields macroscopic properties identical to the more conventional replica method . However, since the replica method was originally devised as a technique to facilitate systemwide averages, it provides much less information on the microscopic conditions of the individual dynamical variables. (b) The diagrammatic approach – To describe the difference between the cavity activation and its generic counterpart of an example, we apply linear response theory and use Green’s function to describe how the influence of the added example propagates through the learning history. The Green’s function is represented by a series of diagrams, whose averages over examples are performed by a set of pairing rules similar to those introduced for Adaline learning , as well as in the dynamics of layered networks . Here we take a further step and use the diagrams to describe the changes from cavity to generic activations, as was done in , rather than the evolution of specific dynamical variables in the case of linear rules . Hence our dynamical equations are widely applicable to any gradient-descent learning rule which minimizes an arbitrary cost function in terms of the activation. It fully takes into account the temporal correlations during learning, and is exact for large networks. The study of learning dynamics should also provide further insights on the steady-state properties of learning. In this respect we will review the cavity approach to the steady-state behavior of learning, and the microscopic variables satisfy a set of TAP equations. The approach is particularly transparent when the energy landscape is smooth, i.e., no local minima interfere with the approach to the steady state. However, the picture is valid only when a stability condition (equivalent to the Almeida-Thouless condition in the replica method) is satisfied. Beyond this regime, local minima begin to appear and the energy landscape is roughened. In this case, a similar set of TAP equations remains valid. The physical picture has been presented in ; a more complete analysis is presented here. The paper is organized as follows. In Section 2 we formulate the dynamics of batch learning. In Section 3 we introduce the cavity method and the dynamical equations for the macroscopic variables. In Section 4 we present simulation results which support the cavity theory. In Sections 5 and 6 we consider the steady-state behaviour of learning and generalize the TAP equations respectively to the pictures of smooth and rough energy landscapes, followed by a conclusion in Section 7. The appendices explain the diagrammatic approach in describing the Green’s function, the fluctuation response relation, and the equations for macroscopic parameters in the picture of rough energy landscapes. ## II Formulation Consider the single layer perceptron with $`N1`$ input nodes $`\{\xi _j\}`$ connecting to a single output node by the weights $`\{J_j\}`$ and often, the bias $`\theta `$ as well. For convenience we assume that the inputs $`\xi _j`$ are Gaussian variables with mean 0 and variance 1, and the output state is a function $`f(x)`$ of the activation $`x`$ at the output node, where $`x=\stackrel{}{J}\stackrel{}{\xi }+\theta `$. The training set consists of $`p\alpha N`$ examples which map inputs $`\{\xi _j^\mu \}`$ to the outputs $`\{S_\mu \}(\mu =1,\mathrm{},p)`$. In the case of random examples, $`S_\mu `$ are random binary variables, and the perceptron is used as a storage device. In the case of teacher-generated examples, $`S_\mu `$ are the outputs generated by a teacher perceptron with weights $`\{B_j\}`$ and often, a bias $`\varphi `$ as well, namely $`S_\mu =f(y_\mu )`$; $`y_\mu =\stackrel{}{B}\stackrel{}{\xi }^\mu +\varphi `$. Batch learning is achieved by adjusting the weights $`\{J_j\}`$ iteratively so that a certain cost function in terms of the activations $`\{x_\mu \}`$ and the output $`S_\mu `$ of all examples is minimized. Hence we consider a general cost function $`E=_\mu g(x_\mu ,y_\mu )`$. The precise functional form of $`g(x,y)`$ depends on the adopted learning algorithm. In previous studies, $`g(x,y)=(Sx)^2/2`$ in Adaline learning , and $`g(x,y)=xS`$ in Hebbian learning . To ensure that the perceptron fulfills the prior expectation of minimal complexity, it is customary to introduce a weight decay term. In the presence of noise, the gradient descent dynamics of the weights is given by $$\frac{dJ_j(t)}{dt}=\frac{1}{N}\underset{\mu }{}g^{}(x_\mu (t),y_\mu )\xi _j^\mu \lambda J_j(t)+\eta _j(t),$$ (1) where the prime represents partial differentiation with respect to $`x`$, $`\lambda `$ is the weight decay strength, and $`\eta _j(t)`$ is the noise term at temperature $`T`$ with $$\eta _j(t)=0\mathrm{and}\eta _j(t)\eta _k(s)=\frac{2T}{N}\delta _{jk}\delta (ts).$$ (2) The dynamics of the bias $`\theta `$ is similar, except that no bias decay should be present according to consistency arguments , $$\frac{d\theta (t)}{dt}=\frac{1}{N}\underset{\mu }{}g^{}(x_\mu (t),y_\mu )+\eta _\theta (t).$$ (3) ## III The Cavity Method Our theory is the dynamical version of the cavity method . It uses a self-consistency argument to consider what happens when a new example is added to a training set. The central quantity in this method is the cavity activation, which is the activation of a new example for a perceptron trained without that example. Since the original network has no information about the new example, the cavity activation is random. Here we present the theory for $`\theta =\varphi =0`$, skipping extensions to biased perceptrons. Denoting the new example by the label 0, its cavity activation at time $`t`$ is $`h_0(t)=\stackrel{}{J}(t)\stackrel{}{\xi }^0`$. For large $`N`$, $`h_0(t)`$ is a Gaussian variable. Its covariance is given by the correlation function $`C(t,s)`$ of the weights at times $`t`$ and $`s`$, that is, $`h_0(t)h_0(s)=\stackrel{}{J}(t)\stackrel{}{J}(s)C(t,s)`$, where $`\xi _j^0`$ and $`\xi _k^0`$ are assumed to be independent for $`jk`$. For teacher-generated examples, the distribution is further specified by the teacher-student correlation $`R(t)`$, given by $`h_0(t)y_0=\stackrel{}{J}(t)\stackrel{}{B}R(t)`$. Now suppose the perceptron incorporates the new example at the batch-mode learning step at time $`s`$. Then the activation of this new example at a subsequent time $`t>s`$ will no longer be a random variable. Furthermore, the activations of the original $`p`$ examples at time $`t`$ will also be adjusted from $`\{x_\mu (t)\}`$ to $`\{x_\mu ^0(t)\}`$ because of the newcomer, which will in turn affect the evolution of the activation of example 0, giving rise to the so-called Onsager reaction effects. This makes the dynamics complex, but fortunately for large $`pN`$, we can assume that the adjustment from $`x_\mu (t)`$ to $`x_\mu ^0(t)`$ is small, and linear response theory can be applied. Suppose the weights of the original and new perceptron at time $`t`$ are $`\{J_j(t)\}`$ and $`\{J_j^0(t)\}`$ respectively. Then a perturbation of (1) yields $$\left(\frac{d}{dt}+\lambda \right)(J_j^0(t)J_j(t))=\frac{1}{N}g^{}(x_0(t),y_0)\xi _j^0+\frac{1}{N}\underset{\mu k}{}\xi _j^\mu g^{\prime \prime }(x_\mu (t),y_\mu )\xi _k^\mu (J_k^0(t)J_k(t)).$$ (4) The first term on the right hand side describes the primary effects of adding example 0 to the training set, and is the driving term for the difference between the two perceptrons. The second term describes the many-body reactions due to the changes of the original examples caused by the added example, and is referred to as the Onsager reaction term. One should note the difference between the cavity and generic activations of the added example. The former is denoted by $`h_0(t)`$ and corresponds to the activation in the perceptron $`\{J_j(t)\}`$, whereas the latter, denoted by $`x_0(t)`$ and corresponding to the activation in the perceptron $`\{J_j^0(t)\}`$, is the one used in calculating the gradient in the driving term of (4). Since their notations are sufficiently distinct, we have omitted the superscript 0 in $`x_0(t)`$, which appears in the background examples $`x_\mu ^0(t)`$. The equation can be solved by the Green’s function technique, yielding $$J_j^0(t)J_j(t)=\underset{k}{}𝑑sG_{jk}(t,s)\left(\frac{1}{N}g_0^{}(s)\xi _k^0\right),$$ (5) where $`g_0^{}(s)g^{}(x_0(s),y_0)`$ and $`G_{jk}(t,s)`$ is the weight Green’s function, which describes how the effects of a perturbation propagates from weight $`J_k`$ at learning time $`s`$ to weight $`J_j`$ at a subsequent time $`t`$. In the present context, the perturbation comes from the gradient term of example 0, such that integrating over the history and summing over all nodes give the resultant change from $`J_j(t)`$ to $`J_j^0(t)`$. For large $`N`$ the weight Green’s function can be found by the diagrammatic approach explained in Appendix A. The result is self-averaging over the distribution of examples and is diagonal, i.e. $`lim_N\mathrm{}G_{jk}(t,s)=G(t,s)\delta _{jk}`$, where $$G(t,s)=G^{(0)}(ts)+\alpha 𝑑t_1𝑑t_2G^{(0)}(tt_1)D_\mu (t_1,t_2)g_\mu ^{\prime \prime }(t_2)G(t_2,s).$$ (6) Here the bare Green’s function $`G^{(0)}(ts)`$ is given by $$G^{(0)}(ts)\mathrm{\Theta }(ts)\mathrm{exp}(\lambda (ts)).$$ (7) $`\mathrm{\Theta }`$ is the step function. $`D_\mu (t,s)`$ is the example Green’s function given by $$D_\mu (t,s)=\delta (ts)+𝑑t^{}D_\mu (t,t^{})g_\mu ^{\prime \prime }(t^{})G(t^{},s).$$ (8) Our key to the macroscopic description of the learning dynamics is to relate the activation of the examples to their cavity counterparts, which is known to be Gaussian. Multiplying both sides of (5) and summing over $`j`$, we have $$x_0(t)h_0(t)=𝑑sG(t,s)g_0^{}(s).$$ (9) In turn, the covariance of the cavity activation distribution is provided by the fluctuation-response relation explained in Appendix B, $$C(t,s)=\alpha 𝑑t^{}G^{(0)}(tt^{})g_\mu ^{}(t^{})x_\mu (s)+2T𝑑t^{}G^{(0)}(tt^{})G(s,t^{}).$$ (10) Furthermore, for teacher-generated examples, its mean is related to the teacher-student correlation given by $$R(t)=\alpha 𝑑t^{}G^{(0)}(tt^{})g_\mu ^{}(t^{})y_\mu .$$ (11) For a given teacher activation $`y`$ of a trained example, the distribution for a set of student activation $`\{x(t)\}`$ of the same example at different times is, in the limit of infinitesimal time steps $`\mathrm{\Delta }t`$, given by $`p(\{x(t)\}\|y)={\displaystyle \frac{1}{\sqrt{\mathrm{det}C}}}{\displaystyle \underset{t}{}}{\displaystyle \frac{dh(t)}{\sqrt{2\pi }}\mathrm{exp}\left\{\frac{1}{2}\underset{t}{}[h(t)R(t)y]C(t,s)^1[h(s)R(s)y]\right\}}`$ (12) $`{\displaystyle \underset{t}{}}\delta \left[x(t)h(t)\mathrm{\Delta }t{\displaystyle \underset{s}{}}G(t,s)g^{}(x(s))\right].`$ (13) This can be written in an integral form which is often derived from path integral approaches, $`p(\{x(t)\}\|y)={\displaystyle \underset{t}{}}{\displaystyle \frac{dh(t)d\widehat{h}(t)}{2\pi }\mathrm{exp}\left\{i𝑑t\widehat{h}(t)[h(t)R(t)y]\frac{1}{2}𝑑t𝑑s\widehat{h}(t)C(t,s)\widehat{h}(s)\right\}}`$ (14) $`{\displaystyle \underset{t}{}}\delta \left[x(t)h(t)\mathrm{\Delta }t{\displaystyle \underset{s}{}}G(t,s)g^{}(x(s))\right].`$ (15) The above distributions and parameters are sufficient to describe the progress of learning. Some common performance measures used for such monitoring purpose include: (a) Training error $`ϵ_t`$, which is the probability of error for the training examples, and can be determined from the distribution $`p(x\|y)`$ that the student activation of a trained example takes the value $`x`$ for a given teacher activation $`y`$ of the same example. (b) Test error $`ϵ_{test}`$, which is the probability of error when the inputs $`\xi _j^\mu `$ of the training examples are corrupted by an additive Gaussian noise of variance $`\mathrm{\Delta }^2`$. This is a relevant performance measure when the perceptron is applied to process data which are the corrupted versions of the training data. When $`\mathrm{\Delta }^2=0`$, the test error reduces to the training error. Again, it can be determined from $`p(x\|y)`$, since the noise merely adds a variance of $`\mathrm{\Delta }^2C(t,t)`$ to the activations. (c) Generalization error $`ϵ_g`$ for teacher-generated examples, which is the probability of error for an arbitrary input $`\xi _j`$ when the teacher and student outputs are compared. It can be determined from $`R(t)`$ and $`C(t,t)`$ since, for an example with teacher activation $`y`$, the corresponding student activation is a Gaussian with mean $`R(t)y`$ and variance $`C(t,t)`$. ## IV Simulation results The success of the cavity approach is illustrated by the many results presented previously for the Adaline rule . This is a common learning rule and bears resemblance with the more common back-propagation rule. Theoretically, its dynamics is particularly convenient for analysis since $`g^{\prime \prime }(x)=1`$, rendering the weight Green’s function time translation invariant, i.e. $`G(t,s)=G(ts)`$. In this case, the dynamics can be solved by Laplace transform. The closed form of the Laplace solution for Adaline learning enables us to examine a number of interesting phenomena in learning dynamics. For example, an overtraining with respect to the generalization error $`ϵ_g`$ occurs when the weight decay is not sufficiently strong, i.e., $`ϵ_g`$ attains a minimum at a finite learning time before reaching a higher steady-state value. Overtraining of the test error $`ϵ_{test}`$ also sets in at a sufficiently weak weight decay, which is approximately proportional to the noise variance $`\mathrm{\Delta }^2`$. We also observe an equivalence between average dynamics and noiseless dynamics, namely that a perceptron constructed using the thermally averaged weights is equivalent to the perceptron obtained at a zero noise temperature. All these results are well confirmed by simulations. Rather than further repeating previous results, we turn to present results which provide more direct support to the cavity method. In the simulational experiment in Fig. 1, we compare the evolution of two perceptrons $`\{J_j(t)\}`$ and $`\{J_j^0(t)\}`$ in Adaline learning. At the initial state $`J_j^0(0)J_j(0)=1/N`$ for all $`j`$, but otherwise their subsequent learning dynamics are exactly identical. Hence the total sum $`_j(J_j^0(t)J_j(t))`$ provides an estimate for the averaged Green’s function $`G(t,0)`$, which gives an excellent agreement with the Green’s function obtained from the cavity method. Using the Green’s function computed from Fig. 1, we can deduce the cavity activation for each example by measuring their generic counterpart from the simulation and substituting back into Eq. (9). As shown in the histogram in Fig. 2(a), the cavity activation distribution agrees well with the Gaussian distribution predicted by the cavity method, with the predicted mean 0 and variance $`C(t,t)`$. Similarly, we show in Fig. 2(b) the distribution of $`h\mathrm{sgn}y`$, i.e., the cavity activation in the direction of the correct teacher output, The cavity method predicts a Gaussian distribution with mean $`\sqrt{2/\pi }R(t)`$ and variance $`C(t,t)2R(t)^2/\pi `$. Again, it agrees well with the histogram obtained from simulation. ## V Steady-state behavior When learning reaches a steady state at $`T=0`$, the cavity and generic activations approach a constant. Hence Eq. (9) reduces to $$x_0h_0=\gamma g^{}(x_0);\gamma =𝑑sG(t,s),$$ (16) where $`\gamma `$ is called the local susceptibility in . Hence $`x_0`$ is a well-defined function of $`h_0`$. Eq. (16) can also be obtained by minimizing the change in the steady-state energy function when example 0 is added, which is $`g(x_0)+(x_0h_0)^2/2\gamma `$, the second term being due to the reaction effects of the background examples. This was shown in for the case of a constant weight magnitude, but the same could be shown for the case of a constant weight decay. A self-consistent expression for $`\gamma `$ can be derived from the steady-state behavior of the Green’s function. Since the system becomes translational invariant in time at the steady state, Eqs. (6) and (8) can be solved by Laplace transform, yielding $`\stackrel{~}{G}(z)=\stackrel{~}{G}^{(0)}(z)+\alpha \stackrel{~}{G}^{(0)}(z)\stackrel{~}{D}_\mu (z)g_\mu ^{\prime \prime }\stackrel{~}{G}(z),`$ (17) $`\stackrel{~}{D}_\mu (z)=1+\stackrel{~}{D}_\mu (z)g_\mu ^{\prime \prime }\stackrel{~}{G}(z),`$ (18) with $`\stackrel{~}{G}^{(0)}(z)=(z+\lambda )^1`$. Identifying $`\stackrel{~}{G}(0)`$ with $`\gamma `$, we obtain $$\gamma =\frac{1}{\lambda }+\frac{\alpha }{\lambda }\frac{\gamma g_\mu ^{\prime \prime }}{1\gamma g_\mu ^{\prime \prime }}.$$ (19) Making use of the functional relation between $`x_\mu `$ and $`h_\mu `$, we have $$\gamma =\frac{1}{\lambda }(1\alpha \chi );\chi =1\frac{x_\mu }{h_\mu },$$ (20) where $`\chi `$ is called the nonlocal susceptibility in . At the steady state, the fluctuation response relations in Eqs. (10) and (11) yield the self-consistent equations for the student-student and teacher-student correlations, $`C\stackrel{}{J}\stackrel{}{J}`$ and $`R\stackrel{}{J}\stackrel{}{B}`$ respectively, namely $$C=\frac{\alpha }{\lambda }g_\mu ^{}x_\mu ;R=\frac{\alpha }{\lambda }g_\mu ^{}y_\mu .$$ (21) Substituting Eqs. (16) and (20), and introducing the cavity activation distributions, we find $`C=(1\alpha \chi )^1\alpha {\displaystyle DyDhP(h\|y)(x(h)h)x(h)},`$ (22) $`R=(1\alpha \chi )^1\alpha {\displaystyle DyDhP(h\|y)(x(h)h)y}.`$ (23) Since $`P(h\|y)`$ is a Gaussian distribution with mean $`Ry`$ and variance $`CR^2`$, its derivatives with respect to $`h`$ and $`R`$ are $`(hRy)P(h\|y)/(CR^2)`$ and $`R(hRy)P(h\|y)/(CR^2)`$ respectively. This enables us to use integration by parts and Eq. (20) for $`\chi `$ to obtain $`C=\alpha {\displaystyle DyDhP(h\|y)(x(h)h)^2},`$ (24) $`R=\alpha \gamma {\displaystyle DyDhP(h\|y)\frac{x}{h}g_{xy}}.`$ (25) Hence we have recovered the macroscopic parameters described by the static version of the cavity method in by considering the steady-state behavior of the learning dynamics. We remark that the saddle point equations in the replica method also produce identical results, although the physical interpretation is less transparent . We can further derive the microscopic equations by noting that at equilibrium for $`T=0`$, Eq. (1) yields $$J_j=\frac{1}{\lambda N}\underset{\mu }{}g_\mu ^{}\xi _j^\mu ,$$ (26) which leads to the set of equations $$x_\nu =\frac{1}{\lambda }\underset{\mu }{}g_\mu ^{}Q_{\mu \nu };Q_{\mu \nu }\frac{1}{N}\underset{j}{}\xi _j^\mu \xi _j^\nu .$$ (27) The TAP equations are obtained by expressing these equations in terms of the cavity activations via Eq. (16), $$h_\nu =\underset{\mu \nu }{}(x(h_\mu )h_\mu )Q_{\mu \nu }+\alpha \chi x(h_\nu ).$$ (28) The iterative solution of the equation set was applied to the maximally stable perceptron, which yielded excellent agreement with the cavity method, provided that the stability condition discussed below is satisfied . However, the agreement is poorer when applied to the committee tree and the pruned perceptron , where the stability condition is not satisfied. To study the stability condition of the cavity solution, we consider the change in the steady-state solution when example 0 is added to the training set. Consider the magnitude of the displaced weight vector $`\mathrm{\Delta }_j(J_j^0J_j)^2`$. Using either the static or dynamic version of the cavity method, we can show that $$\mathrm{\Delta }=\frac{1}{N}\frac{(x_0h_0)^2}{1\alpha \left(1\frac{x_\mu }{h_\mu }\right)^2}.$$ (29) In order that the change due to the added example is controllable, the stability condition is thus $$\alpha \left(1\frac{x_\mu }{h_\mu }\right)^2<1.$$ (30) This is identical to the stability condition of the replica-symmetric ansatz in the replica method, the so-called Almeida-Thouless condition . As a corollary, when a band gap exists in the activation distribution, the stability condition is violated. This is because the function $`x(h)`$ becomes discontinuous in this case, implying the presence of a delta-function component in $`x/h`$. Such is the case in the nonlinear perceptron trained with noisy examples using the backpropagation algorithm . For insufficient examples and weak weight decay, the activation distribution exhibits a gap for the more difficult examples, i.e., when the teacher output $`y`$ and the cavity activation $`h`$ has a large difference. As shown in Fig. 3(a), simulational and theoretical predictions of the activation distributions agree well in the stable regime, but the agreement is poor in the unstable regime shown in Fig. 3(b). Hence the existence of band gaps necessitates the picture of a rough energy landscape, as described in the following section. ## VI The picture of rough energy landscapes To consider what happens beyond the stability regime, one has to take into account the rough energy landscape of the learning space. To keep the explanation simple, we consider the learning of examples generated randomly, the case of teacher-generated examples being similar though more complicated. Suppose that the original global minimum for a given training set is $`\alpha `$. In the picture of a smooth energy landscape, the network state shifts perturbatively after adding example 0, as schematically shown in Fig. 4(a). In contrast, in the picture of a rough energy landscape, a nonvanishing change to the system is induced, and the global minimum shifts to the neighborhood of the local minimum $`\beta `$, as schematically shown in Fig. 4(b). Hence the resultant activation $`x_0^\beta `$ is no longer a well-defined function of the cavity activation $`h_0^\alpha `$. Instead it is a well-defined function of the cavity activation $`h_0^\beta `$. Nevertheless, one may expect that correlations exist between the states $`\alpha `$ and $`\beta `$. Let $`q_0`$ be the correlation between two local minima labelled by $`\beta `$ and $`\gamma `$, i.e. $`\stackrel{}{J}^\beta \stackrel{}{J}^\gamma =q_0`$. Both of them are centred about the global minimum $`\alpha `$, so that $`\stackrel{}{J}^\alpha \stackrel{}{J}^\beta =\stackrel{}{J}^\alpha \stackrel{}{J}^\gamma =\sqrt{q_0q_1}`$, where $`q_1=\stackrel{}{J}^\alpha \stackrel{}{J}^\alpha =\stackrel{}{J}^\beta \stackrel{}{J}^\beta =\stackrel{}{J}^\gamma \stackrel{}{J}^\gamma `$. Since both states $`\alpha `$ and $`\beta `$ are determined in the absence of the added example 0, the correlation $`h_0^\alpha h_0^\beta =\sqrt{q_0q_1}`$ as well. Knowing that both $`h_0^\alpha `$ and $`h_0^\beta `$ obey Gaussian distributions, the cavity activation distribution can be determined if we know the prior distribution of the local minima. At this point we introduce the central assumption in the cavity method for rough energy landscapes: we assume that the number of local minima at energy $`E`$ obeys an exponential distribution $$d\mathrm{}(E)\mathrm{exp}(wE)dE.$$ (31) Similar assumptions have been used in specifying the density of states in disordered systems . Thus the cavity activation distribution is given by $$P(h_0^\beta \|h_0^\alpha )=\frac{G(h_0^\beta \|h_0^\alpha )\mathrm{exp}[w\mathrm{\Delta }E(x(h_0^\beta ))]}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )\mathrm{exp}[w\mathrm{\Delta }E(x(h_0^\beta ))]},$$ (32) where $`G(h_0^\beta \|h_0^\alpha )`$ is a Gaussian distribution with mean $`\sqrt{q_0/q_1}h_0^\alpha `$ and variance $`q_1q_0`$. $`\mathrm{\Delta }E`$ is the change in energy due to the addition of example 0, and is equal to $`g(x_0^\beta )+(x_0^\beta h_0^\beta )^2/2\gamma `$. The weights $`J_j^\beta `$ are given by $$J_j^\beta =\frac{1}{\lambda N}\underset{\mu }{}g^{}(x_\mu ^\beta )\xi _j^\mu .$$ (33) Self-consistent equations for the macroscopic parameters are derived in Appendix C. The results are identical to the first step replica symmetry-breaking solution in the replica method. It remains to check whether the microscopic equations have been modified due to the roughening of the energy landscape. In terms of the generic activations, the microscopic equations are identical to Eq. (27) for each local minimum. In terms of the cavity activations, the TAP equations are again identical to Eq. (28), except that the nonlocal susceptibility $`\chi `$ is now evaluated in the corresponding local minimum. The cavity activation distribution is no longer a Gaussian distribution, but is modified by the density of states in Eq. (32) now. Hence the values of $`\chi `$ and $`\gamma `$ appearing in the TAP equations are no longer identical to the case of restricting learning to a single valley. ## VII Conclusion In summary, we have introduced a general framework for modeling the dynamics of learning based on the cavity method, which is applicable to general learning cost functions, though its tractable solutions are not generally available. We have verified its validity by simulations of the cavity activation distributions. The steady-state behavior is seen to be consistent with the static version of the cavity method in the picture of smooth energy landscapes, which is equivalent to the replica symmetric ansatz in the replica method. This picture is based on the assumption that the dynamics is stable against perturbations, and is manifested in a stability condition equivalent to the Almeida-Thouless condition in the replica method. Beyond the stability regime, rough energy landscapes have to be introduced, but the microscopic TAP equations remain valid. There are two interesting issues concerning the extension of the present work. First, it is interesting to consider how the dynamics is modified in the picture of rough energy landscapes. In this case, aging effects may appear, and the dynamics may not be translationally invariant in time . Second, it is interesting to consider whether the analysis remains tractable for nonlinear learning rules. In general, $`D_\mu (t,s)`$ in (8) has to be expanded as a series. Nevertheless, we have shown that the asymptotic dynamics remains tractable for nonlinear learning rules. For transient dynamics, we may need to consider appropriate approximations. Another applicable area is the case of batch learning with very large learning steps, whose analysis remains simple due to its fast convergence . The method can also be applied to on-line learning of restricted sets of examples. An alternative general theory for learning dynamics is the dynamical replica theory . It yields exact results for Hebbian learning, but for less trivial cases, the analysis is approximate and complicated by the need to solve replica saddle point equations at every learning instant. It is hoped that by adhering to an exact formalism, the cavity method can provide more fundamental insights when extended to multilayer networks. We thank A. C. C. Coolen and D. Saad for fruitful discussions. This work was supported by the Research Grant Council of Hong Kong (HKUST6130/97P and HKUST6157/99P). ## A The Green’s function Substituting Eq. (5) into Eq. (4), we see that the Green’s function satisfies $$\left(\frac{d}{dt}+\lambda \right)G_{jk}(t,s)=\delta (ts)\delta _{jk}+\frac{1}{N}\underset{\mu i}{}\xi _j^\mu g_\mu ^{\prime \prime }(t)\xi _i^\mu G_{ik}(t,s).$$ (A1) Introducing the bare Green’s function $`G^{(0)}(ts)`$ in Eq. (7), $$G_{jk}(t,s)=G^{(0)}(ts)\delta _{jk}+\frac{1}{N}\underset{\mu i}{}𝑑t^{}G^{(0)}(tt^{})\xi _j^\mu g_\mu ^{\prime \prime }(t^{})\xi _i^\mu G_{ik}(t^{},s).$$ (A2) This equation is represented diagrammatically in Fig. 5(a). We use a slanted line to represent an example bit, the top and bottom ends of the line corresponding to the example label and node label respectively. A filled circle represents $`g_\mu ^{\prime \prime }(t)`$. Thin and thick lines represent the bare and dressed Green’s functions $`G^{(0)}(ts)`$ and $`G(t,s)`$ respectively. The iterative solution to Eq. (A2) can be represented by the series of diagrams in Fig. 5(b). It is convenient to concurrently introduce the example Green’s function $`D_\mu (t,s)`$ as shown in Fig. 5(c). The average over the distribution of example inputs is done by pairing of example or node labels and are represented by dashed lines connecting the vertices above or below the solid lines. Pairing of example and node labels yield factors of 1 and $`\alpha `$ respectively. Noting that crossing diagrams do not contribute , the two Green’s functions can be expressed in terms of the self-energies $`\mathrm{\Sigma }`$ and $`\mathrm{\Pi }_\mu `$, via the Dyson’s equations in Fig. 5(d). The self-energies are defined in Fig. 5(e), and are characterized by having the first node or example paired with the last one only. The self-energies can in turn be expressed in terms of the Green’s functions as in Fig. 5(f), thus allowing for self-consistent solutions. After eliminating the self-energies, the results of the diagrammatic analysis are given by Eqs. (6) and (8). ## B The fluctuation response relation In terms of the bare Green’s function, the solution to the dynamical equation Eq. (1) is $$J_j(t)=\frac{1}{N}\underset{\mu }{}𝑑t^{}G^{(0)}(tt^{})g_\mu ^{}(t^{})\xi _j^\mu +𝑑t^{}G^{(0)}(tt^{})\eta _j(t^{}).$$ (B1) Multiplying both sides by $`J_j(s)`$ and summing over $`j`$, we have $$C(t,s)=\alpha 𝑑t^{}G^{(0)}(tt^{})g_\mu ^{}(t^{})x_\mu (s)+𝑑t^{}G^{(0)}(tt^{})\underset{j}{}J_j(s)\eta _j(t^{}).$$ (B2) The correlation between $`J_j(s)`$ and $`\eta _j(t^{})`$ can be considered by comparing the learning process with another one which is noiseless between $`t^{}ϵ`$ and $`t^{}+ϵ`$, but is otherwise identical. Denoting the weight of this alternative process by $`J_j^{\backslash \eta (t^{})}`$, we have $$J_j(s)=J_j^{\backslash \eta (t^{})}(s)+_{t^{}ϵ}^{t^{}+ϵ}𝑑t^{\prime \prime }G(s,t^{\prime \prime })\eta _j(t^{\prime \prime }).$$ (B3) Noting that $`J_j^{\backslash \eta (t^{})}(s)`$ is uncorrelated with $`\eta _j(t^{})`$, and $`\eta _j(t^{\prime \prime })`$ has a delta function correlation with $`\eta _j(t^{})`$ as in Eq. (2), we arrive at Eq. (10). Similarly, multiplying both sides by $`B_j`$ and summing over $`j`$, we arrive at Eq. (11). ## C Macroscopic parameters in rough energy landscapes From Eq. (20), the nonlocal susceptibility is given by $$\chi =𝑑h_0^\alpha G(h_0^\alpha )\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(1x_0^\beta /h_0^\beta )}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}},$$ (C1) where $`G(h_0^\alpha )`$ is a Gaussian with mean 0 and variance $`q_1`$. The local susceptibility $`\gamma `$ is given by $$\gamma =\frac{1}{\lambda (1\alpha \chi )}.$$ (C2) From the fluctuation response relation in Eq. (10), we have $$q_1=\frac{\alpha }{\lambda }𝑑h_0^\alpha G(h_0^\alpha )\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}g^{}(x_0^\beta )x_0^\beta }{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}},$$ (C3) Substituting Eqs. (16) and (20), we find $$(1\alpha \chi )q_1=\alpha 𝑑h_0^\alpha G(h_0^\alpha )\frac{dh_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )x_0^\beta )}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}.$$ (C4) The differentiations of $`G(h_0^\beta \|h_0^\alpha )`$ with respect to $`h_0^\beta `$ and $`h_0^\alpha `$ introduce factors of $`(h_0^\beta \sqrt{q_0/q_1}h_0^\alpha )/(q_1q_0)`$ and $`\sqrt{q_0/q_1}(h_0^\beta \sqrt{q_0/q_1}h_0^\alpha )/(q_1q_0)`$ respectively, and that of $`G(h_0^\alpha )`$ with respect to $`h_0^\alpha `$ introduces $`h_0^\alpha /q_1`$. This allows us to use integration by parts and Eq. (20) for $`\chi `$ to obtain $`q_1=\alpha \left[1+{\displaystyle \frac{w}{\gamma }}(q_1q_0)\right]{\displaystyle 𝑑h_0^\alpha G(h_0^\alpha )\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )^2}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}}`$ (C5) $`+\alpha {\displaystyle \frac{w}{\gamma }}q_0{\displaystyle 𝑑h_0^\alpha G(h_0^\alpha )\left\{\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )^2}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}\left[\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}\right]^2\right\}}.`$ (C6) Next we derive an equation for the interstate overlap $`q_0`$. Consider the steady-state solution of a local minimum $`J_j^\beta `$ given by Eq. (16). Multiplying both sides by the weight vector $`J_j^\gamma `$ at another local minimum and summing over $`j`$, we have $$q_0=\frac{1}{\lambda N}\underset{\mu }{}g^{}(x_\mu ^\beta )x_\mu ^\gamma .$$ (C7) Proceeding as in the case of $`q_1`$, we get $`q_0=\alpha \left[1+{\displaystyle \frac{w}{\gamma }}(q_1q_0)\right]{\displaystyle 𝑑h_0^\alpha G(h_0^\alpha )\left[\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}\right]^2}`$ (C8) $`+\alpha {\displaystyle \frac{w}{\gamma }}q_0{\displaystyle 𝑑h_0^\alpha G(h_0^\alpha )\left\{\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )^2}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}\left[\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}\right]^2\right\}}.`$ (C9) Solving Eqs. (C6) and (C9), $`{\displaystyle 𝑑h_0^\alpha G(h_0^\alpha )\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )^2}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}}={\displaystyle \frac{q_1+\frac{w}{\gamma }(q_1q_0)^2}{\alpha \left[1+\frac{w}{\gamma }(q_1q_0)\right]^2}},`$ (C10) $`{\displaystyle 𝑑h_0^\alpha G(h_0^\alpha )\left[\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}(x_0^\beta h_0^\beta )}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}\right]^2}={\displaystyle \frac{q_0}{\alpha \left[1+\frac{w}{\gamma }(q_1q_0)\right]^2}}.`$ (C11) To determine the distribution of local minima, namely the parameter $`w`$, we introduce a “free energy” $`F(p,N)`$ for $`p`$ examples and $`N`$ input nodes, given by $$d\mathrm{}(E)=\mathrm{exp}[w(F(p,N)E)]dE.$$ (C12) This “free energy” determines the averaged energy of the local minima and should be an extensive quantity, i.e. it should scale as the system size. Cavity arguments enable us to find an expression $`F(p+1,N)F(p,N)`$. When the number of examples increases by 1, the density of states for a given $`h_0^\alpha `$ are related by $$\mathrm{}(E_{p+1},h_0^\alpha )=𝑑E_p\mathrm{}(E_p,h_0^\alpha )𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )\delta (E_{p+1}E_p\mathrm{\Delta }E).$$ (C13) Using Eq. (C12) we obtain, on averaging over $`h_0^\alpha `$, $$F(p+1,N)=F(p,N)\frac{1}{w}𝑑h_0^\alpha G(h_0^\alpha )\mathrm{ln}𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}.$$ (C14) Similarly, we may consider a cavity argument for the addition of one input node, expanding the network size from $`N`$ to $`N+1`$. Skipping the details, the final result is $$F(p,N+1)F(p,N)=\frac{q_0}{2\gamma \left[1+\frac{w}{\gamma }(q_1q_0)\right]}\frac{1}{2w}\mathrm{ln}\left[1+\frac{w}{\gamma }(q_1q_0)\right]+\frac{\lambda }{2}q_1.$$ (C15) Since $`F`$ is an extensive quantity, $`F(p,N)`$ should scale as $`N`$ for a given ratio $`\alpha =p/N`$. This implies $$\frac{F}{N}=\frac{F}{N}=(F(p,N+1)F(p,N))+\alpha (F(p+1,N)F(p,N)).$$ (C16) When $`E=F`$, the density of states reduces to $`O(e^0)`$ and the global minimum is reached. Hence $`{\displaystyle 𝑑h_0^\alpha G(h_0^\alpha )\frac{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}g(x_0^\beta )}{𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}}={\displaystyle \frac{q_0}{2\gamma \left[1+\frac{w}{\gamma }(q_1q_0)\right]}}`$ (C17) $`+{\displaystyle \frac{1}{2w}}\mathrm{ln}\left[1+{\displaystyle \frac{w}{\gamma }}(q_1q_0)\right]+{\displaystyle \frac{\alpha }{w}}{\displaystyle 𝑑h_0^\alpha G(h_0^\alpha )\mathrm{ln}𝑑h_0^\beta G(h_0^\beta \|h_0^\alpha )e^{w\mathrm{\Delta }E}}.`$ (C18) Eqs. (C1), (C2), (C10), (C11) and (C18) form a set of five equations for $`\chi `$, $`\gamma `$, $`q_1`$, $`q_0`$ and $`w`$.
warning/0006/hep-lat0006017.html	ar5iv	text	# Physical Results from Unphysical Simulations ## I Introduction A major obstacle to direct simulations of lattice QCD is the difficulty in simulating with light dynamical quarks. In particular, the up and down quarks must be reached by a chiral extrapolation. In present simulations one must do this extrapolation from roughly $`m_s/2`$, where $`m_s`$ is the physical strange quark mass. This is far from the light quark masses ($`\overline{m}=(m_u+m_d)/2m_s/25`$). The aim of this paper is to provide formulae which can aid in this extrapolation. To do this we use chiral perturbation theory (ChPT) at next-to-leading order (NLO). The parameters of the chiral Lagrangian that enter at this order are the Gasser-Leutwyler (GL) coefficients, $`L_1L_{10}`$. An alternative way to view the extrapolation to QCD is that, by fitting numerical results in a region where quark masses are considerably larger than the physical light quarks, but small enough that NLO chiral perturbation theory is sufficiently accurate, one determines the relevant $`L_i`$. These are physical parameters of QCD, governing many different physical properties (e.g. pion masses and scattering amplitudes). With the $`L_i`$ in hand, one can then extrapolate to QCD, and, in particular, determine the physical light quark masses. For example, as has been stressed in Refs. , determining the combination $`2L_8L_5`$ with only moderate accuracy might allow one to rule out the interesting possibility that $`m_u=0`$. The accuracy of extrapolation depends, of course, on the reliability of NLO chiral perturbation theory. This can be studied by seeing how well the numerical data fit the expected forms, including the curvature predicted by chiral logarithms. An observation of practical importance is that one can make use of partially quenched (PQ) simulations to aid in the extrapolation to QCD . In partially quenched simulations, one changes the mass of the valence quarks (typically reducing them), while holding the dynamical (or “sea”) quark masses fixed. The situation is illustrated schematically in Fig. 1. This leads one into a space of unphysical theories, from which one might expect to obtain only qualitative information about QCD. It turns out, however, that, if all quark masses (valence and sea) are small enough, one can use PQ theories to obtain quantitative information about unquenched theories. Since it is computationally less demanding to reduce valence quark masses, PQ simulations are often used to obtain approximate information on QCD. Our point here is that they can be used to obtain exact information about QCD. This observation follows from the structure of chiral perturbation theory (ChPT) generalized to partially quenched theories—PQChPT . The key point is that there is a subspace of quark masses (corresponding to the diagonal line in Fig. 1) where PQChPT is completely equivalent to chiral perturbation theory for unquenched, QCD-like theories. Since the quark mass dependence in PQChPT is explicit, as in ChPT, it follows that the parameters of the PQ chiral Lagrangian (with 3 light sea quarks) are the same as those of QCD. These parameters do, however, depend on the number of sea quarks, $`N`$. This means that PQ simulations with $`N=3`$ light sea quarks, whatever their precise masses, give information about the parameters of the chiral Lagrangian of QCD. On the other hand PQ simulations with $`N=2`$ or $`N=4`$, the values used in most simulations to date, do not give direct information about QCD, even after extrapolation. These comments motivate the calculation of the NLO chiral corrections to physically interesting quantities in PQ theories. Some results already exist in the literature: those for charged pion masses and decay constants with degenerate sea quarks , heavy-light meson properties , vector and tensor meson properties , baryons masses at large $`N_c`$ , and electroweak amplitudes . We provide here two new results. First, charged pion masses and decay constants are considered for non-degenerate sea quarks (having up to three different masses). This completes the calculations of simple pion properties for any theory that is likely to be simulated. It allows one to determine $`L_{46}`$ and $`L_8`$. Note that, although non-degenerate sea quarks are not necessary in order to extract these GL coefficients, as has been stressed in Refs. , there is no drawback to using them. Indeed, someone might prefer to extrapolate using more “QCD-like” simulations with two degenerate “light” dynamical quarks and one dynamical quark with mass fixed close to the physical strange quark mass. Our formulae apply for such a theory. Our second new result concerns the GL coefficient $`L_7`$. In QCD, this contributes only to neutral meson masses, and does so proportional to $`(m_s\overline{m})^2`$. To determine $`L_7`$ using meson masses from unquenched simulations thus requires non-degenerate quarks. At first sight, PQ simulations do not improve the situation: neutral meson masses are still independent of $`L_7`$ when the sea quarks are degenerate. We find, however, that one can determine $`L_7`$ from the coefficient of the double-pole in the propagators of neutral valence mesons, even for degenerate sea quarks. This is a nice example of the utility of PQ theories. Although the double-pole is itself an indicator that these theories are unphysical, its effects can nevertheless be measured in lattice simulations, and its inferred coefficient turns out to be related to a physical quantity. Throughout our calculations, we treat the $`\eta ^{}`$ as a heavy particle, and integrate it out, following Ref. . This greatly simplifies the resulting expressions, since it removes the dependence on additional $`\eta ^{}`$ coupling constants. It raises, however, two important issues. First, is the $`\eta ^{}`$ heavy enough that removing it is appropriate? The answer depends on the number of light sea quarks, $`N`$, and the number of colors, $`N_c`$. For the physical values of these parameters, $`N=N_c=3`$, we know from experiment that $`M_\eta ^{}1`$GeV. Since this is the scale at which chiral perturbation theory breaks down, it is appropriate to integrate out the $`\eta ^{}`$ for these theories. We stress, however, that our formulae are only valid when the dynamical quarks are light enough that all pseudo-Goldstone bosons, including those composed only of sea quarks, satisfy $`M_{PGB}^2M_\eta ^{}^2`$. For further discussion of this point, and of the limitations of the approach taken here, see Ref. . The second issue is more technical. How does one integrate out the $`\eta ^{}`$ in PQ theories? In this paper we follow Ref. and do this by hand, working only at tree level. This is unsatisfactory, since, as we know from QCD, $`\eta ^{}`$ loops to all orders give contributions of the same order in chiral perturbation theory. In other words, the $`\eta ^{}`$ must be integrated out non-perturbatively. We return to this issue in a companion paper, where we demonstrate that the approach adopted here is equivalent to integrating out the $`\eta ^{}`$ non-perturbatively in the PQ theory . This paper is organized as follows. In the following section we recall the formalism of PQ chiral perturbation theory, and explain our calculational procedure. After presenting a simple form for the neutral meson propagator in Sec. III, we give our results for charged pion properties in Sec. IV. We analyze these results in Sec. V, paying particular attention to the convergence of the chiral expansion and the importance of non-analytic terms. In Sec. VI we explain how to extract $`L_7`$ using PQ theories. We end with some conclusions. Two appendices deal with technical issues in the calculation of the neutral meson propagator. Some parts of this work have been reported previously in Ref. . ## II Theoretical framework We consider partially quenched theories with the following quark complement: $`N_1+N_2+N_3=N`$ sea quarks, $`N_i`$ each of mass $`m_i`$; two valence quarks with masses $`m_A`$ and $`m_B`$; and two corresponding ghost quarks with masses $`m_{\stackrel{~}{A}}=m_A`$ and $`m_{\stackrel{~}{B}}=m_B`$. The ghosts are needed to cancel the determinant arising from the valence quark functional integral . In the chiral limit, this theory has an $`SU(N+2\|2)_LSU(N+2\|2)_R`$ symmetry group . If $`N_1=N_2=N_3=1`$, and $`m_1=m_u`$, $`m_2=m_d`$ and $`m_3=m_s`$, then the sea-quark sector is QCD. Generalizing to arbitrary numbers $`N_i`$ covers most other theories that are likely to be simulated in an effort to shed light on QCD. An important property of PQ theories, which follows trivially from their definition, is that the sea-quark sector decouples from the valence sector. To be precise, all correlation functions composed of only sea-quark fields are identical to those in the unquenched sea-quark theory. There is no “back-reaction” from the valence sector. The same result must also hold for the low energy chiral Lagrangian describing the PQ theory: correlators of pseudo-Goldstone mesons composed of sea quarks should be the same as in the chiral Lagrangian describing the unquenched theory. This was shown to be true in Ref. . In practice, however, one might view all correlation functions that are calculated as being those of valence quarks, and so it is more useful to reformulate this property as follows. When each of the valence quarks is assigned a mass that is equal to one of the sea quark masses, sea and valence quarks become indistinguishable, and the cancelations of the “doubled” quark species against their ghost counterparts trivially render the theory the same as an unquenched theory containing only the “original” sea quarks. That this is so was also shown in Ref. , and it has important consequences in the following. At low energies, the partially quenched chiral effective theory is expressed in terms of the fields <sup></sup><sup></sup>The normalization of $`\mathrm{\Phi }_0`$ is different from that used in Ref. , although it agrees when $`N=3`$ as in QCD. $`\mathrm{\Sigma }`$ $`=`$ $`\mathrm{exp}(2i\mathrm{\Phi }/f),`$ (1) $`\mathrm{\Phi }`$ $`=`$ $`\left(\mathrm{\Phi }_{ab}\right)={\displaystyle \frac{1}{\sqrt{2}}}\left(\pi _{ab}\right),a,b\{A,B,1,2,\mathrm{},N,\stackrel{~}{A},\stackrel{~}{B}\},`$ (2) $`\mathrm{\Phi }_0`$ $`=`$ $`\text{str}\mathrm{\Phi }/\sqrt{N},`$ (3) and the quantity $`\chi =2\mu m=2\mu \text{ diag}(m_A,m_B,\underset{N_1}{\underset{}{m_1,\mathrm{},m_1}},\underset{N_2}{\underset{}{m_2,\mathrm{},m_2}},\underset{N_3}{\underset{}{m_3,\mathrm{},m_3}},m_A,m_B).`$ (4) where $`m`$ is the quark mass matrix. In the following we use the notation $`\chi _A=2\mu m_A`$, $`\chi _1=2\mu m_1`$, etc. The constants $`f`$ and $`\mu `$ are unknown parameters. The fields $`\pi _{ab}`$ describe the pseudo-Goldstone particles of the theory—we refer to them generically as mesons even though some are fermionic. Because of the anomaly, arbitrary functions of the field $`\mathrm{\Phi }_0`$ (the super-$`\eta ^{}`$) can appear in the Lagrangian. The partially quenched chiral effective Lagrangian is expanded in powers of $`ϵ^2M^2/\mathrm{\Lambda }^2p^2/\mathrm{\Lambda }^2`$, where $`M`$ is a typical pseudoscalar meson mass, $`p`$ the momentum, and $`\mathrm{\Lambda }1`$GeV is the scale beyond which the theory breaks down. The parts of the Euclidean Lagrangian contributing to meson masses and decay constants at one-loop order are: $`_{\mathrm{LO}}`$ $`=`$ $`{\displaystyle \frac{f^2}{4}}\text{str}\left(_\mu \mathrm{\Sigma }_\mu \mathrm{\Sigma }^{}\right){\displaystyle \frac{f^2}{4}}\text{str}\left(\chi \mathrm{\Sigma }^{}+\mathrm{\Sigma }\chi \right)+\alpha _\mu \mathrm{\Phi }_0_\mu \mathrm{\Phi }_0+m_0^2\mathrm{\Phi }_0^2`$ (5) $`_{\mathrm{NLO},1}`$ $`=`$ $`L_4\text{str}\left(_\mu \mathrm{\Sigma }_\mu \mathrm{\Sigma }^{}\right)\text{str}\left(\chi \mathrm{\Sigma }^{}+\mathrm{\Sigma }\chi \right)+L_5\text{str}\left[_\mu \mathrm{\Sigma }_\mu \mathrm{\Sigma }^{}\left(\chi \mathrm{\Sigma }^{}+\mathrm{\Sigma }\chi \right)\right]`$ (7) $`L_6\left[\text{str}\left(\chi \mathrm{\Sigma }^{}+\mathrm{\Sigma }\chi \right)\right]^2L_8\text{str}\left(\chi \mathrm{\Sigma }^{}\chi \mathrm{\Sigma }^{}+\mathrm{\Sigma }\chi \mathrm{\Sigma }\chi \right)`$ $`_{\mathrm{NLO},2}`$ $`=`$ $`L_7^{}\left[\text{str}\left(\chi \mathrm{\Sigma }^{}\mathrm{\Sigma }\chi \right)\right]^2+v_1\mathrm{\Phi }_0^2\text{str}\left(_\mu \mathrm{\Sigma }_\mu \mathrm{\Sigma }^{}\right)`$ (9) $`+v_2\mathrm{\Phi }_0^2\text{str}\left(\chi \mathrm{\Sigma }^{}+\mathrm{\Sigma }\chi \right)`$ The coefficients $`\alpha `$, $`m_0`$, $`L_i`$ and $`v_i`$ are further unknown parameters of the low energy theory.We revert to the notation $`\alpha `$ of Ref. , rather than the $`\alpha _\mathrm{\Phi }`$ used in Ref. . The $`L_i`$ depend, in general, on the renormalization scale. The NLO Lagrangian is broken into two parts because flavor off-diagonal mesons receive contributions only from $`_{\mathrm{NL0},1}`$. At this point we can make clear the relationship between the partially quenched chiral Lagrangian and that describing low energy QCD. The latter is obtained by setting $`N=3`$, and “unquenching” — i.e. assigning $`m_A`$ and $`m_B`$ values from $`\{m_1,m_2,m_3\}`$. It follows that the unknown coefficients in $``$ are, for $`N=3`$, identical to those in the QCD chiral Lagrangian. This shows that these constants also govern the chiral behavior of PQ extensions of QCD. In QCD, one can take a further step and “integrate out” the $`\eta ^{}`$. This is appropriate since it is not a pseudo-Goldstone boson, having $`M_\eta ^{}^2m_0^2+O(m)1\mathrm{GeV}^2`$. Technically, the matching between theories with and without the $`\eta ^{}`$ is non-perturbative, since loops involving the $`\eta ^{}`$ are not suppressed by powers of $`M^2`$ or $`p^2`$. Thus in the standard approach one simply writes down the Lagrangian without the $`\eta ^{}`$, and it has the same form as Eqs. (5)-(9), except that $`\mathrm{\Phi }`$ is traceless. It follows that $`\alpha `$, $`m_0`$ and the $`v_i`$ are irrelevant, and the only NLO coefficients are the $`L_i`$. It is in fact in this theory that the $`L_i`$—the Gasser-Leutwyler coefficients—are conventionally defined. In previous work on PQQCD, the step of integrating out the $`\eta ^{}`$ has been done by hand, i.e. at the level of individual diagrams rather than the Lagrangian . We summarize the procedure here—details will become apparent in the following section. Loop diagrams involving the $`\eta ^{}`$ are dropped, since these lead to shifts in the parameters $`L_i`$ which are automatically included if we use the $`L_i`$ from the QCD Lagrangian without the $`\eta ^{}`$. * Couplings special to the $`\eta ^{}`$, such as the $`v_i`$ in Eq. (9), are treated as small, of $`O(ϵ^2)`$, and thus appear only at tree level. The justification for this treatment is that these couplings are suppressed by powers of $`1/N_c`$, in this case $`1/N_c^2`$. * On the other hand, the parameter $`m_0^2`$ is treated non-perturbatively since it is known to be $`\mathrm{\Lambda }^2`$, despite the fact that it is proportional to $`1/N_c`$. In particular, we treat $`M^2/m_0^2`$ as $`O(ϵ^2)`$ (with $`M`$, as above, a typical meson mass). For convenience, we also treat $`\alpha `$ non-perturbatively. While this procedure may be accurate enough for phenomenological purposes, it is theoretically unsatisfactory because the $`\eta ^{}`$ should be integrated out non-perturbatively. As noted in the introduction, we will address this concern in a separate paper . In particular, we will argue that the procedure adopted here in fact leads to results that are equivalent to those obtained from integrating out the $`\eta ^{}`$ non-perturbatively. We close this section by deriving a result needed in Sec. VI. We claimed above that discarding $`\eta ^{}`$ loops allows us to write our results in terms of the standard GL coefficients, i.e. those in an effective Lagrangian without the $`\eta ^{}`$. This is not quite correct. Tree diagrams involving an intermediate $`\eta ^{}`$, which we keep, lead to a shift in $`L_7`$ proportional to $`M^2/m_0^2`$. Thus we are not using the conventional $`L_7`$, but rather that defined in the effective theory containing the $`\eta ^{}`$, which we denote $`L_7^{}`$ \[as anticipated in Eq. (9)\]. In order to express our results in terms of conventional parameters, we need to relate $`L_7^{}`$ to $`L_7`$ within the approximations of our procedure. To determine this relation we need consider only the sea-quark sector, and thus work with the conventional (unquenched) chiral Lagrangian. The $`\eta ^{}`$ field in this theory is the restriction to the sea sector of the “super-$`\eta ^{}`$” field $`\mathrm{\Phi }_0`$. To make the $`\eta ^{}`$ dependence explicit, we decompose $`\mathrm{\Sigma }`$ into pseudo-Goldstone and $`\eta ^{}`$ parts: $`\mathrm{\Sigma }=U\mathrm{exp}({\displaystyle \frac{2i\mathrm{\Phi }_0}{f\sqrt{N}}});USU(N),`$ (10) and substitute into the chiral Lagrangian. The result is the original form with $`\mathrm{\Sigma }U`$ and $`\mathrm{\Phi }_00`$, i.e. the standard QCD chiral Lagrangian, plus the following terms $`_{\mathrm{\Phi }_0}`$ $`=`$ $`c_1\mathrm{\Phi }_0+c_2\mathrm{\Phi }_0^2+c_3(\mathrm{\Phi }_0)^2+O(\mathrm{\Phi }_0^3)`$ (11) $`c_1`$ $`=`$ $`if\text{tr}(\chi U^{}\chi U)/(2\sqrt{N})+O(\chi ^2)`$ (12) $`c_2`$ $`=`$ $`m_0^2+O(\chi )`$ (13) $`c_3`$ $`=`$ $`1+\alpha +O(\chi )`$ (14) Only the leading order terms in the chiral expansion of the coefficients are shown, since higher order terms give contributions to the conventional chiral Lagrangian of orders $`ϵ^6`$ and higher, too high to effect the matching of the GL coefficients which appear at order $`ϵ^4`$. For the same reason, the $`(\mathrm{\Phi }_0)^2`$ term can be dropped. Keeping the $`\eta ^{}`$ in tree graphs amounts to doing the functional integral over $`\mathrm{\Phi }_0`$ keeping only linear and quadratic terms. The result is a contribution to the conventional chiral Lagrangian of order $`ϵ^4`$ with the same form as the $`L_7`$ term: $`{\displaystyle \frac{f^2}{16Nm_0^2}}[\text{tr}(\chi U\chi U^{})]^2(1+O(\chi )).`$ (15) Thus we find, within our approximations, the relation $`L_7=L_7^{}{\displaystyle \frac{f^2}{16Nm_0^2}}.`$ (16) ## III Calculation Expanding $`_{LO}`$ to quadratic order, we obtain the meson propagators. For “charged” mesons (i.e. flavor off-diagonal states) these are $`G_{ab}^C(p){\displaystyle d^4xe^{ipx}\pi _{ab}(x)\pi _{ba}(0)}={\displaystyle \frac{ϵ_b}{p^2+(\chi _a+\chi _b)/2}}(ab),`$ (17) where the signature vector is $`ϵ_a=\{\begin{array}{cc}1& a\{A,B,1,2,3,\mathrm{},N\}\hfill \\ 1& a\{\stackrel{~}{A},\stackrel{~}{B}\}\hfill \end{array}.`$ (20) The propagators for “neutral” (flavor-diagonal) mesons include the contributions of the super-$`\eta ^{}`$ interactions. A general expression has been given in Ref. , but we find it convenient to use an alternative form. The propagator is a matrix acting on the space of neutral meson fields, $`\pi _{aa}`$, $`a=1,N+4`$. In app. A we show that $`G_{ab}^N`$ $``$ $`{\displaystyle d^4xe^{ipx}\pi _{aa}(x)\pi _{bb}(0)}`$ (21) $`=`$ $`{\displaystyle \frac{ϵ_a\delta _{ab}}{p^2+\chi _a}}{\displaystyle \frac{(m_0^2+\alpha p^2)/N}{\left(p^2+\chi _a\right)\left(p^2+\chi _b\right)}}{\displaystyle \frac{\left(p^2+\chi _1\right)\left(p^2+\chi _2\right)\left(p^2+\chi _3\right)}{(1+\alpha )\left(p^2+M_{\pi _0}^2\right)\left(p^2+M_\eta ^2\right)\left(p^2+M_\eta ^{}^2\right)}}.`$ (22) Here the $`\pi _0`$, $`\eta `$ and $`\eta ^{}`$ are neutral mesons in the sea-quark sector. For $`N=3`$ they are the usual neutral mesons of QCD; for $`N>3`$ they are the appropriate generalizations, as explained in the appendix. Their masses are functions of the sea-quark masses, and of the $`N_i`$, as given explicitly in Eqs. (A25)-(A27). The neutral propagator shows explicitly the unphysical nature of the PQ theory. For example, if $`a=b=A`$, the second term has a double-pole at $`p^2=\chi _A`$. These double-poles are absent, however, in the physical, sea-quark, sector. This was shown in Ref. , but is particularly transparent with our result. For example if $`a=b=1`$ (or equivalently if $`a=b=A`$ and $`\chi _A=\chi _1`$), then the $`(p^2+\chi _1)`$ in the numerator reduces the double-pole to a single-pole. For calculations, it is preferable to rewrite the propagator as a sum of (single or double) poles. For simplicity, we discuss the case when $`ab`$ and $`\chi _a\chi _b`$, for then $`G_{ab}^N`$ has only single poles. For uniformity of notation, we introduce the definitions $`\chi _\pi =M_{\pi _0}^2,\chi _\eta =M_\eta ^2,\chi _\eta ^{}=M_\eta ^{}^2.`$ (23) in terms of which $`G_{ab}^N`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{x=a,b,\pi _0,\eta ,\eta ^{}}{}}{\displaystyle \frac{R_x}{p^2+\chi _x}}(ab),`$ (24) $`R_x`$ $`=`$ $`{\displaystyle \frac{(m_0^2\alpha \chi _x)\underset{i=1,3}{}(\chi _i\chi _x)}{(1+\alpha )_{yx}(\chi _y\chi _x)}},`$ (25) where $`x`$ and $`y`$ run over $`a,b,\pi ,\eta ,\eta ^{}`$. The degenerate limit $`\chi _a=\chi _b`$ or $`a=b`$, which, in general, has double poles is straightforward to obtain. At this point we are in a position to integrate out the $`\eta ^{}`$ by hand, bearing in mind that this propagator appears in loops. First, as discussed above, we drop the $`\eta ^{}`$ pole. Second, we expand the residues of the remaining poles in powers of $`\chi _x/m_0^2`$, $`x\eta ^{}`$, and drop all but the leading term. This is justified since we use the neutral propagator in one-loop diagrams, which already give NLO contributions. The only exception is in our discussion of $`L_7`$ in Sec. VI, where the neutral propagator appears at tree-level. With these changes, the residues become $`R_x={\displaystyle \frac{\underset{i=1,3}{}(\chi _i\chi _x)}{_{\genfrac{}{}{0pt}{}{y=a,b,\pi ,\eta }{yx}}(\chi _y\chi _x)}}.`$ (26) Note that both $`m_0^2`$ and $`\alpha `$ disappear from the neutral propagator. ## IV NLO results In this section we calculate the properties of a meson composed of two valence quarks to NLO. We use dimensional regularization, and subtract the poles following the conventions of Ref. . Its mass, $`M_{AB}`$, is obtained from the diagrams of Fig. 2. We find $`M_{AB}^2`$ $`=`$ $`{\displaystyle \frac{\chi _A+\chi _B}{2}}\left(1+\delta _{\mathrm{tree}}^M+\delta _{\mathrm{loop}}^M\right)`$ (27) $`\delta _{\mathrm{tree}}^M`$ $`=`$ $`{\displaystyle \frac{8N}{f^2}}\left(2L_6L_4\right)\overline{\chi }+{\displaystyle \frac{4}{f^2}}\left(2L_8L_5\right)(\chi _A+\chi _B)`$ (28) $`\delta _{\mathrm{loop}}^M`$ $`=`$ $`{\displaystyle \frac{1}{16f^2\pi ^2N}}\left\{R_A\chi _A\mathrm{log}\chi _A+R_B\chi _B\mathrm{log}\chi _B+R_\pi \chi _\pi \mathrm{log}\chi _\pi +R_\eta \chi _\eta \mathrm{log}\chi _\eta \right\}`$ (29) Here $`\overline{\chi }`$ is the average sea-quark mass, $`\overline{\chi }={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1,3}{}}N_i\chi _i,`$ (30) and the residues are defined in Eq. (26). For concreteness we quote two explicit examples $`R_A`$ $`=`$ $`{\displaystyle \frac{\left(\chi _A\chi _1\right)\left(\chi _A\chi _2\right)\left(\chi _A\chi _3\right)}{\left(\chi _A\chi _B\right)\left(\chi _A\chi _\pi \right)\left(\chi _A\chi _\eta \right)}},`$ (31) $`R_\pi `$ $`=`$ $`{\displaystyle \frac{\left(\chi _\pi \chi _1\right)\left(\chi _\pi \chi _2\right)\left(\chi _\pi \chi _3\right)}{\left(\chi _\pi \chi _A\right)\left(\chi _\pi \chi _B\right)\left(\chi _\pi \chi _\eta \right)}}.`$ (32) $`R_B`$ is obtained from $`R_A`$ by interchanging $`A`$ and $`B`$, and $`R_\eta `$ is obtained from $`R_\pi `$ by interchanging $`\pi `$ and $`\eta `$. The renormalization scale is implicit in the logarithms and the $`L_i`$. Using the result $`{\displaystyle \underset{x=A,B,\pi ,\eta }{}}\chi _xR_x=(\chi _A+\chi _B\overline{\chi })+O(\chi ^2/m_0^2)`$ (33) we see that a change in renormalization scale can be absorbed by shifting the $`L_i`$. We have checked (for $`N=3`$) that the scale dependence of the $`L_i`$ in QCD does render $`M_{AB}^2`$ independent of the renormalization scale. To determine the meson decay constant, $`f_{AB}`$, the axial current $`j_{5AB}^\mu `$ is calculated to NLO, and then we evaluate the matrix element $`0\|j_{5AB}^\mu \|\pi _{AB}(p)=i\sqrt{2}f_{AB}p^\mu ,`$ (34) using the diagrams of Fig. 3. The result is $`f_{AB}`$ $`=`$ $`f\left(1+\delta _{\mathrm{tree}}^f+\delta _{\mathrm{VS}\mathrm{loop}}^f+\delta _{\mathrm{VV}\mathrm{loop}}^f\right)`$ (35) $`\delta _{\mathrm{tree}}^f`$ $`=`$ $`{\displaystyle \frac{4N}{f^2}}\overline{\chi }L_4+{\displaystyle \frac{2}{f^2}}(\chi _A+\chi _B)L_5`$ (36) $`\delta _{\mathrm{VS}\mathrm{loop}}^f`$ $`=`$ $`{\displaystyle \underset{i=1,3}{}}N_i{\displaystyle \frac{1}{16\pi ^2f^2}}{\displaystyle \frac{\chi _A+\chi _i}{8}}\mathrm{log}\left({\displaystyle \frac{\chi _A+\chi _i}{2}}\right)+(AB)`$ (37) $`\delta _{\mathrm{VV}\mathrm{loop}}^f`$ $`=`$ $`{\displaystyle \frac{1}{4N}}{\displaystyle \frac{1}{16\pi ^2f^2}}\{D_AD_B`$ (40) $`+{\displaystyle \frac{\mathrm{log}(\chi _A/\chi _B)}{(\chi _A\chi _B)}}\left[\chi _AD_A+\chi _BD_B+(\chi _A\chi _B)^2\right]`$ $`+(\chi _\pi R_\pi (\chi _B\chi _A)[{\displaystyle \frac{\mathrm{log}(\chi _\pi /\chi _A)}{\chi _A\chi _\pi }}{\displaystyle \frac{\mathrm{log}(\chi _\pi /\chi _B)}{\chi _B\chi _\pi }}]+(\pi \eta ))\}`$ where $`D_A={\displaystyle \frac{\underset{i=1,3}{}(\chi _i\chi _A)}{(\chi _\pi \chi _A)(\chi _\eta \chi _A)}},`$ (41) and $`D_B`$ obtained by $`\left(AB\right)`$, are the coefficients of the double poles in the neutral propagators. It is straightforward to see that the scale dependence can be absorbed by shifts in $`L_4`$ and $`L_5`$, and we have checked that these shifts are consistent with standard results for $`N=3`$. As noted in the introduction, our formulae can be used to extract the GL coefficients by fitting to results from simulations. PQ simulations allow one to vary the valence and sea-quark masses independently, and thus to separately determine $`L_{46}`$ and $`L_8`$. In fact, since the NLO analytic terms, Eqs. (28) and (36), depend on the quark masses only through the combinations $`\chi _A+\chi _B`$ and $`\overline{\chi }`$, one need only consider degenerate sea quarks.Note that the NLO analytic dependence of $`M_{AB}^2`$ is not the most general quadratic term, symmetric under $`AB`$, and vanishing when $`\chi _A=\chi _B=0`$. Such a form would contain a term proportional to $`(\chi _A\chi _B)^2`$, which is in fact forbidden by chiral symmetry. Indeed, to extract $`2L_8L_5`$, the combination which determines whether $`m_u=0`$, it is sufficient to use a single sea quark mass (as long as it is light enough that the formulae apply), and vary the valence quark masses. By contrast, with unquenched simulations, one would have to use non-degenerate sea quarks to separately determine all four $`L`$’s. The most likely practical applications of these results are for simulations done with 2 rather than 3 types of nondegenerate sea quarks. For QCD this would correspond to the limit of exact isospin symmetry. The results for this case can be obtained from those given above by carefully taking the limit $`\chi _2\chi _1`$. When two types of quark are degenerate, one of the neutral sea-quark eigenstates becomes an exact flavor non-singlet, and we choose this to be the pion. Thus we have $`\chi _\pi =\chi _1=\chi _2`$. The remaining neutral state has mass $`\chi _\eta =\chi _1+\chi _3\overline{\chi }`$. In this limit, the residues simplify, e.g. $`R_A`$ $``$ $`{\displaystyle \frac{\left(\chi _A\chi _1\right)\left(\chi _A\chi _3\right)}{\left(\chi _A\chi _B\right)\left(\chi _A\chi _\eta \right)}}`$ (42) $`R_\pi `$ $``$ $`0`$ (43) $`R_\eta `$ $``$ $`{\displaystyle \frac{\left(\chi _\eta \chi _1\right)\left(\chi _\eta \chi _3\right)}{\left(\chi _\eta \chi _A\right)\left(\chi _\eta \chi _B\right)}}`$ (44) with a similar cancelation in the $`D_{A,B}`$. With these changes, the formulae given above still hold. We have checked our results in two ways. First, if we take all sea quarks to be degenerate, we obtain the results of Ref. .<sup>§</sup><sup>§</sup>§Except for the following typos in Ref. , pointed out by Jochen Heitger, Rainer Sommer and Hartmut Wittig: in Eqs. (18), (19) and (20) $`\alpha _4`$ should be replaced by $`\alpha _4/2`$. Second, we can consider the unquenched limit. By choosing $`\chi _A`$ and $`\chi _B`$ to be equal to combinations of $`\chi _{13}`$ we obtain the correct one-loop form for the masses of the $`\pi ^+`$, $`K^+`$ and $`K^0`$. The residues $`R_x`$ and $`D_x`$ are singular when pairs of the $`\chi `$’s become degenerate, e.g. $`\chi _A\chi _B`$ and $`\chi _A\chi _\pi `$. As expected, however, these singularities cancel in the full expressions for $`M_{AB}`$ and $`f_{AB}`$, which are analytic functions of the quark masses except in the massless limit. In Ref. , it was emphasized that the one-loop corrections can diverge when the valence-quark masses are sent to zero at fixed sea-quark mass, leading to a breakdown of chiral perturbation theory. This discussion was based on the results for degenerate sea quarks, and we can now generalize it to non-degenerate sea-quarks. For the meson masses, the possibly divergent contribution is $`\delta _{\mathrm{loop}}^M`$, and we see that this only diverges if both $`m_A`$ and $`m_B`$ vanish, in fixed ratio, but not if only one vanishes. For the decay constant the pattern is opposite: $`\delta _{\mathrm{VV}\mathrm{loop}}^f`$ diverges if one of the valence masses vanishes, but not if both do in fixed ratio. This is the same pattern of divergences as for degenerate sea quarks; the sea quark masses only influence the coefficients of the divergent $`\mathrm{log}\chi _A`$ and $`\mathrm{log}\chi _B`$ terms. In the following section we discuss the practical implications of these divergences. It was also noted in Ref. that one can form combinations of the squared meson masses and decay constants from which the analytic correction terms ($`\delta _{\mathrm{tree}}`$) cancel. One thus predicts these combinations in terms of the quark masses and the leading order chiral coefficients $`\mu `$ and $`f`$, up to NNLO corrections. Studying these combinations in simulations allows one to test the applicability of NLO chiral perturbation theory. What we want to point out here is that exactly the same quantities can be used for non-degenerate sea-quarks: the $`L_i`$ still cancel. We do not, however, give the explicit expressions since they are lengthy and unilluminating. Finally, we discuss to what extent our formulae can be applied to lattice results obtained at non-zero lattice spacing. For definiteness, we first consider a calculation using Wilson fermions, in which we work at a fixed bare coupling and vary the bare valence and sea quark masses. Meson properties are of course calculated in lattice units, i.e. we obtain $`M_\pi a`$ and $`f_\pi a`$. There are two effects of working at a non-zero lattice spacing. The first is that chiral symmetry, upon which our calculation is based, is broken explicitly. This symmetry breaking can, however, be incorporated into the chiral Lagrangian framework, as shown in Ref. . The result is that all effects of the explicit symmetry breaking are of $`O(a)`$, except for the additive renormalization of the quark masses. Because of this, our formulae are valid, up to corrections of $`O(a)`$, as long as one uses so-called vector or axial Ward identity quark masses. Note that the corrections of $`O(a)`$ cannot be incorporated into our formulae by simply introducing $`a`$ dependence into the parameters of the chiral Lagrangian parameters, $`f`$, $`\mu `$ and the $`L_i`$. There are additional unknown constants which enter. The second effect is that the lattice spacing itself depends on the quark masses, at fixed bare coupling. This introduces an additional mass dependence into quantities expressed in lattice units. We note, however, that the mass dependence of $`a`$ is a discretization effect of $`O(am)`$ induced by the explicit chiral symmetry breaking . Indeed, for non-perturbatively on-shell-improved Wilson fermions one can calculate this dependence, with the result $$a(m)=a(m=0)[1ca\overline{m}+\mathrm{}],c=\frac{gb_g(g)}{2\beta _{\mathrm{LAT}}(g)},$$ (45) where $`\overline{m}`$ is the average sea-quark mass, $`b_g(g)g^2+O(g^4)`$ is an improvement coefficient introduced in Ref. , and $`\beta _{\mathrm{LAT}}=dg/d\mathrm{ln}ag^3+O(g^5)`$ is the lattice $`\beta `$-function. Note that this $`O(a)`$ effect can be absorbed by shifting the GL coefficients as follows: $$L_{4,6}L_{4,6}\frac{caf^2}{8N\mu }.$$ (46) Alternatively, one could adjust the bare coupling as the quark masses are varied so as to keep the lattice spacing fixed. In summary, our formulae are approximately valid for meson properties expressed in lattice units, with the errors being of $`O(a)`$. Some, but not all, of these discretization errors can be absorbed into $`a`$ dependence of the parameters of the chiral Lagrangian. With staggered or overlap fermions the errors would instead be of $`O(a^2)`$. Since discretization errors can still be substantial at present lattice spacings, it may be better to extrapolate first to $`a=0`$, and then fit to the predicted forms. ## V Behavior of the Chiral Expansion In the framework of chiral perturbation theory (regardless of quenching) one assumes that for any quantity calculated to a given order in the chiral expansion, higher order terms are smaller. Once the unknown couplings and parameters of the theory are determined (e.g. by experiment, or by fitting to lattice data) the consistency of the computation can be checked numerically. This is the main purpose of the current section. We can make this check because, using experimental data, we have reasonable estimates for the actual values of some of the Gasser-Leutwyler coefficients. Thus we can use our results to predict the meson masses and decay constants in PQ simulations. Of course, these predictions are approximate because we only know the $`L_i`$ approximately—this is where the lattice results themselves come in—but they give us a reasonable idea of how the chiral expansion behaves. We consider QCD with exact isospin symmetry, i.e. $`N_1=N_2=N_3=1`$ and $`m_1=m_2=m_u`$ and $`m_3=m_s`$. The meson masses and decay constants then depend on the seven parameters $`f`$, $`\chi _u`$, $`\chi _s`$, $`L_{46}`$ and $`L_8`$. We want to choose these parameters so that the four charged meson quantities $`M_\pi `$, $`M_K`$, $`f_\pi `$ and $`f_K`$ take their experimental values. In order to match the number of parameters and observables, we take as starting values the GL coefficients quoted in Ref. . These are based on experimental results ($`L_5(M_\eta )\mathrm{2.3\; 10}^3,L_8(M_\eta )\mathrm{1.2\; 10}^3`$) and the large $`N_c`$ limit ($`L_4(M_\eta )L_6(M_\eta )0`$). We then take our four free parameters to be $`f`$, $`\chi _u`$, $`\chi _s`$ and the scale, $`\mathrm{\Lambda }_L`$, at which $`L_i`$ take the values just quoted. In effect, this moves us through the space of $`L_i`$ on a particular path, which, when $`\mathrm{\Lambda }_LM_\eta `$, is consistent with our knowledge about the $`L_i`$. We claim no fundamental basis for this path—we use it for simplicity. It allows us to determine the dimensionless quantities $`f/\mathrm{\Lambda }_L`$, $`\chi _u/f^2`$ and $`\chi _s/f^2`$ by fitting the ratios $`M_\pi /f_K,M_K/f_K`$ and $`f_K/f_\pi `$, and then to determine $`\mathrm{\Lambda }_L`$ by requiring, say, $`M_\pi =140`$MeV. We find $`\mathrm{\Lambda }_L=M_\eta (1.0021)`$ (using $`M_\eta =547`$MeV), so that the fitted $`L_i`$’s are, in fact, close to the inputs. The other outputs are $`f=85`$MeV, $`\chi _u/\chi _s=0.044`$ and $`\chi _s=(673\text{MeV})^2`$. We stress that we are not claiming that we have found a unique set of parameters. There is a region in the space of the $`L_i`$ which can describe the experimental observables, and for which it turns out that the chiral expansion is under reasonable control. We have picked, somewhat arbitrarily, one point in this region. We can now explore the behavior of the chiral expansion as a function of the four quantities $`({\displaystyle \frac{\chi _A}{\chi _s}},{\displaystyle \frac{\chi _B}{\chi _s}},{\displaystyle \frac{\chi _1}{\chi _s}},{\displaystyle \frac{\chi _3}{\chi _s}}).`$ (47) We have chosen to normalize the various quark masses relative to the physical strange quark mass, so that a ratio of unity represents the outer limit of where one would expect chiral perturbation theory to be reliable. We consider two types of two-dimensional cross-section of the parameter space: $`(y,y,x,1)`$ and $`(y,1,x,1)`$. In both cases $`y`$ corresponds to a valence mass while $`x`$ is proportional to a sea quark mass. We name these cross-sections the “$`\pi `$-plane” and the “$`K`$-plane” because the unquenched line $`y=x`$ in the former describes a pion-like meson made up of two identical light quarks, whereas, in the latter, $`y=x`$ corresponds to a kaon-like meson. We examine the relative size of the NLO contributions by plotting $`(M_{AB}^2)_{\text{NLO}}/(M_{AB}^2)_{\text{LO}}=(\delta _{\mathrm{tree}}^M+\delta _{\mathrm{loop}}^M)\delta ^M`$ (48) and $`(f_{AB})_{\text{NLO}}/(f_{AB})_{\text{LO}}=\delta _{\mathrm{tree}}^f+\delta _{\mathrm{VS}\mathrm{loop}}^f+\delta _{\mathrm{VV}\mathrm{loop}}^f\delta ^f`$ (49) \[see Eqs. (27), (35)\]. In each plane, these functions are plotted along rays emerging from the origin at angles $`15^{},30^{},45^{},60^{}`$ and $`75^{}`$ with respect to the $`x`$-axis, confined to the unit square (Fig. 4). The $`45^{}`$ line corresponds to unquenched theories, the other lines to PQ theories. The plots are shown in Figs. 5-8. We also show a contour plot of $`\delta ^f`$ in the K-plane in Fig. 9. The first conclusion that can be drawn from these figures is that chiral perturbation theory for QCD to one-loop order is reasonably convergent throughout the PQ region defined by $`\chi _x<\chi _s`$, $`x=A,B,1,3`$. With our parameters, the least convergent quantity is $`f_K`$. We also note that, generally speaking, the expansion is better behaved when one increases the masses along the lines at small angles. In other words, the expansion may be more reliable when valence masses are smaller than sea-quark masses, and vica-versa. This is good news for simulations, since pushing to small valence quark masses is relatively cheap. Of course, as noted in the previous section, the valence quark masses cannot be pushed too low at fixed sea masses. At some point the corrections $`\delta ^M`$ (in the $`\pi `$-plane) and $`\delta ^f`$ (in the K-plane) diverge. This is not apparent, however, from Fig. 9. To observe the divergence we show in Fig. 10 the region close to the $`x`$-axis, which does reveal the expected features. Though of theoretical interest, it seems that the breakdown of the chiral expansion due to enhanced chiral logarithms occurs only in a tiny region of parameter space and is therefore of little or no practical significance. Our final comment concerns the importance of including the non-analytic terms in fits to PQ data. The tree level contributions alone would have produced only straight lines in Figs. 5-8. The curvature seen in the graphs is due to the logarithms originating in loop diagrams. An attempt to model data collected in the heavy quark mass region with only the tree level terms will clearly lead to a significant systematic errors in the extrapolations to QCD values. ## VI Determining $`L_7`$ As noted above, the GL coefficients $`L_4`$, $`L_5`$, $`L_6`$ and $`L_8`$ can be obtained using PQ simulations with degenerate sea quarks. Simulations with non-degenerate sea quarks are useful but not essential. In this section we consider the other coefficient which enters into the NLO expressions for the physical meson masses, namely $`L_7`$. This contributes only to flavor-diagonal mesons, i.e. the $`\eta `$ and $`\pi _0`$ in QCD. Since the contribution is proportional to quark mass differences, the largest effect is on the $`\eta `$ mass. This is conveniently isolated using the violation of the mesonic Gell Mann–Okubo relation $`4m_K^2m_\pi ^23m_\eta ^2`$ $`=`$ $`{\displaystyle \frac{8}{f^2}}{\displaystyle \frac{4(m_K^2m_\pi ^2)^2}{3}}(L_56L_812L_7)`$ (51) $`+{\displaystyle \frac{\left[3m_\eta ^4\mathrm{log}(m_\eta ^2)m_\pi ^4\mathrm{log}(m_\pi ^2)4m_K^4\mathrm{log}(m_K^2)\right]}{16\pi ^2f^2}}.`$ Here we have given the result in the isospin-symmetric limit, i.e. with two degenerate quarks of mass $`\overline{m}=(m_u+m_d)/2`$, and a single strange quark.Note that the one-loop expressions for the dependence of $`m_K`$ and $`m_\pi `$ on quark masses can be obtained from the general results above by setting $`N_1=N_2=N_3=1`$, $`m_2=m_1`$, and choosing the valence masses to be equal to the appropriate sea quark masses. Thus one method for determining $`L_7`$ is to calculate $`m_\eta `$, $`m_K`$ and $`m_\pi `$ in unquenched simulations with non-degenerate sea quarks and fit to the result in Eq. (51). In this section we show how one can obtain $`L_7`$ using degenerate sea quarks by taking advantage of PQ simulations. A clue on how to proceed is provided by the fact that the $`\eta `$ propagator contains contractions in which the quark propagators are disconnected (“hairpin” contractions). Close to the degenerate limit, these contractions are proportional to $`(m_s\overline{m})^2`$, due to cancelations between light and strange quark propagators. Comparing to Eq. (51), and noting that $`m_K^2m_\pi ^2m_s\overline{m}`$, it is plausible that the $`L_7`$ contribution is related to the disconnected contraction, and that by studying this contraction in the PQ theory one can determine $`L_7`$. This is indeed what we find. First, we recall the general form of the tree level propagators. In section III we showed that the propagators for flavor off-diagonal mesons, $`G_{ab}^C`$, have an ordinary single pole at the meson mass \[Eq. (17)\], while $`G_{ab}^N`$, the “neutral” or flavor-diagonal propagators, have a double pole at the same mass if $`a=b`$ or $`m_a=m_b`$ \[Eq. (22)\]. As we show below, these general forms prevail also at NLO: $`G_{AB}^C(p)\|_{m_A=m_B}`$ $`=`$ $`{\displaystyle \frac{𝒵_𝒜}{p^2+M_{AA}^2}}+\text{non-pole},`$ (52) $`G_{AB}^N(p)\|_{m_A=m_B}`$ $`=`$ $`{\displaystyle \frac{𝒵_𝒜𝒟}{(p^2+M_{AA}^2)^2}}+\text{single-pole}+\text{non-pole},`$ (53) $`G_{AA}^N(p)`$ $`=`$ $`G_{AB}^N(p)\|_{m_A=m_B}+G_{AB}^C\|_{m_A=m_B}`$ (54) $`=`$ $`{\displaystyle \frac{𝒵_𝒜𝒟}{(p^2+M_{AA}^2)^2}}+\text{single-pole}+\text{non-pole}.`$ (55) Here $`M_{AA}`$ is a shorthand for $`M_{AB}(m_A=m_B)`$, and its NLO expression is given in Eq. (27). Note that these equations are valid only if $`A`$ and $`B`$ are partially quenched rather than unquenched (i.e. $`m_Am_i`$), and apply only close to the poles at $`p^2=M_{AA}^2`$. In particular, the neutral propagators have additional poles, at the masses of the neutral sea-quark mesons. These are not of interest here. The quantity we propose to use to determine $`L_7`$ is $`𝒟`$, the (suitably normalized) coefficient of the double pole in the neutral propagators. An equivalent definition in terms of lattice observables is $`{\displaystyle \frac{d^3x\pi _{AA}(t,\stackrel{}{x})\pi _{BB}(0)}{d^3x\pi _{AB}(t,\stackrel{}{x})\pi _{BA}(0)}}\|_{m_A=m_B}{\displaystyle \genfrac{}{}{0pt}{}{}{\mathrm{t}\mathrm{}}}{\displaystyle \frac{𝒟t}{2M_{AA}}}.`$ (56) Here the flavor indices are chosen to select the “disconnected” (numerator) and “connected” (denominator) contractions contributing to the $`\pi _{AA}`$ propagator. This choice simplifies the numerical calculation, but is not necessary. It follows from Eq. (55) that the numerator could be replaced by $`\pi _{AA}(t,\stackrel{}{p}=0)\pi _{AA}(0)`$, i.e. the temporal Fourier transform of $`G_{AA}^N`$. This changes the $`t`$independent part of the ratio but leaves $`𝒟`$ unchanged. At large times, the denominator is dominated by the single-pole contribution of the pseudo-Goldstone boson of mass $`M_{AA}`$. Only double-pole contributions to the numerator lead to the ratio growing linearly with $`t`$, and $`𝒟`$ measures their size. We note also that our ratio is the standard one used in studies of artifacts in the quenched theory. We claim that $`𝒟`$ is a “physical” quantity in the PQ theory, on a par with the valence meson masses such as $`M_{AA}`$. In support of this claim, we note that $`𝒟`$ is determined from the long distance properties of correlators, and is independent of the choice of interpolating fields (because it is determined by a ratio). In particular, we claim that $`𝒟`$ can be calculated using the effective low-energy theory with the same level of reliability as the meson masses.<sup>\</sup><sup>\</sup>\**Proving this claim rigorously seems difficult, given the lack of a Hamiltonian formulation of the PQ theory. Our point, however, is that the calculation of $`𝒟`$ is as well controlled as that of meson masses. We stress that the coefficients of the single poles, unlike the double pole, do depend on the choice of interpolating fields and are not quantities which can be predicted using chiral perturbation theory. In PQChPT, we obtain $`𝒟`$ by calculating the propagators $`G_{ab}^C`$ and $`G_{ab}^N`$ at NLO and using Eqs. (52)-(55). The tree level result for $`𝒟`$ can be read off from Eqs. (17) and (22) $`𝒟={\displaystyle \frac{1}{N}}{\displaystyle \frac{(\chi _1\chi _A)(\chi _2\chi _A)(\chi _3\chi _A)}{(\chi _\pi \chi _A)(\chi _\eta \chi _A)}}`$ (57) Note that the residue of the double-pole vanishes whenever the valence quark mass equals any of the sea quark masses. This must be the case since one is then considering a correlator which could be constructed entirely from sea quarks, and thus is physical, and cannot contain double poles. The dependence of $`𝒟`$ on $`L_7`$ begins at one-loop order. We have calculated $`𝒟`$ to this order only for the case of degenerate sea quarks ($`m_1=m_2=m_3`$). Details are given in app. B. The result is $`𝒟`$ $`=`$ $`{\displaystyle \frac{1}{N}}(M_{SS}^2M_{AA}^2){\displaystyle \frac{16}{f^2}}\left(L_7^{}{\displaystyle \frac{f^2}{16Nm_0^2}}+{\displaystyle \frac{L_5}{2N}}\right)(M_{SS}^2M_{AA}^2)^2`$ (59) $`+{\displaystyle \frac{1}{16\pi ^2f^2}}\left({\displaystyle \frac{1}{2}}[\chi _S\chi _A]^2\mathrm{log}\chi _{SA}+\chi _A^2\mathrm{log}(\chi _A/\chi _{SA})+\chi _S^2\mathrm{log}(\chi _S/\chi _{SA})\right),`$ where we use $`S`$ to denote the sea-quark, and $`\chi _{SA}=(\chi _S+\chi _A)/2`$. The first term is the same as the tree-level result with one-loop corrected meson mass-squareds replacing quark masses. The second term is the analytic term containing the $`L_7^{}`$ dependence. As advertised, $`L_7^{}`$ and $`m_0^2`$ appear in the appropriate combination to be combined into the standard $`L_7`$ \[Eq. (16)\]. The logarithmic terms, from wavefunction renormalization and from loop diagrams, combine into a fairly simple form. One check on the result is that the anomalous dimension of $`L_5`$ is such that it cancels the dependence on the choice of the scale in the logarithm (for $`N=3`$, where the anomalous dimension is known). Another is that it vanishes when $`m_1=m_A`$. Thus, from the coefficient of the double pole, one can extract the combination $`2NL_7+L_5`$. Combined with the results of the previous sections this allows a determination of $`L_7`$. ## VII Conclusions Partially quenched theories can play an important role in determining physical parameters of QCD. Our results show how one can use them to simplify the extrapolation to QCD and the determination of the GL parameters $`L_{48}`$. In particular, it is sufficient, though not necessary, to use degenerate sea quarks. One must, however, use three sea quarks. Our approach relies on chiral perturbation theory at next-to-leading order. An important issue is how light the sea quarks need to be for our formulae to be sufficiently accurate. A conservative approach is to work down to masses where the NLO corrections themselves are 10% of the leading order result, so that the missing NNLO terms are very small . This requires working down to $`m_{\mathrm{sea}}m_s/8`$. Another approach is to look at our figures and see what mass range is required to observe the predicted curvature. To do so appears to require working down to at least $`m_{\mathrm{sea}}m_s/4`$. In the end, this issue can be resolved using simulations themselves, including partially quenched simulations, to check the reliability of the NLO predictions. This exercise will also shed light on the question of whether the physical strange quark mass is light enough that NLO chiral perturbation theory is applicable for QCD. Even if the strange quark turns out to be too heavy, the results presented here still apply to the light-quark sector, if we set $`N=2`$. It is often observed that linear fits are adequate for the mass dependence of physical quantities in the range $`m_{\mathrm{sea}}=m_s/2m_s`$. Our results make clear, however, that extending such linear fits to the chiral limit, i.e. leaving out the curvature due to chiral logarithms, can lead to a significant error in the determination of physical parameters. This is most clearly illustrated by Fig. 6. All the formulae we present are for infinite volume, and thus apply for box sizes such that $`M_\pi L1`$. As pion masses decrease, this requires working in boxes of increasing size. This should not, however, be necessary. Chiral perturbation theory remains valid for finite $`M_\pi L`$, as long as $`\mathrm{\Lambda }_{\mathrm{QCD}}L1`$, and can be used to calculate the volume dependence of physical quantities. It would be interesting to do this for the quantities we consider here. Another interesting extension of our calculations is to include the effect of discretization errors within the chiral Lagrangian itself, i.e. to use the appropriate Lagrangian for the lattice theory at non-zero lattice spacing. This would aid in the extrapolation to the continuum limit. Finally, we comment on the theoretical status of our methods. We use a Lagrangian containing the $`\eta ^{}`$ (or more precisely the “super-$`\eta ^{}`$”, $`\mathrm{\Phi }_0`$), and our calculations rely on assumptions about the size of its couplings. One can show, however, that there is a limit in which our calculation is equivalent to that with the $`\eta ^{}`$ integrated out non-perturbatively . In that limit (essentially $`m_0^2\mathrm{}`$) the assumptions that we have made are valid. One might also be concerned about the theoretical foundation of the whole calculation: Is it justified to use chiral perturbation theory for the unphysical PQ theory? We have also made some progress on this issue. One can show that, in certain cases, derivatives of partially quenched quantities with respect to valence quark masses can be exactly related to derivatives of unquenched quantities with respect to sea quark masses . Our NLO results are consistent with these exact relations. ## Acknowledgments We thank David Kaplan and Ann Nelson for useful conversations. This work was supported in part by U.S. Department of Energy Grant No. DE-FG03-96ER40956/A006. ## A Neutral meson propagator at tree level Here we calculate the neutral meson propagator. From $`_{\mathrm{LO}}`$ we find $`G_N^1`$ $`=`$ $`G_0^1+V`$ (A1) $`\left(G_0^1\right)_{ab}`$ $`=`$ $`(p^2+\chi _a)\delta _{ab}ϵ_a`$ (A2) $`V_{ab}`$ $`=`$ $`{\displaystyle \frac{m_0^2+\alpha p^2}{N}}ϵ_aϵ_b.`$ (A3) The full propagator is thus $`G_N=\left(G_0^1+V\right)^1=(1+G_0V)^1G_0.`$ (A4) Because $`V`$ is an outer product, the combination $`G_0V`$ is proportional to a projection operator $`A`$ $``$ $`{\displaystyle \frac{G_0V}{\text{tr}(G_0V)}},A^2=A.`$ (A5) Thus for any function $`f`$, $`f(A)f(0)=A[f(1)f(0)],`$ (A6) and so $`\left(1+G_0V\right)^1`$ $`=`$ $`\left(1+\text{tr}(G_0V)A\right)^1`$ (A7) $`=`$ $`1+A\left[{\displaystyle \frac{1}{1+\text{tr}(G_0V)}}1\right]`$ (A8) $`=`$ $`1{\displaystyle \frac{G_0V}{1+\text{tr}(G_0V)}}.`$ (A9) Inserting this in Eq. (A4), we find $`G_N=G_0{\displaystyle \frac{G_0VG_0}{1+\text{tr}(G_0V)}}.`$ (A10) which reproduces the result of Ref. . The analytic structure of the propagator is not clear from this result. In particular, its diagonal elements $`[G_N]_{aa}`$ appear to contain double poles (from the two factors of $`G_0`$ in the second term). We know, however, that if we restrict ourselves to the physical sea-quark sector then $`G_N`$ cannot contain double poles. Thus there are cancelations hidden in Eq. (A10) which we want to make explicit. To do so we need to introduce the restriction of the various matrices to the sea-sea sector. We denote these restrictions by overbars. We first observe that the previous steps go through identically for the restricted matrices $`\left(\overline{G_N^1}\right)^1=\overline{G}_0{\displaystyle \frac{\overline{G}_0\overline{V}\overline{G}_0}{1+\text{tr}(\overline{G}_0\overline{V})}};`$ (A11) Comparing this with the restriction of Eq. (A10), $`\overline{G}_N=\overline{G}_0{\displaystyle \frac{\overline{G}_0\overline{V}\overline{G}_0}{1+\text{tr}(G_0V)}};`$ (A12) and using the result $`\text{tr}(G_0V)=\text{tr}(\overline{G}_0\overline{V}),`$ (A13) (which follows because the valence and ghost contributions cancel) we find $`\overline{G}_N=\left(\overline{G_N^1}\right)^1.`$ (A14) In other words, restriction to the sea-quark sector commutes with inversion. This result is non-trivial because $`G_N^1`$ is not block diagonal—it encapsulates the lack of feedback from the valence to the sea-quark sector. Returning to the simplification of the propagator, we note that $`det\left[\overline{G}_N^1\right]/det\left[\overline{G}_0^1\right]`$ $`=`$ $`det\left[\overline{G}_0\overline{G}_N^1\right]`$ (A15) $`=`$ $`det\left[1+\text{tr}(\overline{G}_0\overline{V})\overline{A}\right]`$ (A16) $`=`$ $`\mathrm{exp}\text{tr}\mathrm{ln}\left[1+\text{tr}(\overline{G}_0\overline{V})\overline{A}\right]`$ (A17) $`=`$ $`\mathrm{exp}\mathrm{ln}\left[1+\text{tr}(\overline{G}_0\overline{V})\right]`$ (A18) $`=`$ $`1+\text{tr}(\overline{G}_0\overline{V}).`$ (A19) Thus the factor multiplying the double pole in Eq. (A10) can be rewritten as $`\left[1+\text{tr}(G_0V)\right]^1=det\left[\overline{G}_0^1\right]/det\left[\overline{G}_N^1\right].`$ (A20) The determinant in the numerator is simple $`det\left[\overline{G}_0^1\right]=(p^2+\chi _1)^{N_1}(p^2+\chi _2)^{N_2}(p^2+\chi _3)^{N_3}.`$ (A21) Our final task is to evaluate the determinant in the denominator. To do this, we note that $`\overline{G}_N^1=\overline{G_N^1}`$ is block diagonal. The exact $`SU(N_1)\times SU(N_2)\times SU(N_3)`$ flavor symmetry implies that there are, for each sea-quark type $`i`$, $`N_i1`$ flavor non-singlet neutral pions which are eigenvectors of $`\overline{G}_N^1`$. Since $`\overline{V}`$ projects onto flavor singlet states, the corresponding eigenvalues are those of $`\overline{G}_0^1`$, namely $`(p^2+\chi _i)`$. The non-trivial part of $`\overline{G}_N^1`$, to which $`V`$ does contribute, is thus a $`3\times 3`$ block. In QCD, with $`N_i=1`$, this is the entire matrix, and describes the $`\pi _0`$, $`\eta `$ and $`\eta ^{}`$. For convenience, we use these names for general $`N_i`$ as well. We denote the restriction of matrices to this subspace with double bars. A straightforward exercise shows that $`\overline{\overline{G_N^1}}`$ $`=`$ $`SR^TDRS,`$ (A22) $`S`$ $`=`$ $`\text{diag}(1,1,\sqrt{1+\alpha }),`$ (A23) $`D`$ $`=`$ $`\text{diag}(p^2+\chi _\pi ,p^2+\chi _\eta ,p^2+\chi _\eta ^{}),`$ (A24) with $`R`$ a rotation matrix. $`S`$ rescales the singlet field so that it has a canonical kinetic term. The meson mass-squareds are given, up to corrections of size $`\chi ^2/m_0^2`$, by $`\chi _\pi +\chi _\eta `$ $`=`$ $`\chi _1+\chi _2+\chi _3\overline{\chi }`$ (A25) $`\chi _\pi \chi _\eta `$ $`=`$ $`\chi _1\chi _2\chi _3\overline{\chi ^1}`$ (A26) $`\chi _\eta ^{}`$ $`=`$ $`(m_0^2+\overline{\chi })/(1+\alpha )`$ (A27) where $`\overline{\chi }={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1,3}{}}N_i\chi _i,\overline{\chi ^1}={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1,3}{}}N_i\chi _i^1,`$ (A28) are averages over the sea sector. Using these results we find that $`det\left[\overline{G}_N^1\right]=(p^2+\chi _1)^{N_11}(p^2+\chi _2)^{N_21}(p^2+\chi _3)^{N_31}(p^2+\chi _\pi )(p^2+\chi _\eta )(p^2+\chi _\eta ^{})(1+\alpha ).`$ (A29) Inserting this and Eq. (A21) into Eqs. (A20) and (A10) gives the result quoted in the main text, Eq. (22). It is straightforward to generalize this result to an arbitrary number of different sea-quark masses. ## B One-loop calculation of $`𝒟`$ To extract $`𝒟`$ using Eqs. (55) and (52) we need the NLO results for the charged and neutral propagators. We do the calculation only for $`N`$ degenerate sea quarks, which we denote using the label $`S`$ rather than $`1,2,3`$. The NLO calculation of the charged propagator was described in Sec. IV, and the expression for $`M_{AA}^2`$ can be obtained using Eq. (27). Here we also need the wavefunction renormalization factor, for which we find $`𝒵_A=\left[1{\displaystyle \frac{8}{f^2}}(L_4N\chi _S+L_5\chi _A){\displaystyle \frac{N}{3}}{\displaystyle \frac{1}{16\pi ^2f^2}}{\displaystyle \frac{\chi _S+\chi _A}{2}}\mathrm{log}\left({\displaystyle \frac{\chi _S+\chi _A}{2}}\right)\right].`$ (B1) For the neutral propagator we need to generalize the calculation of app. A to NLO. The inverse propagator becomes $`G_N^1=G_0^1+V+\mathrm{\Sigma },`$ (B2) where $`\mathrm{\Sigma }`$ contains the NLO tree and one-loop contributions. The structure of $`\mathrm{\Sigma }`$ is such that the method used to calculate $`G^1`$ in Sec. II does not apply, and we simply invert the matrix by brute force. For our purposes, it is sufficient to consider the restriction of $`G_N`$ to the three-dimensional basis $`(\pi _{AA},\eta ^{},\pi _{\stackrel{~}{A}\stackrel{~}{A}}),\eta ^{}={\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{i=1}{\overset{N}{}}}\pi _{ii}.`$ (B3) This is because, first, we want only the $`\pi _{AA}`$ propagator and so do not need to introduce an additional valence quark; and, second, because we use degenerate sea quarks so that there is no mixing of $`\pi _{AA}`$ with flavor non-singlet neutral sea-quark mesons. The generalization to non-degenerate sea quarks, which involves a larger basis of neutral states, is straightforward in principle, but tedious in practice, and we have not carried it out. The contributions to $`\mathrm{\Sigma }`$ fall into two classes: those that are common to the charged mesons $`\mathrm{\Sigma }_C`$, and those that are special to the neutral mesons $`\mathrm{\Sigma }_N`$. The former can be obtained from the results of Sec. IV, and we find $`G_0^1+\mathrm{\Sigma }_C`$ $`=`$ $`\text{diag}(v,w,v),`$ (B4) $`v`$ $`=`$ $`(p^2+M_{AA}^2)/𝒵_A,`$ (B5) $`w`$ $`=`$ $`(p^2+M_{SS}^2)/𝒵_S.`$ (B6) Here $`M_{SS}^2`$ is the squared mass of flavor non-singlet mesons composed of sea quarks evaluated at NLO, and $`𝒵_S`$ the corresponding wavefunction renormalization. Their expressions can be obtained from those for $`M_{AA}^2`$ and $`𝒵_A`$ by the substitution $`\chi _A\chi _S`$. We now turn to $`\mathrm{\Sigma }_N`$. Because of the graded symmetry, it shares a particular matrix structure with $`V`$, and it is convenient to make this explicit by the following definitions: $`V+\mathrm{\Sigma }_N=\left(\begin{array}{ccc}x& \sqrt{N}y& x\\ \sqrt{N}y& Nz& \sqrt{N}y\\ x& \sqrt{N}y& x\end{array}\right).`$ (B10) The contributions to this matrix from $`V`$ are $`x_V=y_V=z_V=x_0=(m_0^2+\alpha p^2)/N.`$ (B11) Thus we write $`x=x_0+\delta x`$, $`y=x_0+\delta y`$ and $`z=x_0+\delta z`$. The diagrams contributing to $`\mathrm{\Sigma }_N`$ are of the general form shown in Fig. 11. The tree level contributions come from the two-meson vertices included in $`_{NLO,2}`$ \[Eq. (9)\]. The $`v_2`$ term acts as a subleading correction to $`m_0^2`$. Since at LO \[Eq. (57)\], $`𝒟`$ is independent of $`m_0^2`$, it follows that at NLO $`𝒟`$ can at most depend on its leading value. Thus $`v_2`$ contributes only at NNLO. The other tree level diagram comes from the $`L_7`$ term which gives the following contributions: $`\delta ^{(1)}x={\displaystyle \frac{8}{f^2}}2L_7\chi _A^2,\delta ^{(1)}y={\displaystyle \frac{8}{f^2}}2L_7\chi _A\chi _S,\delta ^{(1)}z={\displaystyle \frac{8}{f^2}}2L_7\chi _S^2.`$ (B12) Figure 12 shows the quark line structure of the only one-loop graph involving the lowest order vertices which corresponds to disconnected quark lines, and so does not contribute to $`\mathrm{\Sigma }_C`$. The contribution to $`\mathrm{\Sigma }_N`$ is $`\delta ^{(2)}x`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2f^2}}{\displaystyle \frac{1}{3}}(p^22\chi _A)\chi _A\mathrm{log}\chi _A,`$ (B13) $`\delta ^{(2)}y`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2f^2}}{\displaystyle \frac{1}{3}}(p^2\chi _A\chi _S){\displaystyle \frac{\chi _A+\chi _S}{2}}\mathrm{log}\left({\displaystyle \frac{\chi _A+\chi _S}{2}}\right),`$ (B14) $`\delta ^{(2)}z`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2f^2}}{\displaystyle \frac{1}{3}}(p^22\chi _S)\chi _S\mathrm{log}\chi _S.`$ (B15) Collecting these contributions we end up with $`G^1=\left(\begin{array}{ccc}v+x& \sqrt{N}y& x\\ \sqrt{N}y& w+Nz& \sqrt{N}y\\ x& \sqrt{N}y& v+x\end{array}\right).`$ (B19) The relevant part of the inverse is $`G_{AA}={\displaystyle \frac{1}{v}}{\displaystyle \frac{1}{v^2}}{\displaystyle \frac{xw+N(xzy^2)}{w+Nz}}.`$ (B20) The first term is the one-loop corrected single pole, while the second contains the expected double pole. Inserting this result into the definitions Eqs. (52) and (55), and using Eqs. (B1) and (B5), we can read off the required double-pole coefficient $`𝒟`$ $`=`$ $`𝒵_𝒜\left({\displaystyle \frac{xw+N(xzy^2)}{w+Nz}}\right)\|_{p^2=M_{AA}^2}.`$ (B21) Expanding in powers of $`\chi `$ we find $`𝒟`$ $``$ $`𝒵_𝒜\left\{{\displaystyle \frac{w}{N}}\left(1{\displaystyle \frac{w}{Nx_0}}\right)+(\delta x+\delta z2\delta y)\right\}\|_{p^2=M_{AA}^2}.`$ (B22) Note that $`(x_Vz_Vy_V^2)=0`$, so that a possible contribution proportional to $`m_0^2`$ cancels. Substituting and rearranging, we find the answer Eq. (59).
warning/0006/nucl-th0006052.html	ar5iv	text	# STRANGE HADRONIC STELLAR MATTER WITHIN THE BRUECKNER-BETHE-GOLDSTONE THEORY ## 1 Neutron stars within the BBG approach The nuclear matter equation of state (EoS) is the fundamental input for building models of neutron stars. These compact objects, among the densest in the universe, are indeed characterized by values of the density which span from the iron density at the surface up to eight-ten times normal nuclear matter density in the core. Therefore a detailed knowledge of the equation of state over a wide range of densities is required $`^\mathrm{?}`$. This is a very hard task from the theoretical point of view. In fact, whereas at densities close to the saturation value the matter consists mainly of nucleons and leptons, at higher densities several species of particles may appear due to the fast rise of the nucleon chemical potentials. In our work we perform microscopic calculations of the nuclear matter EoS containing fractions of $`\mathrm{\Lambda }`$ and $`\mathrm{\Sigma }^{}`$ hyperons in the framework of the Brueckner-Hartree-Fock (BHF) scheme $`^\mathrm{?}`$. The BHF approximation, with the continuous choice for the single particle potential, reproduces closely the many-body calculations up to the three hole-line level. In this approach, the basic input is the two-body interaction. We chose the Paris and the Argonne $`v_{18}`$ potential for the nucleon-nucleon (NN) part, whereas the Nijmegen soft-core model has been adopted for the nucleon-hyperon (NY) potential. No hyperon-hyperon interaction is taken into account, since no robust experimental data are available yet. For more details, the reader is referred to ref. $`^\mathrm{?}`$ and references therein. However, as commonly known, all many-body methods fail to reproduce the empirical nuclear matter saturation point $`\rho _0=0.17fm^3`$. This drawback is commonly corrected by introducing three-body forces (TBF’s) among nucleons. In our approach we have included a contribution containing a long range two-pion exchange attractive part and an intermediate range repulsive part $`^\mathrm{?}`$. This allows the correct reproduction of the saturation point. In figure 1 we show the chemical composition of $`\beta `$-stable and asymmetric nuclear matter containing hyperons (panel (a)) and the corresponding equation of state (panel (b)). The shown calculations have been performed using the Paris potential. We observe that hyperon formation starts at densities $`\rho 23`$ times normal nuclear matter density. The $`\mathrm{\Sigma }^{}`$ baryon appears earlier than the $`\mathrm{\Lambda }`$, in spite of its larger mass, because of the negative charge. The appearance of strange particles has two main consequences, i) an almost equal percentage of nucleons and hyperons are present in the stellar core at high densities and ii) a strong deleptonization of matter, since it is energetically convenient to maintain charge neutrality through hyperon formation than $`\beta `$-decay. The equation of state is displayed in panel (b). The dotted line represents the case when only nucleons and leptons are present in stellar matter, whereas the solid line shows the case when hyperons are included as well. In the latter case the equation of state gets very soft, since the kinetic energy of the already present baryonic species is converted into masses of the new particles, thus lowering the total pressure. This fact has relevant consequences for the structure of the neutron stars. ## 2 Equilibrium configurations of neutron stars We assume that a star is a spherically symmetric distribution of mass in hydrostatic equilibrium. The equilibrium configurations are obtained by solving the Tolman-Oppenheimer-Volkoff (TOV) equations $`^\mathrm{?}`$ for the pressure $`P`$ and the enclosed mass $`m`$, $`{\displaystyle \frac{dP(r)}{dr}}`$ $`=`$ $`{\displaystyle \frac{Gm(r)\rho (r)}{r^2}}{\displaystyle \frac{\left[1+\frac{P(r)}{\rho (r)}\right]\left[1+\frac{4\pi r^3P(r)}{m(r)}\right]}{1\frac{2Gm(r)}{r}}},`$ (1) $`{\displaystyle \frac{dm(r)}{dr}}`$ $`=`$ $`4\pi r^2\rho (r),`$ (2) being $`G`$ the gravitational constant (we assume $`c=1`$). Starting with a central mass density $`\rho (r=0)\rho _c`$, we integrate out until the pressure on the surface equals the one corresponding to the density of iron. This gives the stellar radius $`R`$ and the gravitational mass is then $$M_Gm(R)=4\pi _0^R𝑑rr^2\rho (r).$$ (3) For the outer part of the neutron star we have used the equations of state by Feynman-Metropolis-Teller $`^\mathrm{?}`$ and Baym-Pethick-Sutherland $`^\mathrm{?}`$, and for the medium-density regime we use the results of Negele and Vautherin $`^\mathrm{?}`$. For density $`\rho >0.08\mathrm{fm}^3`$ we use the microscopic equations of state obtained in the BHF approximation described above. For comparison, we also perform calculations of neutron star structure for the case of asymmetric and $`\beta `$-stable nucleonic matter. The results are plotted in Fig.2. We display the gravitational mass $`M_G`$ (in units of the solar mass $`M_o`$) as a function of the radius $`R`$ (panel (a)) and central baryon density $`n_c`$ (panel (b)). We note that the inclusion of hyperons lowers the value of the maximum mass from about 2.1 $`M_o`$ down to 1.26 $`M_o`$. This value lies below the value of the best observed pulsar mass, PSR1916+13, which amounts to 1.44 solar masses. However the observational data can be fitted if rotations are included, see dotted line in panel (b). In this case only equilibrium configurations rotating at the Kepler frequency $`\mathrm{\Omega }_K`$ are shown. In conclusion, the main finding of our work is the surprisingly low value of the maximum mass of a neutron star, which hardly comprises the observational data. This fact indicates how sensitive the properties of the neutron stars are to the details of the interaction. In particular our result calls for the need of including realistic hyperon-hyperon interactions. ## References
warning/0006/math0006016.html	ar5iv	text	# The Calibration Method for Free Discontinuity Problems ## 1 Introduction In De Giorgi introduced the name free discontinuity problems to denote a wide class of minimum problems for functionals of the form $$F(u):=_{\mathrm{\Omega }S_u}f(x,u(x),u(x))𝑑x+_{S_u}\psi (x,u^+(x),u^{}(x),\nu _u(x))𝑑^{n1},$$ (1) where $`\mathrm{\Omega }`$ is a given bounded domain in $`𝐑^n`$ with Lipschitz boundary, $`f:\mathrm{\Omega }\times 𝐑\times 𝐑^n[0,+\mathrm{}]`$ and $`\psi :\mathrm{\Omega }\times 𝐑\times 𝐑\times 𝐒^{n1}[0,+\mathrm{}]`$ are given Borel functions, $`𝐒^{n1}=\{v𝐑^n:\|v\|=1\}`$, $`^{n1}`$ is the $`(n1)`$-dimensional Hausdorff measure, and the unknown function $`u:\mathrm{\Omega }𝐑`$ is assumed to be regular out of a (partially regular) singular set $`S_u`$ of dimension $`n1`$, with unit normal $`\nu _u`$, on which $`u`$ admits unilateral traces $`u^+`$ and $`u^{}`$. The main feature of these problems is that the shape and location of the discontinuity set $`S_u`$ are not prescribed. Thus minimizing $`F`$ means optimizing both the function $`u`$ and the singular set $`S_u`$, which is indeed often regarded as an independent unknown. These problems have an increasing importance in many branches of applied analysis, such as image processing (Mumford-Shah functional for image segmentation) and fracture mechanics (Griffith’s criterion and Barenblatt cohesive zone model). The Mumford-Shah functional was introduced in in the context of a variational approach to image segmentation problems (for which we refer to ). It can be written, in dimension $`n`$, as $$F_g^{\alpha ,\beta }(u):=_{\mathrm{\Omega }S_u}\|u(x)\|^2𝑑x+\alpha ^{n1}(S_u)+\beta _{\mathrm{\Omega }S_u}\|u(x)g(x)\|^2𝑑x,$$ (2) where $`g`$ is a given function in $`L^{\mathrm{}}(\mathrm{\Omega })`$ (interpreted as the grey level of the image to be analysed), and $`\alpha >0`$ and $`\beta 0`$ are constants. When $`n=2`$ (the only case considered in image processing), the singular set $`S_u`$ of a minimizer $`u`$ of $`F_g^{\alpha ,\beta }`$ is interpreted as the set of the most relevant segmentation lines of the image. Using different classes of infinitesimal variations, one can show that every minimizer must satisfy certain equilibrium conditions, which could be globally called Euler-Lagrange equations for $`F_g^{\alpha ,\beta }`$. For instance, $`u`$ must satisfy the equation $`\mathrm{\Delta }u=\beta (ug)`$ on $`\mathrm{\Omega }S_u`$, with Neumann boundary conditions on $`S_u\mathrm{\Omega }`$. Moreover, there is a link between the mean curvature of $`S_u`$ (where defined) and the traces of $`u`$ and $`u`$ on the two sides of $`S_u`$; for instance, when $`\beta =0`$, the mean curvature of $`S_u`$ must be equal to the difference of the squares of the norms of the traces of $`u`$. Additional conditions have been derived for the two-dimensional case. We refer the reader to and for a precise description of these equilibrium conditions. However, since $`F_g^{\alpha ,\beta }`$ is not convex, all conditions which can be derived by infinitesimal variations are necessary for minimality, but never sufficient. The purpose of this note is precisely to present a sufficient condition for minimality (Theorem 3.1 for $`F_g^{\alpha ,\beta }`$ and Theorem 3.4 for $`F`$), and give a few applications (Examples 4.14.8). Detailed proofs and further results will be given in the forthcoming paper . ## 2 Notation and preliminaries For a complete mathematical treatment of the minimum problems for the functional $`F`$ considered in (1), we use the space $`SBV(\mathrm{\Omega })`$ of special functions of bounded variation, introduced by De Giorgi and Ambrosio in . A self-contained presentation of this space can be found in the recent book , which contains also the complete proof of the existence of a minimizer $`u`$ of $`F_g^{\alpha ,\beta }`$, and of the partial regularity of the corresponding singular set $`S_u`$ (the regularity of $`u`$ on $`\mathrm{\Omega }S_u`$ follows from the standard theory of elliptic equations). We recall that for every $`uSBV(\mathrm{\Omega })`$ the approximate upper and lower limits $`u^+(x)`$ and $`u^{}(x)`$ at a point $`x\mathrm{\Omega }`$ are defined by $$u^\pm (x):=\pm inf\{t𝐑:\underset{\rho 0+}{lim}\rho ^n^n(\{\pm u>t\}B_\rho (x))=0\},$$ where $`B_\rho (x)`$ is the open ball with centre $`x`$ and radius $`\rho `$. The singular set (or jump set) of $`u`$ is defined by $`S_u:=\{x\mathrm{\Omega }:u^{}(x)<u^+(x)\}`$. It is known that $`S_u`$ is countably $`(^{n1},n1)`$-rectifiable and that there exists a Borel measurable function $`\nu _u:S_u𝐒^{n1}`$ such that for $`^{n1}`$-a.e. $`xS_u`$ we have $$\underset{\rho 0+}{lim}\frac{1}{\rho ^n}_{B_\rho ^\pm (x)}\|u(y)u^\pm (x)\|𝑑y=0,$$ (3) where $`B_\rho ^\pm (x):=\{yB_\rho (x):\pm (yx)\nu _u(x)>0\}`$ and $``$ denotes the scalar product in $`𝐑^n`$ (see \[7, Theorem 4.5.9\]). Condition (3) says that $`\nu _u(x)`$ points from the side of $`S_u`$ corresponding to $`u^{}(x)`$ to the side corresponding to $`u^+(x)`$. The gradient $`Du`$ of $`u`$ is a measure that can be decomposed as the sum of two measures $`Du=D^au+D^su`$, where $`D^au`$ is absolutely continuous and $`D^su`$ is singular with respect to the Lebesgue measure $`^n`$. The density of $`D^au`$ with respect to $`^n`$ is denoted by $`u`$. Since $`uSBV(\mathrm{\Omega })`$, for every Borel set $`B`$ in $`\mathrm{\Omega }`$ we have $$(Du)(B)=_Bu(x)𝑑x+_{BS_u}(u^+(x)u^{}(x))\nu _u(x)𝑑^{n1}.$$ The graph of $`u`$ is defined as $$\mathrm{\Gamma }_u:=\{(x,t)\mathrm{\Omega }\times 𝐑:u^{}(x)tu^+(x)\}.$$ The characteristic function of the subgraph $`\{(x,t)\mathrm{\Omega }\times 𝐑:tu(x)\}`$ is denoted by $`1_u`$. It is defined by $`1_u(x,t):=1`$ if $`tu(x)`$, and $`1_u(x,t):=0`$ if $`t>u(x)`$. It belongs to $`SBV(\mathrm{\Omega }\times 𝐑)`$ and its gradient $`D1_u`$ is a measure concentrated on $`\mathrm{\Gamma }_u`$. ## 3 The main results We fix an open subset $`U`$ of $`\mathrm{\Omega }\times 𝐑`$ of the form $$U:=\{(x,t)\mathrm{\Omega }\times 𝐑:\tau _1(x)<t<\tau _2(x)\},$$ (4) where $`\tau _1`$ and $`\tau _2`$ are two continuous functions on $`\overline{\mathrm{\Omega }}`$ such that $`\mathrm{}\tau _1(x)\tau _2(x)+\mathrm{}`$ for every $`x\overline{\mathrm{\Omega }}`$. Let $`F`$ be the functional introduced in (1). We say that a function $`uSBV(\mathrm{\Omega })`$, with graph $`^n`$-contained in $`U`$ (i.e., $`^n(\mathrm{\Gamma }_uU)=0`$), is a Dirichlet $`U`$-minimizer of $`F`$, if $`F(u)F(v)`$ for every $`vSBV(\mathrm{\Omega })`$ with the same trace as $`u`$ on $`\mathrm{\Omega }`$ and with graph $`^n`$-contained in $`U`$. If the inequality $`F(u)F(v)`$ holds for every $`vSBV(\mathrm{\Omega })`$ with graph $`^n`$-contained in $`U`$, we say that $`u`$ is a $`U`$-minimizer of $`F`$. We omit $`U`$ when $`U=\mathrm{\Omega }\times 𝐑`$. The symbol $`\varphi `$ will always denote a bounded Borel measurable vectorfield defined on $`U`$ with values in $`𝐑^{n+1}=𝐑^n\times 𝐑`$, with components $`\varphi ^x𝐑^n`$ and $`\varphi ^t𝐑`$. The divergence of $`\varphi `$ is then $`\mathrm{div}\varphi (x,t)=\mathrm{div}_x\varphi ^x(x,t)+_t\varphi ^t(x,t)`$. We begin with a theorem concerning the functional $`F_g^{\alpha ,\beta }`$ introduced in (2). ###### Theorem 3.1 Let $`uSBV(\mathrm{\Omega })`$ with graph $`^n`$-contained in $`U`$. Assume that there exists a bounded vectorfield $`\varphi `$ of class $`C^1`$ on $`U`$ with the following properties: $`\frac{1}{4}\|\varphi ^x(x,t)\|^2\varphi ^t(x,t)+\beta \|tg(x)\|^2`$ for $`^n`$-a.e. $`x\mathrm{\Omega }`$ and for every $`\tau _1(x)<t<\tau _2(x)`$; $`\varphi ^x(x,u(x))=2u(x)`$ and $`\varphi ^t(x,u(x))=\|u(x)\|^2\beta \|u(x)g(x)\|^2`$ for $`^n`$-a.e. $`x\mathrm{\Omega }`$; $`\left\|{\displaystyle _{t_1}^{t_2}}\varphi ^x(x,t)𝑑t\right\|\alpha `$ for $`^{n1}`$-a.e. $`x\mathrm{\Omega }`$ and for every $`\tau _1(x)<t_1<t_2<\tau _2(x)`$; $`{\displaystyle _{u^{}(x)}^{u^+(x)}}\varphi ^x(x,t)𝑑t=\alpha \nu _u(x)`$ for $`^{n1}`$-a.e. $`xS_u`$; $`\mathrm{div}\varphi (x,t)=0`$ for every $`(x,t)U`$. Then $`u`$ is a Dirichlet $`U`$-minimizer of $`F_g^{\alpha ,\beta }`$. If, in addition, $`\varphi ^x(x,t)`$ satisfies the boundary condition $`\underset{(y,s)(x,t)}{lim}\varphi ^x(y,s)\nu (x)=0`$ for $`^{n1}`$-a.e. $`x\mathrm{\Omega }`$ and for $`^1`$-a.e. $`t[\tau _1(x),\tau _2(x)]`$, where $`\nu (x)`$ is the outer unit normal to $`\mathrm{\Omega }`$, then $`u`$ is a $`U`$-minimizer of $`F_g^{\alpha ,\beta }`$. A vectorfield $`\varphi `$ which satisfies conditions $`(\mathrm{a1})`$$`(\mathrm{c1})`$ of Theorem 3.1 is called a calibration for the functional $`F_g^{\alpha ,\beta }`$ on $`U`$. If $`\varphi `$ satisfies also $`(\mathrm{c2})`$, it is called a Neumann calibration. Theorem 3.1 is an immediate consequence of the following lemmas. ###### Lemma 3.2 Let $`\varphi `$ be a vectorfield which satisfies conditions $`(\mathrm{a1})`$ and $`(\mathrm{b1})`$ of Theorem 3.1. Then for every $`uSBV(\mathrm{\Omega })`$ with graph $`^n`$-contained in $`U`$ we have $$F_g^{\alpha ,\beta }(u)_U\varphi d(D1_u).$$ (5) Moreover, equality holds in (5) for a given $`u`$ if and only if conditions $`(\mathrm{a2})`$ and $`(\mathrm{b2})`$ of Theorem 3.1 are satisfied. The next lemma is a consequence of the divergence theorem. ###### Lemma 3.3 Suppose that $`\varphi `$ is of class $`C^1`$ and that $`\mathrm{div}\varphi =0`$ on $`U`$. Then $$_U\varphi d(D1_u)=_U\varphi d(D1_v)$$ (6) for every pair of functions $`u`$, $`v`$ in $`BV(\mathrm{\Omega })`$ with the same trace on $`\mathrm{\Omega }`$ and with graphs $`^n`$-contained in $`U`$. If, in addition, $`\varphi `$ satisfies condition $`(\mathrm{c2})`$ of Theorem 3.1, then (6) holds for every pair of functions $`u`$, $`v`$ in $`BV(\mathrm{\Omega })`$ with graphs $`^n`$-contained in $`U`$. As a matter of fact, the method of calibrations can be easily adapted to the functional $`F`$ defined in (1). ###### Theorem 3.4 Let $`uSBV(\mathrm{\Omega })`$ with graph $`^n`$-contained in $`U`$. Assume that there exists a bounded vectorfield $`\varphi `$ of class $`C^1`$ on $`U`$ with the following properties: $`\varphi ^x(x,t)v\varphi ^t(x,t)+f(x,t,v)`$ for $`^n`$-a.e. $`x\mathrm{\Omega }`$, for every $`\tau _1(x)<t<\tau _2(x)`$, and for every $`v𝐑^n`$; $`\varphi ^x(x,u(x))u(x)=\varphi ^t(x,u(x))+f(x,u(x),u(x))`$ for $`^n`$-a.e. $`x\mathrm{\Omega }`$; $`\nu {\displaystyle _{t_1}^{t_2}}\varphi ^x(x,t)𝑑t\psi (x,t_1,t_2,\nu )`$ for $`^{n1}`$-a.e. $`x\mathrm{\Omega }`$, for every $`\tau _1(x)<t_1<t_2<\tau _2(x)`$, and for every $`\nu 𝐒^{n1}`$; $`\nu _u(x){\displaystyle _{u^{}(x)}^{u^+(x)}}\varphi ^x(x,t)𝑑t=\psi (x,u^{}(x),u^+(x),\nu _u(x))`$ for $`^{n1}`$-a.e. $`xS_u`$; $`\mathrm{div}\varphi (x,t)=0`$ for every $`(x,t)U`$. Then $`u`$ is a Dirichlet $`U`$-minimizer of $`F`$. If $`\varphi ^x(x,t)`$ satisfies also the boundary condition $`(\mathrm{c2})`$ of Theorem 3.1, then $`u`$ is a $`U`$-minimizer of $`F`$. ###### Remark 3.5 We note that in Theorem 3.4 there is no regularity or convexity hypothesis on $`f`$ or $`\psi `$. If $`f^{}(x,t,v^{})`$ is the the convex conjugate of $`f(x,t,v)`$ with respect to $`v`$, condition $`(\mathrm{a1})`$ is equivalent to $`f^{}(x,t,\varphi ^x(x,t))\varphi ^t(x,t)`$ for $`^n`$-a.e. $`x\mathrm{\Omega }`$ and for every $`\tau _1(x)<t<\tau _2(x)`$. If this condition is satisfied, and $`f(x,t,v)`$ is convex and differentiable with respect to $`v`$, then condition $`(\mathrm{a2})`$ is equivalent to $`\{\begin{array}{cc}\varphi ^x(x,u(x))=_vf(x,u(x),u(x))\hfill & \\ \varphi ^t(x,u(x))=f^{}(x,u(x),\varphi ^x(x,u(x)))\hfill & \end{array}`$ for $`^n`$-a.e. $`x\mathrm{\Omega }`$. ###### Remark 3.6 In Theorems 3.1 and 3.4 the hypothesis that $`\varphi `$ is of class $`C^1`$ is too strong for many applications. It is used only in Lemma 3.3 and it can be relaxed in several ways (see for details). For instance, one may consider piecewise $`C^1`$ vectorfields, which may be discontinuous along sufficiently regular interfaces. In this case the divergence-free condition $`(\mathrm{c1})`$ must be understood in the distributional sense, i.e., the pointwise divergence vanishes (where defined) and the normal component of $`\varphi `$ is continuous across the discontinuity surfaces. ## 4 Some examples The following examples show that the calibration method is very flexible, and can be used to prove the minimality of a given function $`u`$ in many different situations. In the first examples we will consider only the “homogeneous” functional $`F^\alpha :=F_g^{\alpha ,0}`$, in which the lower order term $`\beta _\mathrm{\Omega }\|ug\|^2𝑑x`$ vanishes. ###### Example 4.1 (Affine function in one dimension) Let $`n:=1`$, $`\mathrm{\Omega }:=]0,a[`$, and $`u(x):=\lambda x`$, with $`\lambda >0`$. It is easy to see that $`u`$ is a Dirichlet minimizer of $`F^\alpha `$ if and only if $`a\lambda ^2\alpha `$. In this case a calibration is given by the piecewise constant function $$\varphi (x,t):=\{\begin{array}{cc}(2\lambda ,\lambda ^2),\hfill & \text{if }\frac{\lambda }{2}xt\frac{\lambda }{2}(x+a)\text{,}\hfill \\ (0,0),\hfill & \text{otherwise.}\hfill \end{array}$$ (7) Another calibration is given by $$\varphi (x,t):=\{\begin{array}{cc}(2\frac{t}{x},(\frac{t}{x})^2),\hfill & \text{if }0t\lambda x\text{,}\hfill \\ (2\frac{\lambda at}{ax},(\frac{\lambda at}{ax})^2),\hfill & \text{if }\lambda xt\lambda a\text{,}\hfill \\ (0,0),\hfill & \text{otherwise.}\hfill \end{array}$$ (8) If $`a\lambda ^2>\alpha `$, then the function $`u(x):=\lambda x`$ is not a Dirichlet minimizer of $`F^\alpha `$, but it is still a Dirichlet $`U`$-minimizer with $$U:=\{(x,t)]0,a[\times 𝐑:\lambda x\frac{\alpha }{4\lambda }<t<\lambda x+\frac{\alpha }{4\lambda }\}.$$ A calibration on $`U`$ is given by $`\varphi (x,t):=(2\lambda ,\lambda ^2)`$. ###### Example 4.2 (Jump in one dimension) Let $`n:=1`$, $`\mathrm{\Omega }:=]0,a[`$, $`u(x):=0`$ for $`0<x<c`$, and $`u(x):=h`$ for $`c<x<a`$, with $`0<c<a`$ and $`h>0`$. It is easy to see that $`u`$ is a Dirichlet minimizer of $`F^\alpha `$ if and only if $`a\alpha h^2`$. In this case two different calibrations are given by (7) and (8) with $`\lambda =\sqrt{\alpha }/\sqrt{a}`$. Suppose now that $`a\alpha >h^2`$. Let $`\epsilon >0`$ be a constant such that $`2\epsilon +\sqrt{2\alpha \epsilon }h`$, let $$\tau _1(x)=\{\begin{array}{cc}\epsilon ,\hfill & \text{if }xc\text{,}\hfill \\ \epsilon +\frac{h}{\epsilon }(xc),\hfill & \text{if }cxc+\epsilon \text{,}\hfill \\ h\epsilon ,\hfill & \text{if }c+\epsilon x\text{,}\hfill \end{array}$$ let $`\tau _2(x)=\tau _1(x+\epsilon )+2\epsilon `$, and let $`U`$ be the open set defined by (4). Then $`u`$ is a Dirichlet $`U`$-minimizer of $`F^\alpha `$, and a calibration on $`U`$ is given by the piecewise constant function $$\varphi (x,t):=\{\begin{array}{cc}(2\lambda ,\lambda ^2),\hfill & \text{if }c\epsilon <x<c+\epsilon \text{ and}\hfill \\ & \epsilon +\frac{\lambda }{2}(xc+\epsilon )<t<\epsilon +\frac{\lambda }{2}(xc+\epsilon )+\frac{\alpha }{2\lambda }\text{,}\hfill \\ (0,0),\hfill & \text{otherwise,}\hfill \end{array}$$ where $`\lambda >0`$ is any constant such that $`\epsilon +\epsilon \lambda +\frac{\alpha }{2\lambda }h\epsilon `$, for instance $`\lambda =\sqrt{\alpha }/\sqrt{2\epsilon }`$. ###### Example 4.3 (Harmonic function) Let $`\mathrm{\Omega }`$ be a bounded domain in $`𝐑^n`$, $`n`$ arbitrary, and let $`u`$ be a harmonic function on $`\mathrm{\Omega }`$. As pointed out by Chambolle , $`u`$ is a Dirichlet minimizer of $`F^\alpha `$ if $$\underset{\mathrm{\Omega }}{osc}u\underset{\mathrm{\Omega }}{sup}\|u\|\alpha ,$$ (9) where $`osc_\mathrm{\Omega }u:=sup_\mathrm{\Omega }uinf_\mathrm{\Omega }u`$. Note that for $`n=1`$ this condition reduces to the constraint $`a\lambda ^2\alpha `$ of Example 4.1. Inspired by the one dimensional case (see (7)), we construct the calibration $$\varphi (x,t):=\{\begin{array}{cc}(2u(x),\|u(x)\|^2),\hfill & \text{if }\frac{1}{2}(u(x)+m)t\frac{1}{2}(u(x)+M)\text{,}\hfill \\ (0,0),\hfill & \text{otherwise,}\hfill \end{array}$$ (10) where $`m:=inf_\mathrm{\Omega }u`$ and $`M:=sup_\mathrm{\Omega }u`$. Another calibration (see (8)) is given by $$\varphi (x,t):=\{\begin{array}{cc}(2\frac{tm}{u(x)m}u(x),(\frac{tm}{u(x)m})^2\|u(x)\|^2),\hfill & \text{if }mtu(x)\text{,}\hfill \\ (2\frac{Mt}{Mu(x)}u(x),(\frac{Mt}{Mu(x)})^2\|u(x)\|^2),\hfill & \text{if }u(x)tM\text{,}\hfill \\ (0,0),\hfill & \text{otherwise.}\hfill \end{array}$$ If (9) is not satisfied, $`u`$ is still is a Dirichlet $`U`$-minimizer of $`F^\alpha `$, for $$U:=\{(x,t)\mathrm{\Omega }\times 𝐑:u(x)\frac{\alpha }{4}\|u(x)\|^1<t<u(x)+\frac{\alpha }{4}\|u(x)\|^1\},$$ (11) and a calibration in $`U`$ is given by $`\varphi (x,t):=(2u(x),\|u(x)\|^2)`$. ###### Example 4.4 (Pure jump) Let $`n2`$ and let $`\mathrm{\Omega }:=]0,a[\times V`$, where $`V`$ is a bounded domain in $`𝐑^{n1}`$ with Lipschitz boundary. Denoting the first coordinate of $`x`$ by $`x_1`$, let $`u(x):=0`$ for $`0<x_1<c`$, and $`u(x):=h`$ for $`c<x_1<a`$, with $`0<c<a`$ and $`h>0`$. Using the results of Example 4.2 it is easy to see that $`u`$ is a Dirichlet minimizer of $`F^\alpha `$ if $`a\alpha h^2`$. In this case two different calibrations can be constructed in the following way: the projection of these calibrations onto the $`(x_1,t)`$-plane are given by (7) and (8), with $`\lambda =\sqrt{\alpha }/\sqrt{a}`$ and $`x`$ replaced by $`x_1`$, while all other components of these calibrations vanish. If $`a\alpha >h^2`$, it may happen that $`u`$ is still a Dirichlet minimizer of $`F^\alpha `$. For instance, if $`n=2`$ and $`V=]0,b[`$, with $`b\alpha \pi 2h^2`$, a different calibration has been constucted in . Therefore $`u`$ is a Dirichlet minimizer of $`F^\alpha `$ even if $`a\alpha `$ is very large with respect to $`h^2`$, provided that $`b\alpha `$ is small enough. Arguing as in the last part of Example 4.2 one can prove that for every $`a`$ and $`V`$ there exists an open set $`U`$ of the form (4), containing $`\mathrm{\Gamma }_u`$, such that $`u`$ is a Dirichlet $`U`$-minimizer of $`F^\alpha `$. ###### Example 4.5 (Triple junction) Let $`n:=2`$, let $`\mathrm{\Omega }:=B(0,r)`$ be the open ball with radius $`r>0`$ centered at the origin, and let $`u`$ be given, in polar coordinates, by $`u(\rho ,\theta ):=a`$ for $`0\theta <\frac{2}{3}\pi `$, $`u(\rho ,\theta ):=b`$ for $`\frac{2}{3}\pi \theta <\frac{4}{3}\pi `$, and $`u(\rho ,\theta ):=c`$ for $`\frac{4}{3}\pi \theta <2\pi `$, where $`a`$, $`b`$, and $`c`$ are distinct constants. Thus $`S_u`$ is given by three line segments meeting at the origin with equal angles. If $$2\alpha r\mathrm{min}\{\|ab\|^2,\|bc\|^2,\|ca\|^2\},$$ (12) then $`u`$ is a Dirichlet minimizer of $`F^\alpha `$. To construct a calibration, it is not restrictive to assume $`a<b=0<c`$. Inspired by the one dimensional case described in Example 4.2, we take $`e_\pm :=(\pm \sqrt{3}/2,1/2)`$, and $`\lambda >0`$ such that $`\frac{\lambda r}{2}+\frac{\alpha }{\lambda }\mathrm{min}\{a,c\}`$ (which is possible by (12)), and we define the calibration by $$\varphi (x,t):=\{\begin{array}{cc}(\lambda e_+,\lambda ^2/4),\hfill & \text{if }\frac{\lambda }{4}(r+xe_+)t\frac{\lambda }{4}(r+xe_+)+\frac{\alpha }{\lambda }\text{,}\hfill \\ (\lambda e_{},\lambda ^2/4),\hfill & \text{if }\frac{\lambda }{4}(r+xe_{})\frac{\alpha }{\lambda }t\frac{\lambda }{4}(r+xe_{})\text{,}\hfill \\ (0,0),\hfill & \text{otherwise.}\hfill \end{array}$$ (13) If $`\alpha r`$ is much larger than $`\mathrm{min}\{\|ab\|^2,\|bc\|^2,\|ca\|^2\}`$, it is easy to construct a comparison function $`v`$ with the same boundary values as $`u`$ and such that $`F^\alpha (v)<F^\alpha (u)`$. This shows that in this case $`u`$ is not a Dirichlet minimizer. However, for every value of the parameters $`\alpha `$, $`r`$, $`a`$, $`b`$, $`c`$, one can construct a suitable neighbourhood $`U`$ of the graph $`\mathrm{\Gamma }_u`$, of the form (4), such that a variant of (13) is a calibration in $`U`$, and therefore $`u`$ is a Dirichlet $`U`$-minimizer of $`F^\alpha `$. We refer to for the details. We consider now the functional $`F_g^{\alpha ,\beta }`$, with $`\beta >0`$. ###### Example 4.6 (Solution of the Neumann problem) Let $`\mathrm{\Omega }`$ be a bounded open set in $`𝐑^n`$ with boundary of class $`C^{1,\epsilon }`$ for some $`\epsilon >0`$, and let $`u`$ be the solution of the Neumann problem $$\{\begin{array}{cc}\mathrm{\Delta }u=\beta (ug)\hfill & \text{on }\mathrm{\Omega }\text{,}\hfill \\ \frac{u}{\nu }=0\hfill & \text{on }\mathrm{\Omega }\text{,}\hfill \end{array}$$ (14) with $`\beta >0`$ and $`gL^{\mathrm{}}(\mathrm{\Omega })`$. Assume that condition (9) of Example 4.3 is satisfied. Then $`u`$ is a minimizer of $`F_g^{\alpha ,\beta }`$. If the strict inequality holds in (9), then $`u`$ is the unique minimizer. A Neumann calibration $`\varphi (x,t)`$ is given by $$\{\begin{array}{cc}(0,\beta \|\frac{m}{2}\frac{u(x)}{2}\|^2\beta \|\frac{m}{2}+\frac{u(x)}{2}g(x)\|^2),\hfill & \text{if }t\frac{u(x)}{2}<\frac{m}{2}\text{,}\hfill \\ (2u(x),\|u(x)\|^2\beta \|tg(x)\|^2+\beta \|tu(x)\|^2),\hfill & \text{if }\frac{m}{2}t\frac{u(x)}{2}\frac{M}{2}\text{,}\hfill \\ (0,\beta \|\frac{M}{2}\frac{u(x)}{2}\|^2\beta \|\frac{M}{2}+\frac{u(x)}{2}g(x)\|^2),\hfill & \text{if }\frac{M}{2}<t\frac{u(x)}{2}\text{,}\hfill \end{array}$$ where $`m:=inf_\mathrm{\Omega }u`$ and $`M:=sup_\mathrm{\Omega }u`$. If (9) is not satisfied, $`u`$ is still is a $`U`$-minimizer of $`F_g^{\alpha ,\beta }`$, where $`U`$ is the open set defined by (11). A Neumann calibration on $`U`$ is given by $$\varphi (x,t):=(2u(x),\|u(x)\|^2\beta \|tg(x)\|^2+\beta \|tu(x)\|^2).$$ The hypothesis that $`\mathrm{\Omega }`$ is of class $`C^{1,\epsilon }`$ is used only to obtain the boundary condition $`(\mathrm{c2})`$ of Theorem 3.1, which, in this case, becomes $$\underset{yx}{lim}u(y)\nu (x)=0\text{for }^{n1}\text{-a.e. }x\mathrm{\Omega }.$$ (15) It is clear that (15) is still true if for $`^{n1}`$-a.e. $`x\mathrm{\Omega }`$ there exists an open neighbourhood $`V_x`$ of $`x`$ in $`𝐑^n`$ such that $`V_x\mathrm{\Omega }`$ is a manifold of class $`C^{1,\epsilon }`$ (see \[2, Theorem 7.5.2\]). Therefore the result of this example is true also when $`\mathrm{\Omega }`$ is polyhedral. In the next examples we construct a calibration for $`F_g^{\alpha ,\beta }`$ when the parameter $`\beta `$ is large enough. ###### Example 4.7 (Smooth $`g`$ and large $`\beta `$) Let $`\mathrm{\Omega }`$ be a bounded open set in $`𝐑^n`$ with smooth boundary, and let $`gC^2(\overline{\mathrm{\Omega }})`$. There exists a constant $`\beta _00`$, depending on $`g`$ and $`\alpha `$, such that for every $`\beta >\beta _0`$ the solution $`u`$ of the Neumann problem (14) of Example 4.6 is the unique minimizer of $`F_g^{\alpha ,\beta }`$. A Neumann calibration is constructed in . This shows that the minimizer of $`F_g^{\alpha ,\beta }`$ is smooth, provided that $`g`$ is smooth and $`\beta `$ is large enough. Therefore the solution of the image segmentation problem ($`n=2`$) based on the minimization of $`F_g^{\alpha ,\beta }`$ has an empty set of segmentation lines if the “grey level” function $`g`$ is smooth and the parameter $`\beta `$ in the fidelity term $`\beta _\mathrm{\Omega }\|ug\|^2𝑑x`$ is large. ###### Example 4.8 (Function $`g`$ with only two values) Let $`\mathrm{\Omega }`$ be an open set in $`𝐑^n`$ and let $`E`$ be a compact set contained in $`\mathrm{\Omega }`$ with boundary of class $`C^2`$. Let $`g(x):=a`$ for $`xE`$ and $`g(x):=b`$ for $`x\mathrm{\Omega }E`$, with $`ab`$. There exists a constant $`\beta _00`$, depending on $`g`$ and $`\alpha `$, such that for every $`\beta >\beta _0`$ the function $`u:=g`$ is the unique minimizer of $`F_g^{\alpha ,\beta }`$. To construct a calibration, it is not restrictive to assume $`a<b`$. We take a $`C^1`$ vectorfield $`v:\mathrm{\Omega }𝐑^n`$ with compact support in $`\mathrm{\Omega }`$ such that $`\|v(x)\|1`$ for every $`x\mathrm{\Omega }`$ and $`v(x)`$ is the outer unit normal to $`E`$ for every $`xE`$. Then we set $`\varphi ^x(x,t)=\sigma (t)v(x)`$, where $`\sigma `$ is a fixed positive smooth function with integral equal to $`\alpha `$ and support contained in $`]a,b[`$. We see that conditions $`(\mathrm{b1})`$, $`(\mathrm{b2})`$, and $`(\mathrm{c2})`$ of Theorem 3.1 are satisfied by construction. It remains to choose $`\varphi ^t`$ so that $`(\mathrm{a1})`$, $`(\mathrm{a2})`$, and $`(\mathrm{c1})`$ hold. Condition $`(\mathrm{a2})`$ forces us to set $`\varphi ^t(x,t)=0`$ for $`t=g(x)`$, while $`(\mathrm{c1})`$ gives $`_t\varphi ^t(x,t)=\sigma (t)\mathrm{div}_xv(x)`$. These two conditions determine $`\varphi ^t(x,t)`$ at every point $`(x,t)`$. It is then easy to see that $`(\mathrm{a1})`$ holds if $`\beta `$ is large enough. We refer to for the details. This example shows that, if $`gSBV(\mathrm{\Omega })`$ has only two values, and $`S_g`$ is smooth enough, then the minimizer of the Mumford-Shah functional $`F_g^{\alpha ,\beta }`$ reconstructs $`g`$ exactly, when $`\beta `$ is large enough. Recently the following question has been studied by using the calibration method: is it true that a function $`u`$ is a (Dirichlet) minimizer of $`F_g^{\alpha ,\beta }`$, if it satisfies the Euler-Lagrange equations and the domain $`\mathrm{\Omega }`$ is sufficiently small? For the moment we have only a partial answer. In we have considered the case where $`n:=2`$ and $`S_u`$ is a line segment joining two points of the boundary of $`\mathrm{\Omega }`$. If $`u`$ satisfies the Euler-Lagrange equations for the “homogeneous functional” $`F^\alpha :=F_g^{\alpha ,0}`$, then for every $`x_0S_u`$ there exists an open neighbourhood $`\mathrm{\Omega }_0`$ of $`x_0`$, contained in $`\mathrm{\Omega }`$, such that $`u`$ is a Dirichlet minimizer of $`F^\alpha `$ in $`\mathrm{\Omega }_0`$. The minimality is proved by constructing a complicated calibration on $`\mathrm{\Omega }_0\times 𝐑`$. This result has been extended in to the case where $`S_u`$ is an analytic curve joining two points of $`\mathrm{\Omega }`$. The (more difficult) construction of the calibration presented in this paper shows that one can take the same set $`\mathrm{\Omega }_0`$ for every $`x_0S_u`$; in other words, one can take as $`\mathrm{\Omega }_0`$ a suitable tubular neighbourhood of $`S_u`$. Moreover, it is proved in that an additional condition on $`u`$ and $`S_u`$ implies that $`u`$ is a Dirichlet $`U`$-minimizer for a suitable open neighbourhood $`U`$ of the graph $`\mathrm{\Gamma }_u`$. A counterexample (where $`S_u`$ is a line segment joining two points of $`\mathrm{\Omega }`$) shows that this is not always true when $`u`$ is just a solution of the Euler-Lagrange equations with $`S_u\mathrm{}`$, in contrast to the case $`S_u=\mathrm{}`$ (see Example 4.6).
warning/0006/hep-ph0006275.html	ar5iv	text	# Cosmology and Hierarchy in Stabilized Warped Brane Models ## Abstract We examine the cosmology and hierarchy of scales in models with branes immersed in a five-dimensional curved spacetime subject to radion stabilization. When the radion field is time-independent and the inter-brane spacing is stabilized, the universe can naturally find itself in the radiation-dominated epoch. This feature is independent of the form of the stabilizing potential. We recover the standard Friedmann equations without assuming a specific form for the bulk energy-momentum tensor. In the models considered, if the observable brane has positive tension, a solution to the hierarchy problem requires the presence of a negative tension brane somewhere in the bulk. We find that the string scale can be as low as the electroweak scale. In the situation of self-tuning branes where the bulk cosmological constant is set to zero, the brane tensions have hierarchical values. In the case of a polynomial stabilizing potential no new hierarchy is created. preprint: MADPH–00-1184 hep-ph/0006275 I. Introduction. It has recently been realized that the string scale can be much lower than the Planck scale and even close to the electroweak scale . A low string scale provides new avenues on solving the hierarchy problem . The argument resides in the fact that a low string scale ($`M_X`$) may result in the apparent size of the Planck scale ($`M_{Pl}`$) due to the existence of a large volume ($`R^n`$) of compact extra dimensions, $`M_{Pl}^2M_X^{n+2}R^n`$ . Such a scenario may lead to a rich phenomenology at low energies and is thus testable at collider experiments . Recently, a model involving just one extra dimension with a background $`AdS_5`$ metric was proposed by Randall and Sundrum (see also Ref. ). In this scenario, two branes (one with positive tension and the other with negative tension) are located at the fixed points of an $`S^1/𝐙_\mathrm{𝟐}`$ orbifold in a bulk with negative cosmological constant. An exponential hierarchy between the physical scales on the two branes is generated due to the curved spacetime, providing an explanation for the large hierarchy between the weak and Planck scales. The model is amenable to a holographic interpretation motivated by string theories . However, the Randall-Sundrum model has some drawbacks. First, a perfect fine-tuning among the brane tensions and the bulk cosmological constant is needed to guarantee a static solution for the warped metric of spacetime. Mechanisms for stabilizing the brane locations via interactions between a bulk scalar field (called the radion) and the branes were suggested in and elaborated in , where an elegant solution that accounts for the back-reaction of the scalar profile on the geometry is outlined. An overall fine-tuning equivalent to setting the four-dimensional cosmological constant to zero is still present. Second, it was found that the brane world may not lead to the standard cosmology . In particular, the Hubble parameter $`H`$ was found to be proportional to the matter density $`\rho `$ , in contradiction with the usual $`H\sqrt{\rho }`$ behavior. Although this can be remedied by a fine-tuned cancellation between the brane tension and the bulk cosmological constant , only with a negative energy density on the observable brane is the standard cosmology recovered . Much attention has been devoted to studying cosmology without an explicit stabilization mechanism . The connection between radion stabilization and cosmology was explored in , and the standard cosmology can be obtained if the radion is time-dependent . There have been attempts to modify the Randall-Sundrum model so that the observable brane has positive tension , since the localization of matter and gauge fields on positive tension branes is well understood in string theory. An initial study with two positive tension branes was carried out , incorporating localized gravity in a noncompact $`AdS_5`$ geometry . The necessary hierarchy can be generated between the Planck and electroweak scales by placing the hidden and observable branes at specific locations in the infinite dimension. We will refer to this as the Lykken-Randall model. In this letter we examine the cosmology and hierarchy in models with radion stabilization. We shall adopt the formalism of Ref. to stabilize the brane separation. We call this the Solution Generating Technique. We generalize it to the case of branes with arbitrary tensions and no relation between the metrics on either side of the branes. Using this technique we study the cosmology resulting from radion stabilization and find a cancellation between the bulk cosmological constant and the brane tension. This approach is rather different in philosophy from the one adopted in Refs. , where the bulk energy-momentum tensor extracted from the linearized field equations is chosen specifically to get the conventional cosmology. We obtain the important result that the stabilization of the inter-brane spacing can naturally lead to the cosmology of the radiation-dominated universe. Specifically, this is a consequence of requiring consistency in the equation of motion of the radion before and after perturbing the solutions by placing matter energy density on the observable brane. We argue that to obtain the complete evolution of our universe, a time-dependence of the radion field on the brane should be introduced. We speculate on the existence of some new dynamics that causes the transition from one epoch to the next. Guided by cosmology, we explore the consequences on the hierarchy between the Planck and electroweak scales. As concrete examples, we study two classes of models. The first is the “self-tuning” model, motivated by recent attempts to solve the cosmological constant problem . A dilaton-like coupling of a bulk scalar field with a brane was shown to result in a vanishing bulk cosmological constant. The result persists irrespective of the tension on the brane. This feature is referred to as self-tuning. The other model involves a Higgs-like radion potential, similar to that in . In both models, the cosmology on the observable brane is independent of the configuration of branes and the potential that leads to radion stabilization. All that is required is some stabilizing potential and that the observable brane has positive tension. To generate the hierarchy of scales, at least one hidden brane with negative tension is required. The latter cannot be positioned at an orbifold fixed point. The self-tuning brane model has the following properties: 1. The model illustrates the unique minimal configuration from which the hierarchy of scales can be obtained without fine-tuning. There are two positive tension branes (one of which is the observable brane), and one negative tension brane. The values of the brane tensions become hierarchical. 2. A dilatonic coupling between the branes and the bulk stabilizes the inter-brane spacings. Thus the radion may be identified with the dilaton. 3. The same coupling ensures self-tuning of the branes to be flat and the bulk cosmological constant to be zero. The tree-level contribution to the four-dimensional cosmological constant is eliminated. We will truncate the space to avoid curvature singularities. The model with a Higgs-like radion potential possesses the following properties: 1. The minimal set-up to generate the scale hierarchy requires one positive tension observable brane and one negative tension hidden brane. 2. The extra dimension is linearly infinite with finite proper volume. 3. The radion field is unbounded above and leads to the bulk cosmological constant being unbounded below. In Section II we present the Solution Generating Technique. Section III is devoted to studying the cosmology in a general setting. In Section IV we demonstrate a realization of a Lykken-Randall-like model with self-tuning branes. In Section V we perform a similar analysis but with a polynomial superpotential. We conclude in Section VI. II. Solution Generating Technique. It is possible to use a gauged supergravity-inspired approach to reduce the nonlinear classical field equations of brane models with scalar-tensor gravity to a system of decoupled first order differential equations . Using the technique of , the brane spacing can be stabilized. We assume the scalar field to be static to prevent the four-dimensional Planck mass from being time-dependent. The formalism is independent of whether the fifth dimension is compact or noncompact. We shall present the arguments for the case of a noncompact dimension. We assume the presence of three 3-branes in the space $`(\mathrm{},\mathrm{})`$, at $`y_0=0`$, $`y_1`$ and $`y_2`$ with the observable brane located at $`y_0`$. We will refer to the branes at $`y_1`$ and $`y_2`$ as hidden branes. The four-dimensional metric on the brane labelled by its position $`y_i`$ is $`g_{\mu \nu }^{(i)}(x^\mu )g_{\mu \nu }(x^\mu ,y=y_i)`$, where $`g_{AB}`$ is the five-dimensional metric and $`A,B=0,1,2,3,5`$ and $`\mu ,\nu =0,1,2,3`$. We use the metric signature $`(,+,+,+,+)`$. The five-dimensional gravitational action including a scalar field $`\varphi (y)`$ is $`S=S_{Gravity}+S_{Brane}`$ with $`S_{Gravity}`$ $`=`$ $`{\displaystyle d^4x𝑑y\sqrt{g}\{\frac{1}{2\kappa ^2}R\frac{1}{4\kappa ^2}_A\varphi ^A\varphi \mathrm{\Lambda }(\varphi )\}},`$ (1) $`S_{Brane}`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{2}{}}}{\displaystyle d^4x𝑑y\sqrt{g^{(i)}}V_i(\varphi )\delta (yy_i)},`$ (2) where $`\kappa ^2=8\pi G_N^{(5)}=M_X^3`$ is the five-dimensional coupling constant of gravity, $`M_X`$ is the Planck scale in five dimensions, $`R`$ is the curvature scalar and $`V_i(\varphi _i)`$ is the tension of the brane at $`y_i`$. $`\mathrm{\Lambda }(\varphi )`$ is the potential of the field $`\varphi `$ in the bulk and is interpreted as the cosmological constant although it has a $`\varphi `$-dependence. We allow it to be discontinuous at the branes, but continuous in each section. We write $`\mathrm{\Lambda }(\varphi )`$ as $`\mathrm{\Lambda }_0(\varphi )`$ if $`y<0`$, $`\mathrm{\Lambda }_1(\varphi )`$ if $`\mathrm{\hspace{0.17em}0}<y<y_1`$ as $`\mathrm{\Lambda }_2(\varphi )`$ if $`y_1<y<y_2`$ and $`\mathrm{\Lambda }_3(\varphi )`$ if $`y>y_2`$ (see Fig. 1). The five-dimensional Einstein equations arising from the above action are $`G_{AB}R_{AB}{\displaystyle \frac{1}{2}}g_{AB}R=\kappa ^2T_{AB}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(_A\varphi _B\varphi {\displaystyle \frac{1}{2}}g_{AB}(\varphi )^2)\kappa ^2g_{AB}\mathrm{\Lambda }(\varphi )`$ (4) $`\kappa ^2{\displaystyle \underset{i=0}{\overset{2}{}}}V_i(\varphi )\sqrt{{\displaystyle \frac{g^{(i)}}{g}}}g_{\mu \nu }^{(i)}\delta _A^\mu \delta _B^\nu \delta (yy_i),`$ where $`R_{AB}`$ is the five-dimensional Ricci tensor. The most general five-dimensional metric that respects four-dimensional Poincaré symmetry is $$ds^2=e^{2A(y)}\eta _{\mu \nu }dx^\mu dx^\nu +(dy)^2.$$ (5) The factor $`e^{2A(y)}`$ is commonly called a “warp factor”. The equation of motion of $`\varphi `$ is $`\varphi ^{\prime \prime }+4A^{}\varphi ^{}`$ $`=`$ $`2\kappa ^2\left({\displaystyle \frac{\mathrm{\Lambda }(\varphi )}{\varphi }}+{\displaystyle \underset{i=0}{\overset{2}{}}}{\displaystyle \frac{V_i(\varphi )}{\varphi }}\delta (yy_i)\right),`$ (6) and the Einstein equations can be written as $$A^{\prime \prime }=\frac{1}{6}\varphi ^2\frac{\kappa ^2}{3}\underset{i=0}{\overset{2}{}}V_i(\varphi )\delta (yy_i),A^2=\frac{1}{24}\varphi ^2\frac{\kappa ^2}{6}\mathrm{\Lambda }(\varphi ).$$ (7) Here a prime denotes a derivative with respect to $`y`$. The jumps corresponding to the presence of the branes are $$A^{}\|_{y_iϵ}^{y_i+ϵ}=\frac{\kappa ^2}{3}V_i(\varphi _i),\varphi ^{}\|_{y_iϵ}^{y_i+ϵ}=2\kappa ^2\frac{V_i(\varphi )}{\varphi }\|_{\varphi =\varphi _i},$$ (8) where $`\varphi _i\varphi (y_i)`$. Let $`W(\varphi )`$ be any sectionally continuous function (which we call the superpotential), with sectional functions $`W_i(\varphi )`$ defined analogous to $`\mathrm{\Lambda }_i(\varphi )`$ (see Fig. 1). Taking $`2\kappa ^2\mathrm{\Lambda }(\varphi )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{W(\varphi )}{\varphi }}\right)^2{\displaystyle \frac{1}{3}}W(\varphi )^2,`$ (9) it is possible to show that a solution to the equations, $$\varphi ^{}=\frac{W(\varphi )}{\varphi },A^{}=\frac{1}{6}W(\varphi ),$$ (10) subject to the constraints, $$W(\varphi )\|_{y_iϵ}^{y_i+ϵ}=2\kappa ^2V_i(\varphi _i),\frac{W(\varphi )}{\varphi }\|_{y_iϵ}^{y_i+ϵ}=2\kappa ^2\frac{V_i(\varphi )}{\varphi }\|_{\varphi =\varphi _i}$$ (11) is also a solution to the system of equations ($`\text{6}\text{8}`$). By solving Eq. (10) in the bulk and applying boundary conditions on the branes, we can determine the locations of the branes and hence their separation. Note that up to an arbitrary function $`_{n=2}^{\mathrm{}}\gamma _n(\varphi \varphi _i)^n`$, the brane tension is completely determined by the superpotential and the value of $`\varphi `$ on the brane, $$2\kappa ^2V_i(\varphi )=W_{i+1}(\varphi _i)W_i(\varphi _i)+\left(\frac{}{\varphi }(W_{i+1}(\varphi )W_i(\varphi ))\right)\|_{\varphi =\varphi _i}(\varphi \varphi _i).$$ (12) The solution involves fine-tuning even though it appears not to be the case. There are six constraints arising from the jump conditions on the three branes, but only five integration constants; the equations of motion and the jumps depend only upon $`A^{}(y)`$ and $`A^{\prime \prime }(y)`$ thereby rendering the value of $`A`$ on one of the branes irrelevant . III. Cosmology: General Considerations. Our starting point is the most general five-dimensional metric that preserves three-dimensional rotational and translational invariance. Thus, we adopt the cosmological principle of isotropy and homogeneity on the observable brane. However, $`y`$-dependence is maintained in the metric tensor since isotropy is broken in the fifth dimension due to the presence of branes. We consider the metric to be of the form $$ds^2=n^2(\tau ,y)d\tau ^2+a^2(\tau ,y)d𝐱^2+b^2(\tau ,y)dy^2.$$ (13) The Einstein tensor for this metric is given by , $`G_{00}`$ $`=`$ $`3\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\right){\displaystyle \frac{n^2}{b^2}}\left({\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{b^{}}{b}}\right)\right)\right\},`$ (14) $`G_{ij}`$ $`=`$ $`{\displaystyle \frac{a^2}{b^2}}\delta _{ij}\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+2{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^{}}{b}}\left({\displaystyle \frac{n^{}}{n}}+2{\displaystyle \frac{a^{}}{a}}\right)+2{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{n^{\prime \prime }}{n}}\right\}`$ (16) $`+{\displaystyle \frac{a^2}{n^2}}\delta _{ij}\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+2{\displaystyle \frac{\dot{n}}{n}}\right)2{\displaystyle \frac{\ddot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\left(2{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{n}}{n}}\right){\displaystyle \frac{\ddot{b}}{b}}\right\},`$ $`G_{05}`$ $`=`$ $`3\left({\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{\dot{a}^{}}{a}}\right),`$ (17) $`G_{55}`$ $`=`$ $`3\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^2}{n^2}}\left({\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{\dot{n}}{n}}\right)+{\displaystyle \frac{\ddot{a}}{a}}\right)\right\}.`$ (18) where a dot denotes a derivative with respect to $`\tau `$. Note that the time-independent solution of the previous section corresponds to $`a(\tau ,y)=n(\tau ,y)=e^{A(y)}`$ and $`b(\tau ,y)=1`$. We will maintain the assumption that the stabilizing potential is static, $`\dot{b}=0`$, and without loss of generality we set $`b=1`$. The energy-momentum tensor can be decomposed into two parts: a contribution from fields on the observable brane, $`\stackrel{~}{T}_{}^{A}{}_{B}{}^{}=diag(\rho _0,p_0,p_0,p_0,0)\delta (y)`$, and the contribution $`\stackrel{ˇ}{T}_{}^{A}{}_{B}{}^{}`$ of all other sources, i.e. bulk fields and matter on the other brane: $`T_{}^{A}{}_{B}{}^{}=\stackrel{~}{T}_{}^{A}{}_{B}{}^{}+\stackrel{ˇ}{T}_{}^{A}{}_{B}{}^{}`$. In the time-independent case, $`\stackrel{~}{T}_{}^{A}{}_{B}{}^{}=diag(V_0,V_0,V_0,V_0,0)\delta (y)`$. The jump conditions on the observable brane are $$\frac{[a^{}]}{a_0}=\frac{\kappa ^2}{3}\rho _0\mathrm{and}\frac{[n^{}]}{n_0}=\frac{\kappa ^2}{3}\left(3p_0+2\rho _0\right),$$ (19) where $`[a^{}]=a^{}(+ϵ)a^{}(ϵ)`$, and functions with the subscript $`0`$ are evaluated on the observable brane. On taking the jump of the (0,5) component of Einstein’s equations, these conditions lead to the energy conservation equation, $$\dot{\rho _0}+3(p_0+\rho _0)\frac{\dot{a_0}}{a_0}=0.$$ (20) which is independent of $`\dot{b}`$. Taking the jump of the (5,5) component of Einstein’s equations, we get $$3p_0\frac{a^{}}{a_0}=\rho _0\frac{n^{}}{n_0}+\frac{[\stackrel{ˇ}{T}_{55}]}{a_0},$$ (21) where $$2a^{}=a^{}(+ϵ)+a^{}(ϵ).$$ (22) We choose $`n_0=1`$, which amounts to identifying $`\tau `$ with time in conventional cosmology. Evaluating the (5,5) component of Einstein’s equations on either side of the observable brane and adding, we obtain $`\left({\displaystyle \frac{\dot{a_0}}{a_0}}\right)^2+{\displaystyle \frac{\ddot{a_0}}{a_0}}={\displaystyle \frac{\kappa ^4}{36}}\rho _0(\rho _0+3p_0){\displaystyle \frac{\kappa ^2}{3}}\stackrel{ˇ}{T}_{55}+{\displaystyle \frac{a^{}^2}{a_0^2}}(1+{\displaystyle \frac{3p_0}{\rho _0}}){\displaystyle \frac{a^{}[\stackrel{ˇ}{T}_{55}]}{\rho _0a_0^2}}.`$ (23) Usually, motivated by Hořava-Witten supergravity , a $`𝐙_\mathrm{𝟐}`$ symmetry is imposed on the solutions. We simplify the above result by requiring the solutions to Einstein’s equations to obey a “$`Z_2`$ symmetry” in the neighborhood of the observable brane. By this we simply mean $$W(\varphi (+ϵ))=W(\varphi (ϵ)),$$ (24) because it leads to the warp factor being symmetric on either side of the observable brane. By a contextual abuse of terminology, we will call this a “local $`Z_2`$ symmetry”. (Of course, in no sense are we gauging the symmetry). To create a distinction, we will reserve the bold font for the “global” $`𝐙_\mathrm{𝟐}`$ symmetry. We are therefore left with the following Friedmann-like equation: $$\left(\frac{\dot{a_0}}{a_0}\right)^2+\frac{\ddot{a_0}}{a_0}=\frac{\kappa ^4}{36}\rho _0(\rho _0+3p_0)\frac{\kappa ^2}{3}\stackrel{ˇ}{T}_{55},$$ (25) derived in . Let us emphasize the two assumptions on which our results will hinge. They are: 1. The extra dimension is assumed to be stable before studying cosmology. 2. The solutions satisfy a $`Z_2`$ symmetry in the immediate neighborhood of the observable brane, Eq. (24). In finding the static solution of the previous section, we ignored the matter energy densities on the branes by assuming that they are negligible in comparison to the brane tensions. We now include their contribution as a perturbation to the “matter-less” solution. Thus, we can study the resulting cosmology by making the ansatz $$\rho _0=\rho +V_0,p_0=pV_0,$$ (26) and Eq. (25) becomes $`\left({\displaystyle \frac{\dot{a_0}}{a_0}}\right)^2+{\displaystyle \frac{\ddot{a_0}}{a_0}}={\displaystyle \frac{\kappa ^4}{18}}V_0^2{\displaystyle \frac{1}{18}}W(\varphi _0)^2+{\displaystyle \frac{\kappa ^4}{36}}V_0(\rho 3p){\displaystyle \frac{\kappa ^4}{36}}\rho (\rho +3p).`$ (27) Here we have used $$\kappa ^2\stackrel{ˇ}{T}_{55}=\frac{1}{6}W(\varphi _0)^2,$$ (28) which is obtained by inserting (9) and (10) into (4). Notice that $`W(\varphi _0)`$ is proportional to $`A^{}(0)`$ and is therefore not well-defined. However, $`W(\varphi _0)^2`$ is well-defined. On account of the local $`Z_2`$ symmetry and Eq. (12), we have $$W(\varphi (+ϵ))^2=W(\varphi (ϵ))^2=\kappa ^4V_0^2.$$ (29) The definition of $`W(\varphi _0)^2`$ is analogous to Eq. (22), and we obtain $$W(\varphi _0)^2=\kappa ^4V_0^2.$$ (30) This result relies heavily on the existence of the local $`Z_2`$ symmetry. The first two terms on the right hand side of Eq. (27) cancel out and we are left with $$\left(\frac{\dot{a_0}}{a_0}\right)^2+\frac{\ddot{a_0}}{a_0}=\frac{\kappa ^4}{36}V_0(\rho 3p)\frac{\kappa ^4}{36}\rho (\rho +3p).$$ (31) The leading term on the right-hand side reproduces the standard cosmology if we make the identification, $`\kappa ^4V_0=6/M_{Pl}^2`$. It is essential for the observable brane to possess positive tension to arrive at the correct Friedmann equations in spite of an explicit radion stabilization. Furthermore, a specific form of the bulk energy-momentum tensor was not chosen to implement the cancellation. Beyond the conventional fine-tuning one does not need additional machinery to obtain the usual cosmology. In introducing the perturbation (26), it is no longer obvious that the solutions remain consistent. Let us consider the equation of motion of $`\varphi `$. On the observable brane it is $$\varphi ^{\prime \prime }+\left(3\frac{a^{}}{a}+\frac{n^{}}{n}\right)\varphi ^{}=2\kappa ^2\frac{}{\varphi }(\mathrm{\Lambda }(\varphi )+V_0(\varphi )),$$ (32) where all quantities are evaluated at $`y_0`$. We started with just tension on the branes and demanded that the radion stabilize the configuration of branes. Having obtained this static solution we then proceeded to consider the effect of matter on the observable brane. Let us require that the positions of the branes be unchanged by appealing to the stability of such a scenario. This is equivalent to saying that $`\varphi `$ and $`\mathrm{\Lambda }(\varphi )`$ are unchanged before and after the introduction of the matter energy density. Then consistency requires, $$\left(3\frac{a^{}}{a}+\frac{n^{}}{n}\right)\|_{0,Static}=\left(3\frac{a^{}}{a}+\frac{n^{}}{n}\right)\|_{0,Perturbed}.$$ (33) Again, using the local $`Z_2`$ symmetry and the jump conditions (19), we get $$\frac{\kappa ^2}{2}V_0\frac{\kappa ^2}{6}V_0=\frac{\kappa ^2}{2}(\rho +V_0)+\frac{\kappa ^2}{6}(3p+2\rho V_0),$$ (34) which leads to the condition for a radiation-dominated (RD) universe, $$\rho =3p.$$ (35) The interpretation of the above constraint is interesting. When matter on the observable brane is radiation, the inter-brane spacing is identical to the case when there is no matter on the brane. Conversely, when the brane location is unaffected by the matter-perturbation, the universe is RD. This observation is consistent with the fact that the radion couples to the trace of the energy-momentum tensor , which is zero for radiation. It may be possible to identify the process of radion stabilization with inflation and reheating and the time at which the inter-brane spacing becomes stable marks the end of reheating. The RD universe then ensues. To study cosmology at lower temperatures, we need the radion to be time-dependent. The radion can be written as $$\varphi (\stackrel{}{x},t,y)=\varphi (\stackrel{}{x},t)\varphi (y),$$ (36) the form of which encodes the requirement that the bulk remain static (so as to maintain a non-fluctuating Planck scale). In our previous analysis we have set $`\varphi (\stackrel{}{x},t)=1`$. The Solution Generating Technique is still applicable provided the derivatives of $`\varphi (\stackrel{}{x},t)`$ with respect to $`\stackrel{}{x}`$ and $`t`$ are negligible. It is conceivable that $`\varphi (\stackrel{}{x},t)`$ plays a role in the evolution from a RD universe to a matter-dominated universe and is perhaps responsible for the transition to an accelerating universe as in quintessence models . IV. Self-tuning Flat Branes. In this section we study the case where the superpotential takes on the form of the tree-level dilaton coupling. This will lead us to the case of a vanishing bulk cosmological constant . The resulting $`\varphi `$ has two singularities at finite distances on either side of the observable brane. We will describe how these singularities can be dealt with. Consider a superpotential of the exponential form $`W_i(\varphi )`$ $`=`$ $`\omega _ie^{\beta \varphi },`$ (37) for which $$12\kappa ^2\mathrm{\Lambda }_i(\varphi )=(3\beta ^22)\omega _i^2e^{2\beta \varphi }.$$ (38) For $`\beta ^2=2/3`$, we have the important result that $`\mathrm{\Lambda }=0`$ . Henceforth, we restrict ourselves to this choice. When we have not committed to the sign of $`\beta `$, we will leave it explicit in the equations. With $`\beta ^2=2/3`$ the branes are flat and will remain so, independent of the matter on them (hence the expression “self-tuning flat branes”). As we pointed out in Sec. II, the tension on the branes is fixed by the superpotential. As long as the tension satisfies Eq. (12), it is irrelevant what functional form it takes, i.e. the particular form of the tension we choose is simply a calculational device with no bearing on the physics. We take $`12\kappa ^2V_i(\varphi )`$ $`=`$ $`(\omega _{i+1}\omega _i)e^{\beta \varphi }.`$ (39) Then it can be shown that $$\varphi (y)=\frac{1}{\beta }\mathrm{ln}[(\underset{i=0}{\overset{2}{}}k_i\|yy_i\|+k_cy)+c],A(y)=\frac{1}{6\beta }\varphi (y)+h,$$ (40) where $`k_i=(\omega _{i+1}\omega _i)/3`$ and $`k_c=(\omega _0+\omega _3)/3`$. Here $`c`$ is a constant that can be fixed by the boundary conditions on any brane. Conventionally, the constant $`h`$ is set to zero, but its value does not affect the hierarchy of scales. Let us call the argument of the logarithm in Eq. (40), $`f(y)`$. Then, $$f(y)e^{\beta \varphi (y)}=\frac{2}{3}\omega _iy+c_i\mathrm{for}y_{i1}yy_i,$$ (41) where $`c_i`$ is a constant of integration. Note that when $`\omega _i`$ is positive (negative), $`e^{\beta \varphi }`$ falls (rises) linearly. When $`f`$ vanishes, $`\varphi `$ diverges and the warp factor vanishes. The space collapses to a point at these locations and these points can be identified with horizons. Assume that the horizons at $`y_{min}`$ and $`y_{max}`$ are such that $`y_{min}<0`$ and $`y_{max}>y_2`$. We will truncate the space at the horizons. There is much debate about the justification of this procedure if one hopes to solve the cosmological constant problem. It has been claimed that the four-dimensional cosmological constant vanishes only when the singularities contribute to the vacuum energy . If negative tension branes are introduced at the singularities, it is possible to set the four-dimensional cosmological constant to zero, but not without fine-tuning . Attempts have been made to find bulk potentials such that self-tuning remains while simultaneously removing the singularities. It has been found that if the singularities are removed, gravity is no longer localized and the four-dimensional Planck scale diverges. We will content ourselves with having $`\mathrm{\Lambda }=0`$ as an improvement to the cosmological constant problem. We assume the presence of some dynamics at the singularities that does not affect the global properties of the solution and may resolve the problem of fine-tuning. So that $`\varphi `$ be well-defined and accommodate our assumptions, we must impose the constraint, $$\underset{i=0}{\overset{2}{}}k_i>\|k_c\|\begin{array}{cc}\omega _3>0\mathrm{iff}k_c<0,\omega _0<0\mathrm{iff}k_c>0.\hfill & \end{array}$$ (42) Now we can calculate the four-dimensional Planck scale in terms of the five-dimensional Planck scale $`M_X`$, $$M_{Pl}^2=M_X^3e^{2A(y)}𝑑y$$ (43) to be $$M_{Pl}^2=M_X^3e^{2h}\left[(\frac{1}{\omega _1}\frac{1}{\omega _0})e^{\varphi _0/\beta }+(\frac{1}{\omega _2}\frac{1}{\omega _1})e^{\varphi _1/\beta }+(\frac{1}{\omega _3}\frac{1}{\omega _2})e^{\varphi _2/\beta }\right].$$ (44) We recall that the electroweak scale ($`M_{EW}`$) can be generated from the five-dimensional Planck scale $`M_X`$ via the square root of the warp factor , $$M_{EW}M_Xe^{A(0)}=M_Xe^{\varphi _0/(6\beta )}.$$ (45) Two branes geometry: Let us specialize to the case of just two branes. The formulae for the case of three branes apply by simply dropping the terms corresponding to the extra indices. We choose $$W_1(\varphi )=W_0(\varphi )=\omega _1e^{\beta \varphi },W_2(\varphi )=\omega _2e^{\beta \varphi }.$$ (46) As required by cosmology, $`W(\varphi )^2`$ has a local $`Z_2`$ symmetry about the observable brane. From Eq. (39), we need to impose $`\omega _1>0`$ so that the observable brane has positive tension. The four-dimensional Planck scale is $$M_{Pl}^2=M_X^3\left[\frac{2}{\omega _1}e^{\varphi _0/\beta }+(\frac{1}{\omega _2}\frac{1}{\omega _1})e^{\varphi _1/\beta }\right].$$ (47) From Eq. (41), one can readily see that $`e^{\varphi _0/\beta }>e^{\varphi _1/\beta }`$. Since $`\omega _0=\omega _1<0`$, the constraint from Eq. (42) requires $`\omega _2>\omega _1`$. This implies that the hidden brane has positive tension. Therefore, the second term in Eq. (47) makes a negative contribution to $`M_{Pl}^2`$. The first term must be the dominant contribution for $`M_{Pl}^2>0`$ and $`\varphi _0/\beta \varphi _1/\beta `$ must be satisfied, where “$``$” implies a hierarchy of at most two orders of magnitude. We would like the fundamental parameters $`M_X`$ and $`\omega _i`$ to be roughly of the same order of magnitude to avoid a fine-tuned hierarchy of scales. We obtain $$M_{Pl}M_Xe^{\varphi _0/(2\beta )},M_{Pl}/M_{EW}e^{\varphi _0/(3\beta )}.$$ (48) We need $`\varphi _0/(3\beta )37`$ to get the correct hierarchy with $`M_{Pl}10^{19}`$ GeV and $`M_{EW}10^3`$ GeV. This leads to $`M_X10^5`$ GeV. Interpreting $`M_X`$ as the string scale is now impossible. The difficulty arises because both the Planck and electroweak scales are determined by the value of $`\varphi `$ on the observable brane. Any self-tuning brane model with only two branes shares this problem; $`\varphi `$ will always be required to have its maximum value on the observable brane because of the local $`Z_2`$ symmetry. Three branes geometry: We construct a model with three branes in which the Planck scale will be generated by the value of $`\varphi `$ on a neighboring brane. We investigate the superpotential $$W_1(\varphi )=W_0(\varphi )=\omega _1e^{\beta \varphi },W_2(\varphi )=\omega _2e^{\beta \varphi },W_3(\varphi )=\omega _3e^{\beta \varphi }.$$ (49) Notice that $`W(\varphi )^2`$ has a local $`Z_2`$ symmetry. We leave the sign of the tensions of the hidden branes unspecified for the time being. The Planck scale is given by Eq. (44) with $`\omega _0=\omega _1`$. Again, from Eq. (41), $`e^{\varphi _0/\beta }>e^{\varphi _1/\beta }`$. The only way of getting $`e^{\varphi _2/\beta }>e^{\varphi _0/\beta }`$ is by choosing $`\omega _2<0`$. From Eq. (39) we can see that the brane at $`y_1`$ must have negative tension. Since there are no more branes in the bulk, so that $`e^{\varphi /\beta }0`$ at $`y_{max}`$, we must have $`\omega _3>0`$. Thus, the brane at $`y_2`$ has positive tension. This situation calls for two hidden branes, one with positive and the other with negative tension in a unique configuration. The configuration of branes, the profiles of $`e^{\beta \varphi }`$ and the warp factor $`e^{2A(y)}=\sqrt{e^{\beta \varphi }}`$ are shown in Fig. 2. The model resembles the “$`++`$” model of Ref. , which however is not derived by imposing constraints from radion stabilization and has a compactification on $`S^1/𝐙_\mathrm{𝟐}`$ . If we assume $`M_X`$ and $`\omega _i`$ to be of the same order of magnitude and $`\varphi _2/\beta \varphi _1/\beta ,\varphi _0/\beta `$, then $$M_{Pl}M_Xe^{\varphi _2/(2\beta )},M_{Pl}/M_{EW}e^{(\varphi _2\varphi _0/3)/(2\beta )}.$$ (50) To obtain the correct hierarchy we must have $`(\varphi _2\varphi _0/3)/(2\beta )37`$. By choosing appropriate values of $`\varphi _0`$ and $`\varphi _2`$, we are able to generate the hierarchy between $`M_{Pl}`$ and $`M_{EW}`$ for essentially any value of $`M_X`$ in between. For illustration, we present two particularly interesting examples. First, we can achieve $`M_XM_{Pl}`$ by taking $`\varphi _20`$ (or any other value that yields $`e^{\varphi _2}𝒪(1)`$), corresponding to $`\varphi _0/\beta 220`$. At the other extreme, we can obtain $`M_XM_{EW}`$ by taking $`\varphi _00`$, corresponding to $`\varphi _2/\beta 75`$. We mention in passing that the solution that leads to $`M_XM_{EW}`$ is slightly less fine-tuned in terms of the difference in $`\|\varphi _i\|`$ than the one leading to the high energy string scale. A lighter string scale is preferred in this sense. A couple of points are noteworthy. The solution presented represents the unique minimal configuration that allows for the generation of the hierarchy of scales without fine-tuning. The negative tension brane must lie between two positive tension branes. In the case at hand, it is not possible to place the negative tension brane at the fixed point of an orbifold. The only possible discrete symmetry that can be imposed on $`R^1`$ is $`𝐙_\mathrm{𝟐}`$. If we considered the orbifold $`R^1/𝐙_\mathrm{𝟐}`$ with a fixed point at $`y_1`$, we would not be able to satisfy $`e^{\varphi _2/\beta }>e^{\varphi _0/\beta }`$. Since the negative tension brane is not at an orbifold fixed point, the radion may have a problem with positivity of energy . This is an unpleasant circumstance but nevertheless, we assume the model to be theoretically feasible. More problematic is the introduction of a new hierarchy problem. By inspecting Eq. (39) it can be seen that due to the exponential dependence of the brane tensions on $`\varphi `$, a large hierarchy is generated between the values of the tensions for even moderately different values of $`\varphi `$. V. Polynomial Superpotentials. Here we consider the type of superpotential that leads to the stabilization mechanism suggested in . In it was demonstrated that a quadratic superpotential results in the mechanism of . We have studied all geometries with two positive tension branes (with bounded and unbounded $`\varphi `$) and numerically scanned the parameter space. As in the model of the previous section, we find that it is not possible to generate the appropriate scale hierarchy with only positive tension branes. We therefore study a model with two branes where the hidden brane has negative tension. The first column of Table I shows our particular choice of the polynomial superpotential, which is guided by cosmology discussions with a $`Z_2`$ symmetry. For simplicity, we have multiplied $`y`$ by $`M_X`$ to make it dimensionless. The second and third columns of Table I present the static solution to Einstein’s equations, where $`a_0`$ is an irrelevant integration constant which we set to zero. We find it necessary for the radion to be unbounded for $`y>y_1`$. This leads to $`\mathrm{\Lambda }(\varphi )`$ being unbounded below. This is often seen in $`AdS`$ supergravity and offers no threat to the model . In the region $`y<0`$, the radion may or may not be bounded without affecting the hierarchy. We will choose it to be unbounded in both regions, thus making the value of $`\varphi _0`$ a global minimum. Figure 3 illustrates the profiles of $`\varphi `$ and $`A(y)`$ in the bulk. The location of the hidden brane is $$y_1=\mathrm{ln}\left(\frac{\varphi _1}{\varphi _0}\right)^{\frac{1}{2}}.$$ (51) The electroweak scale is $$M_{EW}M_Xe^{A(0)}=M_Xe^{\frac{\varphi _0^2}{24}}.$$ (52) The Planck scale is given by $$\left(2\frac{M_{Pl}}{M_X}\right)^2=\left(\frac{\varphi _0^2}{12}\right)^{\frac{\xi }{12}}\mathrm{\Gamma }(\frac{\xi }{12},\frac{\varphi _1^2}{12},\mathrm{})+\left(\frac{\varphi _0^2}{12}\right)^{\frac{\eta }{12}}\left[\mathrm{\Gamma }(\frac{\eta }{12},\frac{\varphi _0^2}{12},\frac{\varphi _1^2}{12})+\mathrm{\Gamma }(\frac{\eta }{12},\frac{\varphi _0^2}{12},\mathrm{})\right],$$ (53) where the generalized incomplete gamma function is $`\mathrm{\Gamma }(a,x,y)_x^yt^{a1}e^t𝑑t`$. Consistency conditions imposed by positivity of the tension of the observable brane and the profile of the radion are $`\eta <\varphi _0^2<\varphi _1^2`$. When the correct hierarchy is generated, by far the dominant contribution to $`M_{Pl}`$ comes from the first term on the right-hand-side of Eq. (53). This term is the integral over the space $`y>y_1`$. The condition under which this integral dominates is $`\varphi _0^2<\varphi _1^2<\xi `$. Then $`\eta <\xi `$, and the brane at $`y_1`$ has negative tension. If we choose $`\eta >\xi `$, the brane will have positive tension, but the desired hierarchy of scales cannot be obtained. It is not possible to place the negative tension brane at an orbifold fixed point because the space beyond $`y_1`$ is crucial for generating the scale hierarchy. We can again obtain a solution with a string scale anywhere between $`M_{EW}`$ and $`M_{Pl}`$. As an explicit realization that solves the hierarchy problem, consider the following choice of parameters: $`\eta =12,\varphi _0^2=24,\varphi _1^2=100,\xi =450`$. The largest hierarchy among these parameters is only $`𝒪(10)`$. With the above choice, $`M_{Pl}10^{15}M_X`$ and $`M_{EW}10^1M_X`$. VI. Conclusion. We have studied the cosmology and hierarchy in models with branes immersed in a five-dimensional curved spacetime subject to radion stabilization. We found that when the radion field is time-independent and the inter-brane spacing is stabilized, consistent solutions that reproduce the conventional cosmological equations can naturally lead to a radiation-dominated universe. This feature is independent of the form of the stabilizing potential. The only assumption made is that the warp factor is symmetric on either side of the observable brane. Guided by constraints on the stabilizing superpotential imposed by cosmology, we proceeded to consider solutions to the hierarchy problem. We insisted that the observable brane have positive tension and considered a noncompact fifth dimension. We examined two classes of models— an exponential and a polynomial superpotential. We find that these scenarios generically require at least one hidden brane with negative tension to get the correct hierarchy. This brane cannot be located at the fixed point of an orbifold. In both models, the correct hierarchy between the electroweak and Planck scales can be obtained for any value of the string scale, including the interesting result $`M_XM_{EW}`$, without fine-tuning. The exponential superpotential leads to the interesting case of a vanishing bulk cosmological constant, referred to as a self-tuning brane model. As in , we needed to truncate the space to avoid curvature singularities. In this model, generating the hierarchy of scales results in the brane tensions becoming hierarchical. In the case of a polynomial superpotential no new hierarchy is created. Acknowledgments. We thank C. Goebel for discussions. This work was supported in part by a DOE grant No. DE-FG02-95ER40896, in part by the Wisconsin Alumni Research Foundation, and in part by the Fermi National Accelerator Laboratory, which is operated by the Universities research Association, Inc., under contract No. DE-AC02-76CHO3000.
warning/0006/hep-ph0006243.html	ar5iv	text	# HIGH ENERGY SCATTERING FROM THE 𝐴⁢𝑑⁢𝑆/𝐶⁢𝐹⁢𝑇 CORRESPONDENCEaafootnote aTalk presented at the DIS00 workshop, Liverpool, April 2000. ## 1 Introduction Historically, the interpretation of strong interactions in terms of a string theory has raised much hope $`^\mathrm{?}`$ but was a deceiving adventure. Indeed, while the Veneziano (resp. Shapiro-Virasoro) amplitudes for Reggeon (resp. Pomeron) exchanges were very promising and at the root of the open (resp. closed) string theories, problems arise when looking for the internal consistency of the whole scheme in our 4-dimensional world: a quantum anomaly requires 26 or 10 dimensions, gravitons and zero-mass vectors unavoidably appear in the spectrum of strong interaction states. So a stringy description remains an open problem for $`QCD_4.`$ Recently, the proposal of an AdS/CFT correspondence $`^\mathrm{?}`$ seems to overcome some difficulties met during the last 30 years. In very brief terms (for a extended review, see $`^\mathrm{?}`$) the idea is to unify a “microscopic” and a “microscopic” description of a configuration of $`N1`$ three-branes in the so-called Type II-B string theory in 10 dimensions. In the “microscopic” description, the system gives rise to a 4-dimensional $`SU(N)`$ gauge theory (with $`𝒩=4`$ supersymmetries), while in the “macroscopic” one, it is the source of a gravitational background equipped with a $`AdS_5S_5`$ metric with the physical Minkowski space lying at the boundary of $`AdS_5`$. A duality property is conjectured between the 4-dimensional $`SU(N)`$ gauge theory at strong coupling and the gravitational background at weak coupling. Interestingly enough, the dynamical rôle of the fifth dimension in $`AdS_5`$ is crucial for the validity of the correspondence. The case of a gauge theory with $`𝒩=4`$ supersymmetries corresponds to a 4-dimensional, non-confining, conformal field theory. The conjecture could be enlarged to confining theories without supersymmetry (e.g. see $`^\mathrm{?}`$) by introducing a “horizon” scale in the 5-th dimension. The term “horizon” comes from the consideration of a black hole metric in the bulk of $`AdS`$ space in order to break supersymmetry. Even if the exact dual of $`QCD_4`$ has not yet been identified, these dualities give a laboratory framework for gauge observables at strong coupling. For instance, the Wilson area criterion for confinement can be explicitely verified $`^\mathrm{?}`$. ## 2 High energy amplitudes: obervables and results Let us briefly outline the derivation of our papers $`^{\mathrm{?},\mathrm{?}}`$. Scattering amplitudes in the high energy limit (and small momentum transfer) can be conveniently expressed in terms of a correlator of Wilson loops $`^\mathrm{?}`$. $$A(s,q^2)=2isd^2x_{}e^{iqx_{}}\frac{W_1W_2}{W_1W_2}1$$ (1) where the two tilted Wilson loops follow elongated trajectories along the light-cone direction with transverse separation $`a`$ and a tilting angle $`\theta `$ around the impact parameter axis. This corresponds to the scattering of colorless quark-antiquark pairs of mass $`ma^1`$. Indeed, the geometrical parameters of the configuration can be related to the energy scales by analytic continuation $`\mathrm{cos}\theta \mathrm{cosh}\chi \frac{1}{\sqrt{1v^2}}=\frac{s}{2m^2}1`$ where $`\chi =\frac{1}{2}\mathrm{log}\frac{1+v}{1v}`$ is the Minkowski angle (rapidity) between the two lines, and $`v`$ is the relative velocity. The results we obtain distinguish between the large and the small impact parameter kinematics. At large impact parameter, we could use $`^\mathrm{?}`$ the supergravity approximation of the appropriate type II-B string theory, since the fields, in particular gravity, are weak. We computed the exchange contribution of all zero-mode fields between the two separated $`AdS_5`$ surfaces whose geometry is fixed by area minimization with the two initial Wilson loops at the boundary. Looking to the contribution of the various fields (dilaton, Kaluza-Klein scalars, antisymmetric tensors mixed with Ramond-Ramond forms and the graviton) we find a hierarchy of real phase-shifts $`\delta (b)\mathrm{log}\frac{W_1W_2}{W_1W_2}`$ contributing to elastic scattering at large impact parameter only. Indeed, this hierarchy is different from the static Wilson loop correlator $`^\mathrm{?}`$, since the graviton is dominant and not the Kaluza-Klein scalars. The potential problems with unitarity are avoided, since the weak field approximation appears to be valid only at very large impact parameter $`L/as^{2/7}`$ where the scattering amplitudes are purely elastic. Note, however, the retentivity of the gravitational interaction, which is still mysterious in the general context of the AdS/CFT correspondence where the decoupling from gravity is expected. In a second paper $`^\mathrm{?}`$, we addressed the problem of small impact parameter and the origin of inelasticity, i.e. imaginary contributions to the phase shift. We concentrated on a situation where the difficulty with supergravity field exchanges does not arise, since there exists a single connected minimal surface which gives the dominant contribution to the scattering amplitude in the strong coupling regime. This allows us to extend our study to small impact parameters, where inelastic channels are expected to play an important rôle. Moreover, it is possible to investigate both cases of conformal (non-confining) and confining cases by considering the appropriate geometries in $`AdS`$ spaces. Our goal was to understand the rôle of confining geometries in the characteristic features of scattering amplitudes at high energy. The main expected feature is Reggeization, i.e. the determination of the amplitudes by singularities (poles and cuts) in the complex plane of the crossed channel partial waves, moving with $`tq^2.`$ Our main result $`^\mathrm{?}`$ is that high energy amplitudes are governed by the geometry of minimal surfaces, generalizing the helicoid in different $`AdS`$ geometries with the elongated tilted Wilson loops at the boundary. Indeed, the confining geometries have the remarkable properties to admit approximately flat configurations near the horizon scale in the fifth dimension and thus the tilting angle induces (approximate) helicoidal solutions for the minimal surface problem. For this solution and after analytic continuation, one finds a Regge singularity corresponding to a linear double Regge pole trajectory with intercept one $$\alpha (t)=1+\frac{R_0^2}{4\sqrt{2g_{YM}^2N}}t,$$ (2) where $`R_0`$ is the horizon scale and $`g_{YM},`$ the gauge theory coupling. The results in confining geometries for impact parameter larger or of the order of $`R_0`$ can be contrasted with the conformal (non-confining) $`AdS_5S_5`$ case which, using an asymptotic evaluation (the mathematical knowledge on minimal surfaces embedded in $`AdS`$ spaces is yet limited!), leads to amplitudes with flat trajectories of the type $$A(s,t)is^{1+\frac{2\pi ^4}{\mathrm{\Gamma }(1/4)^4}\frac{\sqrt{2g_{YM}^2N}}{2\pi }}t^{1\frac{F(\pi /2)}{2\pi }\frac{\sqrt{2g_{YM}^2N}}{2\pi }},$$ (3) where $`F(\pi /2).3\pi `$ comes from an anomalous dimension computed in $`^\mathrm{?}`$. However, even in the confining cases, it may of course happen that the impact parameter distance between the two Wilson lines becomes much smaller than $`R_0.`$ In this case (see Fig. 1) the minimal surface problem becomes less affected by the black hole geometry and will just probe the small $`z`$ region of the geometry. The precise behaviour at these shorter distances will depend on the type of gauge theory and, in particular, on the small $`z`$ limit of the appropriate metric. In fact this limit resembles the original $`AdS_5\times S^5`$ geometry. Note that the same behaviour can be equivalently obtained through rescaling, by keeping the impact parameter fixed and putting the scale $`R_0\mathrm{}`$. The conformal behaviour of the amplitude (3) may thus give a hint of the small impact parameter limit also present in the physical confining case of $`QCD_4,`$ and thus a kind of hard-soft transition in impact parameter. ## 3 Outlook Using the $`AdS/CFT`$ correspondence, we found a relation between high-energy amplitudes in gauge theories at strong coupling and minimal surfaces generalizing the helicoid in various $`AdS`$ backgrounds. We considered three cases: (i) flat metric approximation of an $`AdS`$ black hole metric giving rise to Regge amplitudes with linear trajectories, (ii) an approximate evaluation for the conformal $`AdS_5\times S^5`$ geometry leading to flat Regge trajectories and (iii) evidence for a transition, in a confining theory, from behaviour of type (i) to (ii) when the impact parameter decreases below the interpolation scale set by the horizon radius. It would be quite useful to supplement the approximations made in our investigations by an evaluation of the string fluctuation pattern around the classical configurations we analyzed, in order to have a more precise determination of the predictions based on the $`AdS/CFT`$ correspondence. After that, we will be able to discuss the validity and usefulness of this stimulating conjecture in a deeper way. ## Acknowledgments This work was done in tight collaboration with R. Janik whose name should be associated to all the results mentioned in this contribution. This work was supported in part by the EU Fourth Framework Programme ‘Training and Mobility of Researchers’, Network ‘Quantum Chromodynamics and the Deep Structure of Elementary Particles’, contract FMRX-CT98-0194 (DG 12 - MIHT). ## References
warning/0006/hep-th0006064.html	ar5iv	text	# Theorem 1 SHEP 00-04 A gauge invariant exact renormalization group II Tim R. Morris Department of Physics, University of Southampton, Highfield, Southampton SO17 1BJ, UK Abstract A manifestly gauge invariant and regularized renormalization group flow equation is constructed for pure $`SU(N)`$ gauge theory in the large $`N`$ limit. In this way we make precise and concrete the notion of a non-perturbative gauge invariant continuum Wilsonian effective action. Manifestly gauge invariant calculations may be performed, without gauge fixing, and receive a natural interpretation in terms of fluctuating Wilson loops. Regularization is achieved by covariant higher derivatives and by embedding in a spontaneously broken $`SU(N\|N)`$ supergauge theory; the resulting heavy fermionic vectors are Pauli-Villars fields. We prove the finiteness of this method to one loop and any number of external gauge fields. A duality is uncovered that changes the sign of the squared coupling constant. As a test of the basic formalism we compute the one loop $`\beta `$ function, for the first time without any gauge fixing, and prove its universality with respect to cutoff function. hep-th/0006064 June, 2000. 1. Introduction In ref. we presented a gauge invariant Wilsonian RG (renormalization group) , formulated directly in the continuum. This formulation was shown to have many attractive features, in particular the fact that manifest gauge invariance can be maintained at all stages of the calculation and thus also in the solution for the effective action $`S`$, no gauge fixing or ghosts being required, and the equations may be reinterpreted in terms of fluctuations in the natural gauge invariant order parameters, namely Wilson loops. However, the formulation presented was not sufficient to regularise all ultra-violet divergences. In this paper we solve this problem whilst preserving all these attractive aspects . Our formulation thus furnishes for the first time a precise and concrete realisation of the notion of a non-perturbative gauge invariant continuum Wilsonian effective action . In recent years there has been substantial progress in solving supersymmetric gauge theories by computing just such an effective action, even though this object has never been defined. (Only certain general properties were required.) Whilst we concentrate here solely on pure Yang-Mills theory, we see no essential difficulty in generalising the flow equations to include fermions and scalars and indeed spacetime supersymmetry. It is clear then that our framework can underpin these ideas . The regularisation employed in ref. arises essentially from an effective cutoff function which is gauge covariantized. Similarly to gauge covariant higher derivative regularisation, this is not sufficient to regulate all ultra-violet divergences. One loop divergences slip through . In standard perturbation theory, this problem has been cured by supplementing the higher derivative regularisation with a system of Pauli-Villars regulator (PV) fields, the action being bilinear in these fields so that they provide, on integrating out, the missing one loop counterterms<sup>1</sup> and of course other finite contributions . This solution turns out to be unwieldy, but worse, here the property of being bilinear in the PV fields is not preserved by the flow: as the gauge field is integrated out higher-point PV interactions are generated. Instead, we uncovered a system of regulating fields that is more natural from the exact RG point of view, particularly so in the Wilson loop picture , as we first reported in ref. . We have gradually realised that hidden in this formulation are supermatrices and a spontaneously broken local $`SU(N\|N)`$ (in unitary gauge). We use this insight to give a concise and complete exposition of the formulation sketched in ref. . As in ref. , we will concentrate on the gauge group $`SU(N)`$ in the large $`N`$ limit. All the ideas adapt to finite $`N`$ and indeed other gauge groups, except that the embedding in the appropriate supergauge group should be formulated in such a way as to make this connection more manifest. The disadvantage of the regularisation framework reported here and in ref. , is that it was developed intuitively from the bottom up, without us being aware of the underlying local $`SU(N\|N)`$ structure. Whilst many aspects fell out correctly nevertheless, the formulation given in ref. is limited to one loop. Complete regularisation should be achieved in a manifestly local $`SU(N\|N)`$ framework, for reasons that we will outline later. The full exposition of this latter formulation is however left for a future paper. Such a framework may of course be used independently of the Wilsonian RG, and provides a novel and elegant four dimensional ‘physical’<sup>2</sup> in the usual sense that it directly suppresses higher momentum modes regularisation for gauge theory which, as we have already intimated, appears to generalise straightforwardly to (chiral) fermions, spacetime supersymmetry and so on. Subtleties in its precise definition and properties are discussed later and in the conclusions, however a full treatment is left to a later paper . One fascinating property reported here,<sup>3</sup> and paranthetically in ref. is a duality that effectively exchanges the squared coupling constant $`g^2`$ with $`g^2`$. At the moment it is not clear to us whether this duality survives in a manifestly local $`SU(N\|N)`$ framework. We comment further in the conclusions. Let us stress however that there are two main threads in this paper. On the one hand we introduce this natural gauge invariant regularisation, as described above. On the other hand, we go on to use it to repair the divergences reported in ref. , and thus develop a consistent calculational framework in which manifest gauge invariance can be maintained at all stages. The fact that the $`SU(N)`$ gauge invariance is in this way explicitly preserved, thus with no gauge fixing or BRST ghosts, results in elegant and highly constrained relations. One important consequence is that there is no wavefunction renormalization, the only quantity requiring renormalization being the coupling constant . Manifest gauge invariance is also a necessary component of our PV regularisation scheme, as we will show. Non-perturbative approaches to non-Abelian gauge theory that proceed by gauge fixing, must face up to the challenging problem of Gribov copies . Here, these problems are entirely avoided . Indeed we may turn the issue around, and use the present formulation, since it is already well defined without gauge fixing, to investigate explicitly the (quantum) consequences of gauge fixing and Gribov copies directly in the continuum. In view of the novelty of the construction presented here and in ref. , a basic test of the formalism is surely desirable. We compute concrete expressions for all the elementary vertices for two certain choices of covariantization but for general cutoff functions. We then use one of these to derive, directly in the continuum , the classical values of the two, three and four point vertices in $`S`$ and the one-loop contribution to the two-point vertex. From this we compute the one-loop $`\beta `$ function. Throughout the calculation, we use entirely general cutoff functions,<sup>4</sup> up to some basic criteria on normalisation and ultraviolet decay rates and maintain manifest gauge invariance. The fact that we obtain the result $`\beta _1=\frac{11}{3}\frac{N}{(4\pi )^2}`$, independently of the choice of cutoff functions is encouraging confirmation that the expected universality of the continuum limit has been successfully incorporated. The fact that it agrees with the usual perturbative result, demonstrates for the first time explicitly that the one-loop $`\beta `$ function is free from Gribov problems, as expected. The paper is structured as follows. In sec. 2, after preliminary definitions, we state the flow equation in superfield notation and show that it is manifestly gauge invariant, leaves the partition function invariant, and recall the property of quasilocality and that the flow corresponds to integrating out . This last feature relies on showing that the integrals are indeed ultraviolet regularised, which is established for physical one-loop vertices in sec. 9 (see below). In setting up the definitions, we also show that the present formulation is a reformulation of that sketched in ref. , and show why its form follows essentially from spontaneously broken $`SU(N\|N)`$. Full and partial supermatrix differentials, the resulting trapped $`\sigma _3`$s, covariantizations, super ‘wines’, cutoff functions $`c`$ and $`\stackrel{~}{c}`$, and the renormalization condition are all introduced here. The essence of the calculation can be followed by reading this section, sec. 3 (which demonstrates that the isolation of perturbative contributions follows in the same way as for the pure gauge case ), and secs. 7 and 8 that cover respectively the tree-level and one-loop calculation. However, to arrive at the equations in secs. 7 and 8, one needs to extract the Feynman rules from the flow equation. We do this for the reader in sec. 4. This also serves to fix the nomenclature for the Feynman rules and to explain carefully the precise relation between the flow equation, the Wilson loop diagrams, and their Feynman diagrammatic expansion. Underlying these rules are the wine vertices themselves. Although these are defined implicitly in sec. 2, see also refs. , clearly we need explicit expressions for concrete calculations. These are derived in sec. 5, for two forms of covariantization. For the form that we will use here for calculation, we also need to resolve their values at certain special momenta. In this section we also work out the wine vertices’ large momentum behaviour. This is not needed for the explicit tree and one-loop calculations, but it is needed at a very rough level to establish finiteness at one loop (again see below). We pause here to work out their large momentum behaviour in much more detail than required, because it is elegant and interesting, falls out with little effort, and may later be important. In sec. 6 we uncover and describe the symmetries of the flow equation. On the one hand this helps to understand the equation and the novelties of the underlying $`SU(N\|N)`$ regularization at a deeper level, in particular we show that charge conjugation invariance and fermion number mix to form a $`Z_4`$, that the superfields are only pseudoreal, and demonstrate the existence of a duality that in a sense exchanges $`g^2`$ with $`g^2`$. On the other hand, a precise delineation of all symmetries is needed to constrain certain ‘counterterms’ in the tree-level calculations in sec. 7. Indeed in sec. 7, we will see that the classical vertices suffer a form of divergence as a result of certain freedoms in the Pauli-Villars sector. These in turn lead to the introduction of some new parameters $`\gamma `$. We compute only those vertices that we need for the one-loop calculation in sec. 8. We also streamline the calculation a little, by borrowing some general results from sec. 9. Sec. 8 starts by showing how the $`\beta `$ function is determined, in principle non-perturbatively, following ref. . We then specialize to the concrete one-loop calculation. We explain in particular how we handle the calculation for general cutoff functions by integrating by parts so as to lower the degree of differentiation, whilst bringing all terms to a canonical algebraic form. In this way we reduce the computation to a set of boundary terms in $`D=4`$ dimensions which however depend on the power with which the cutoff functions decay. In fact as we already demonstrated in refs. , see also , the result is ambiguous due to certain total derivative terms that integrate to finite surface terms. Keeping $`D4`$ allows all us to discard all such terms, and we find that as $`D4`$ we recover the famous result for $`\beta _1`$, as already indicated. In section 9 we give a proof of finiteness for all one-loop physical vertices (i.e. with no external PV fields). We show at the end of the section where the difficulties lie in ensuring finiteness for a larger set of diagrams in this ‘unitary gauge’ formulation. The first part – up to Lemma 3 – explains in broad outline why our exact RG has these finiteness properties; it is these considerations that motivated the form of the exact RG. Although the rest of section 9 stands apart from the paper and can be skipped on first reading, they contain the reasons for some finer details in our flow equation. Finally in sec. 10, we present our conclusions, some comparisons with earlier attempts at a manifestly gauge invariant calculations, and indicate future directions. 2. SU(N$`\|`$N), Pauli-Villars regularisation, and the exact RG We formulate the approach in $`D`$ Euclidean dimensions, specializing to $`D=4`$ only when required. In ref. we defined the Pauli-Villars (PV) regulator fields as follows. We took the tensor product of the gauge group with itself, writing $`SU_1(N)\times SU_2(N)`$ to distinguish the two groups. Here we write the Hermitian generators of the two groups similarly as $`\tau _1^a`$ and $`\tau _2^a`$, where each set of generators are orthonormalised as $`\mathrm{tr}(\tau ^a\tau ^b)=\frac{1}{2}\delta ^{ab}`$. Their associated gauge fields are $`A_\mu ^1`$, the physical gauge field, and $`A_\mu ^2`$, which is unphysical and part of the regularisation scheme. (The fields are valued in the Lie algebra i.e. $`A_\mu ^iA_{}^{i}{}_{\mu }{}^{a}\tau _i^a`$, $`i=1,2`$.) We introduced a fermionic Pauli-Villars (PV) field $`(B_\mu )_{j_2}^{i_1}`$ and its complex conjugate $`(\overline{B}_\mu )_{j_1}^{i_2}`$. As indicated, $`B`$ transforms as the fundamental of $`SU_1(N)`$ and as the complex conjugate fundamental of $`SU_2(N)`$. Finally, we introduced two bosonic real scalar adjoint PV fields $`C^1`$ and $`C^2`$. The $`A^i`$’s are massless (of course) while $`B`$ and the $`C^i`$ have masses at the effective cutoff $`\mathrm{\Lambda }`$. Originally we figured out the statistics, representation content, and the interactions of these fields intuitively in such a way as to ensure finiteness of the quantum corrections in the RG flow equation. We found that there was very little freedom in the choice of interactions if this was to be achieved. The majority of the third lecture in ref. was devoted to describing the construction from this point of view, and will not be repeated here. We now understand these choices in terms of spontaneously broken $`SU(N\|N)`$. Let the supergauge field of $`SU(N\|N)`$ be $`𝒜_\mu `$, which we write in supermatrix form, i.e. as a matrix representation with bosonic diagonal elements, and fermionic off-diagonal elements: $$𝒜_\mu =\left(\begin{array}{cc}A_\mu ^1& B_\mu \\ \overline{B}_\mu & A_\mu ^2\end{array}\right).$$ Since $`𝒜`$ is valued in the graded Lie algebra of $`SU(N\|N)`$, we require that its supertrace vanishes: $`\mathrm{str}𝒜=0`$, where the supertrace is defined by $$\mathrm{str}𝐗=\mathrm{str}\left(\begin{array}{cc}X^{11}& X^{12}\\ X^{21}& X^{22}\end{array}\right)=\mathrm{tr}X^{11}\mathrm{tr}X^{22}$$ for any supermatrix $`𝐗`$. (The extra signs incurred on commutation mean that for such matrices it is only the supertrace that is cyclically symmetric, and thus in particular $`SU(N\|N)`$ invariant.) We also introduce the superscalar $$𝒞=\left(\begin{array}{cc}C^1& D\\ \overline{D}& C^2\end{array}\right),$$ but do not require $`\mathrm{str}𝒞=0`$. Note that the conditions $`\mathrm{tr}C^1=\mathrm{tr}C^2=0`$ are thus not imposed as in ref. . Similarly, the conditions $`\mathrm{tr}A_\mu ^1=\mathrm{tr}A_\mu ^2=0`$ are not directly imposed. Only the constraint $`\mathrm{tr}A_\mu ^1\mathrm{tr}A_\mu ^2=0`$ has been applied so far. Actually we can, and here will, require that the $`U(1)`$ components of the supermatrix algebra, $`\mathrm{tr}(A_\mu ^1+A_\mu ^2)`$ and $`\mathrm{tr}(C^1+C^2)`$, are also absent, by a suitable modification of the matrix commutator representation of the super-Lie product in $`SU(N\|N)`$ . (Although it is not required for the rest of the paper, it may help to pause on these points. Bars’ observation is that the super-Lie product may be equally well represented by a bilinear (anti)symmetric “$``$” bracket $$[,]_\pm ^{}=[,]_\pm \frac{1}{2N}\mathrm{tr}[,]_\pm ,$$ where $`[,]_\pm `$ is applied to the supergenerators and is a commutator or anticommutator as appropriate. This effectively removes the unit matrix as a representative of the algebra, and thus its bosonic $`U(1)`$ subgroup. In fact we will handle these $`U(1)`$ components in a different way in our proof that spontaneously broken $`SU(N\|N)`$ acts as a regulator . For the present paper it is only important to note that these modifications are subleading in $`N`$ and thus vanish in the $`N=\mathrm{}`$ limit that we will take shortly. Let us stress that there are in any case no $`U(1)`$ components in the present formulation . In this paper, we use $`SU(N\|N)`$ to interpret and systemize our earlier results , and in this interpretation we may understand the missing $`U(1)`$ factors in the way that we have just described.) We choose the Lagrangian so that $`𝒞`$ picks up the expectation value $`<𝒞>\mathrm{\Lambda }\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$, breaking $`SU(N\|N)`$ spontaneously to $`SU(N)\times SU(N)`$. In unitary gauge the Goldstone modes $`D`$ vanish (eaten by the $`B_\mu `$’s), leaving the massive ‘Higgs’ $`C^i`$ and massive vector fermions $`B_\mu `$. In this way, we elegantly recover exactly the spectrum of fields introduced in ref. . Remarkably, with covariant higher derivative regularisation, we also recover, up to some small details, the $`\widehat{S}`$ interactions used in ref. for the one-loop $`\beta `$ function calculation! Note that the only remaining massless fields are the gauge fields of $`SU(N)\times SU(N)`$, which are of course charge neutral under each others gauge group. Thus at energies much less than $`\mathrm{\Lambda }`$ we are left only with these gauge fields which decouple into the required physical $`SU(N)`$ Yang-Mills theory, and a copy which as we will see has the opposite sign squared coupling constant and is thus unphysical. Not surprisingly, the $`SU(N\|N)`$ theory described above has very good ultra-violet behaviour. Technically this arises in the unbroken theory, because any quantum correction involving $`\mathrm{tr}\mathrm{\hspace{0.17em}1}=N`$ in the $`SU(N)`$ theory, here involves $`\mathrm{str}\mathrm{\hspace{0.17em}1}=0`$. (We will refer to this as ‘the supertrace mechanism’; of course at the level of component fields it arises through exact cancellation between bosonic and fermionic degrees of freedom.) Indeed in the large $`N`$ limit the symmetric phase of the supertheory thus has no quantum corrections at all. With covariant higher derivative regularisation, we expect to be able to ensure that the remaining corrections even at $`N\mathrm{}`$ are finite. However, the full development of these investigations is left for the future. Here, and from now on, we will take the $`N=\mathrm{}`$ limit and work solely with the theory described in ref. , which is closely related to spontaneously broken $`SU(N\|N)`$ in unitary gauge as described above (and further outlined below). As we will show, it appears that we cannot ensure the full regularisation of the theory this way except to one loop with external gauge fields. We interpret this as symptoms of the differences described above and of the expected poor ultraviolet behaviour of the unitary gauge. It is helpful to introduce the $`2N\times 2N`$ ‘Pauli’ matrices $$\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\text{and}\text{ }\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$ (in terms of which for example we have $`\mathrm{str}𝐗=\mathrm{tr}\sigma _3𝐗`$) and to separate the bosonic and fermionic parts of the superfields, thus $$𝐀_\mu =\left(\begin{array}{cc}A_\mu ^1& 0\\ 0& A_\mu ^2\end{array}\right),𝐁_\mu =\left(\begin{array}{cc}0& B_\mu \\ \overline{B}_\mu & 0\end{array}\right),𝐂=\left(\begin{array}{cc}C^1& 0\\ 0& C^2\end{array}\right),\omega =\left(\begin{array}{cc}\omega ^1& 0\\ 0& \omega ^2\end{array}\right),$$ where $`\omega `$ is the unbroken $`SU(N)\times SU(N)`$ gauge transformation. Defining now $$𝐃_\mu =_\mu i𝐀_\mu \text{and}\text{ }_\mu =_\mu i𝒜_\mu ,$$ both $`𝐀`$, and $`𝒜=𝐀+𝐁`$, gauge transform under $`SU(N)\times SU(N)`$: $$\delta 𝐀_\mu =𝐃_\mu \omega :=[𝐃_\mu ,\omega ]\text{,}\text{ }\delta 𝒜_\mu =_\mu \omega [_\mu ,\omega ].$$ The formalism of ref. can be lifted to the supertheory as follows. Functional derivatives are defined with respect to the supertrace, thus $$\begin{array}{cc}\hfill \frac{\delta }{\delta 𝐀_\mu }& =\left(\begin{array}{cc}\delta /\delta A_\mu ^1& 0\\ 0& \delta /\delta A_\mu ^2\end{array}\right),\frac{\delta }{\delta 𝐁_\mu }=\left(\begin{array}{cc}0& \delta /\delta \overline{B}_\mu \\ \delta /\delta B_\mu & 0\end{array}\right)\hfill \\ \hfill \frac{\delta }{\delta 𝐂}& =\left(\begin{array}{cc}\delta /\delta C^1& 0\\ 0& \delta /\delta C^2\end{array}\right),\frac{\delta }{\delta 𝒜}=\frac{\delta }{\delta 𝐀}+\frac{\delta }{\delta 𝐁},\hfill \end{array}$$ (where the adjoint derivatives may be defined by $`\delta /\delta A^i=2\tau _i^a\delta /\delta A_{}^{i}{}_{}{}^{a}`$, $`i=1,2`$, as in ref. . At large $`N`$, $`\delta /\delta C^i`$ will also be thought of as an adjoint derivative, i.e. we take the $`C^i`$ to be traceless. Again this restriction makes no difference in the large $`N`$ limit, but it is convenient for the $`C^i`$ terms to inherit in this way the conventions associated with the $`\tau _i^a`$ being normalised to $`1/2`$.) These transform homogeneously; their other important properties are easier to state by ignoring the $`x`$ dependence and $`\mu `$ index (as in ref. ). Let $`𝐌`$, $`𝐘`$ and $`𝐙`$ also be supermatrix representations as defined above. If $`s(𝐌)`$ is a (bosonic) function of $`𝐌`$ such that $$\delta s(𝐌)=\mathrm{str}\delta 𝐌𝐘,$$ then, in precise analogy with ref. , $$\frac{s}{𝐌}=𝐘.$$ (Here in the adjoint parts we use the completeness relation for $`SU(N)`$, which leads to the same relation as the unrestricted functional derivative, up to $`1/N`$ corrections that ensure the coefficient matrix is projected onto its traceless part . These corrections can be neglected since, as stated above, we are only interested from here on, in the $`N=\mathrm{}`$ limit.) This leads to ‘supersowing’ under the supertrace $$\mathrm{str}𝐗\frac{s}{𝐌}=\mathrm{str}\mathrm{𝐗𝐘},$$ and to ‘supersplitting’ $$\delta 𝐗=𝐘\delta 𝐌𝐙\mathrm{str}\frac{}{𝐌}𝐗=\mathrm{str}𝐘\mathrm{str}𝐙.$$ If however $`𝐌`$ is only block diagonal like $`𝐀`$ or $`𝐂`$, equivalently purely bosonic, then $`𝐘`$ in (2.1) and (2.1) is replaced by $`\mathrm{d}_+𝐘`$, the block diagonal (a.k.a. bosonic) part of $`𝐘`$. On the other hand, if $`𝐌`$ is only block off-diagonal, like $`𝐁`$, equivalently purely fermionic, then $`𝐘`$ in (2.1)–(2.1) is replaced by $`\mathrm{d}_{}𝐘`$, the block off-diagonal (a.k.a. fermionic) part of $`𝐘`$. The projectors $`\mathrm{d}_\pm `$ may be expressed as: $$\mathrm{d}_\pm 𝐘=\frac{1}{2}\left(𝐘\pm 𝐘^{}\right),$$ where $`𝐘^{}=\sigma _3𝐘\sigma _3`$ has opposite sign fermionic components compared to $`𝐘`$. Either from this or directly, we have the identities $$\mathrm{str}𝐗\mathrm{d}_\pm 𝐘=\mathrm{str}(\mathrm{d}_\pm 𝐗)𝐘=\mathrm{str}\mathrm{d}_\pm 𝐗\mathrm{d}_\pm 𝐘.$$ Summarising, for a ‘partial’ supermatrix $`𝐌=\mathrm{d}_\pm 𝐌`$, (2.1) is altered to $$\mathrm{str}𝐗\frac{s}{𝐌}=\mathrm{str}𝐗\mathrm{d}_\pm 𝐘.$$ These alterations and the relations (2.1) are obvious in the bosonic / fermionic language. Similarly, for a partial supermatrix $`𝐌=\mathrm{d}_\pm 𝐌`$, the splitting relation (2.1) is altered: $$\mathrm{str}\frac{}{𝐌}𝐗=\frac{1}{2}\left(\mathrm{str}𝐘\mathrm{str}𝐙\pm \mathrm{str}\sigma _3𝐘\mathrm{str}\sigma _3𝐙\right),$$ as follows most readily by writing $`\delta 𝐗=𝐘\mathrm{d}_\pm (\delta 𝐌)𝐙`$, after which the ‘full’ supersplitting relation (2.1) may be used. This ‘broken’ supersplitting relation (which in a local $`SU(N\|N)`$ formulation would be seen to arise from insertions of $`<𝒞>\mathrm{\Lambda }\sigma _3`$) provides the only reason why there are any quantum corrections at all in the large $`N`$ limit. As discussed in ref. , the quantum corrections split open the traces, and in the large $`N`$ limit, they survive only in the terms where all fields vacate one of the two traces so as to leave $`\mathrm{tr}\mathrm{\hspace{0.17em}1}=N`$. In the present case only the broken terms survive, through $`\mathrm{str}\sigma _3=2N`$; the unbroken quantum corrections have no way to survive because vacating a supertrace leaves behind $`\mathrm{str}\mathrm{\hspace{0.17em}1}`$ which vanishes. Expanding (2.1) in $`N\times N`$ block components, in the bosonic case we have $`\mathrm{tr}Y^{11}\mathrm{tr}Z^{11}+\mathrm{tr}Y^{22}\mathrm{tr}Z^{22}`$, and in the fermionic case $`\mathrm{tr}Y^{11}\mathrm{tr}Z^{22}\mathrm{tr}Y^{22}\mathrm{tr}Z^{11}`$. Note the minus signs in the latter, as expected for a fermionic loop. Recall that a central construction of the gauge invariant flow equation in ref. , is the ‘wine’, the covariantization of a (smooth) momentum space kernel $`W_pW(p^2/\mathrm{\Lambda }^2)`$. In position space we write the kernel as $$W_{xy}\frac{d^Dp}{(2\pi )^D}W(p^2/\mathrm{\Lambda }^2)\mathrm{e}^{ip.(xy)}.$$ For two $`N\overline{N}`$ representations of the gauge group $`SU(N)`$, $`v(y)`$ and $`u(x)`$, we write $$\begin{array}{c}\text{ }u\{W\}v=(2.1)\text{ }\hfill \\ \text{ }\underset{m,n=0}{\overset{\mathrm{}}{}}d^Dxd^Dyd^Dx_1\mathrm{}d^Dx_nd^Dy_1\mathrm{}d^Dy_mW_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(x_1,\mathrm{},x_n;y_1,\mathrm{},y_m;x,y)\text{ }\hfill \\ \text{ }\mathrm{tr}\left[u(x)A_{\mu _1}(x_1)\mathrm{}A_{\mu _n}(x_n)v(y)A_{\nu _1}(y_1)\mathrm{}A_{\nu _m}(y_m)\right],\text{ }\hfill \end{array}$$ where without loss of generality we may insist that $`\{W\}`$ satisfies $`u\{W\}vv\{W\}u`$. This equation defines the wine vertices $`W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}`$, given the method of covariantization . As in , we write the $`m=0`$ vertices (where there is no second product of gauge fields), more compactly as $$W_{\mu _1\mathrm{}\mu _n}(x_1,\mathrm{},x_n;x,y)W_{\mu _1\mathrm{}\mu _n,}(x_1,\mathrm{},x_n;;x,y),$$ while the $`m=n=0`$ term is just the original kernel (2.1), i.e. $$W_,(;;x,y)W_{xy}.$$ For two supermatrix representations $`𝐯(y)`$ and $`𝐮(x)`$ we define in precise analogy, the supergauge invariant $$\begin{array}{c}\text{ }𝐮\{W\}𝐯=(2.2)\text{ }\hfill \\ \text{ }\underset{m,n=0}{\overset{\mathrm{}}{}}d^Dxd^Dyd^Dx_1\mathrm{}d^Dx_nd^Dy_1\mathrm{}d^Dy_mW_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(x_1,\mathrm{},x_n;y_1,\mathrm{},y_m;x,y)\text{ }\hfill \\ \text{ }\mathrm{str}\left[𝐮(x)𝒜_{\mu _1}(x_1)\mathrm{}𝒜_{\mu _n}(x_n)𝐯(y)𝒜_{\nu _1}(y_1)\mathrm{}𝒜_{\nu _m}(y_m)\right].\text{ }\hfill \end{array}$$ Note that the superwine’s vertices $`W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}`$ are the same as those in (2.1). Equation (2.2) corresponds precisely to the ‘peppering’ prescription in ref. ! We restrict, as in ref. , to covariantizations which traverse back and forth along a coincident Wilson line, i.e. such that they may be represented as $$𝐮\{W\}𝐯=d^Dxd^Dy𝒟𝒞_{xy}\mathrm{str}𝐮(x)\mathrm{\Phi }[𝒞_{xy}]𝐯(y)\mathrm{\Phi }^1[𝒞_{xy}].$$ Here the measure $`𝒟𝒞_{xy}`$ over curves $`𝒞_{xy}`$ from $`x`$ to $`y`$, is normalised by $$𝒟𝒞_{xy}1=W_{xy},$$ and encodes our choice of covariantization. The (super) Wilson lines are defined by the path ordered exponential $$\begin{array}{cc}\hfill \mathrm{\Phi }[𝒞_{xy}]& =P\mathrm{exp}i_{𝒞_{xy}}𝑑z^\mu 𝒜_\mu (z),\hfill \\ & =1i_0^1𝑑\tau \dot{z}^\mu 𝒜_\mu (z)_0^1𝑑\tau _2_0^{\tau _2}𝑑\tau _1\dot{z}𝒜(\tau _1)\dot{z}𝒜(\tau _2)+\mathrm{},\hfill \end{array}$$ where we have parametrized $`𝒞_{xy}`$ by $`z^\mu (\tau )`$, $`\tau [0,1]`$, $`z(0)=x`$, $`z(1)=y`$. Since this definition is independent of the parametrization of the path, the same is true of the measure $`𝒟𝒞_{xy}`$, without loss of generality. We also choose the measure to be Lorentz covariant, satisfy the exchange symmetry below (2.1), and to be smooth in momentum space, i.e. to yield vertices that are Taylor expandable to all orders in momenta . One example that satisfies all these criteria is to utilise the momentum representation to write : $$𝐮\{W\}𝐯=\mathrm{str}d^Dx𝐮(x)W(^2/\mathrm{\Lambda }^2)𝐯(x).$$ This is the covariantization that we will use in this paper to calculate the one-loop $`\beta `$ function. \[We will see later that it corresponds to coincident lines (2.3).\] Another example is simply to use two straight super Wilson lines : $$𝐮\{W\}𝐯=d^Dxd^DyW_{xy}\mathrm{str}𝐮(x)\mathrm{\Phi }[ł_{xy}]𝐯(y)\mathrm{\Phi }^1[ł_{xy}],$$ $`ł_{xy}`$ being the straight line between $`x`$ and $`y`$. We will compute concrete formulae for the wine vertices in both covariantizations. Introducing the superfield strength $`_{\mu \nu }=i[_\mu ,_\nu ]`$, we write the ‘seed’ action as $$\widehat{S}=\frac{1}{2}_{\mu \nu }\{c^1\}_{\mu \nu }+\mathrm{\Lambda }^2_\mu \{\stackrel{~}{c}^1\}_\mu +\sigma \mathrm{\Lambda }^4𝐂\{1\}𝐂.$$ The first term, which we will refer to as $`\widehat{S}_A`$, is simply the $`\widehat{S}`$ of ref. translated to superfields (corresponding again to peppering ), in particular $`c(p^2/\mathrm{\Lambda }^2)>0`$ is a smooth cutoff function satisfying $`c(0)=1`$ (in fact without loss of generality) and $`c(x)0`$ as $`x\mathrm{}`$.<sup>5</sup> As in refs. , $`x`$ used as a generic argument for these functions, should not be confused with the position $`x`$ and the position space kernel $`c_{xy}`$, defined as in (2.1). In addition we have two further terms for $`𝐁`$ and $`𝐂`$, where $$_\mu =𝐁_\mu +_\mu 𝐂,𝐂\{1\}𝐂\mathrm{str}d^Dx𝐂^2(x),$$ and $`\stackrel{~}{c}(p^2/\mathrm{\Lambda }^2)>0`$ is another smooth cutoff profile whose properties are recalled below. (2.3) is nothing but the seed action of ref. , translated into this language. Note however that such terms are expected in spontaneously broken $`SU(N\|N)`$; thus $`\widehat{S}_A`$ is just the higher-derivative regularised pure gauge part, in particular the wrong sign $`A_\mu ^2`$ action is now seen to be a consequence of the superfield (in particular supertrace) structure; the second term, which we name $`\widehat{S}_B`$, collects the remaining kinetic pieces after expanding around $`<𝒞>`$ and imposing unitary gauge,<sup>6</sup> apart from one significant difference described below in particular providing the mass term for $`𝐁`$; the last term, $`\widehat{S}_C`$, is the mass term for the ‘Higgs’. (Here $`\sigma >0`$ a free parameter. Note that the $`𝐂`$ field we discussed earlier appears here as $`\mathrm{\Lambda }𝐂`$. This change of variables is required for the finiteness of the flow equation, as explained in sec. 9.) Recall that the point of $`𝐂`$ was to cancel the longitudinal divergences from $`𝐁_\mu `$ . In the $`SU(N\|N)`$ language this cancellation is unsurprising because the longitudinal part of $`𝐁_\mu `$ is nothing but $`𝐃`$, the eaten fermionic partner of $`𝐂`$. Actually for $`\widehat{S}_B`$ to correspond exactly to spontaneously broken $`SU(N\|N)`$, the $`𝐁`$-$`𝐂`$ cross-terms should contain an insertion of $`i\sigma _3`$ arising from an explicit $`<𝒞>`$. This in turn leads to some differences in the discrete symmetries here, as explained in sec. 5. Also, the self-interactions of the Higgs are missing: these are not needed for the regularisation (to one loop with only external gauge fields). We require $`\stackrel{~}{c}(x)0`$ as $`x\mathrm{}`$, to regularise $`𝐂`$ propagation, while for the $`\stackrel{~}{c}^1`$ term not to disturb the high energy behaviour of the transverse part of $`𝐁_\mu `$ we clearly require $$c/(x\stackrel{~}{c})0\text{as}\text{ }x\mathrm{}.$$ (We require $`\stackrel{~}{c}_0\stackrel{~}{c}(0)>0`$ in order for it to act as a mass for $`𝐁_\mu `$ and positive kinetic term for $`𝐂`$ but we do not require $`\stackrel{~}{c}_0=1`$.) More precise requirements on the UV asymptotics of $`c`$ and $`\stackrel{~}{c}`$ will be needed and are derived in sec 9. We introduce three new kernels via $$K(x)=\frac{d}{dx}\left(\frac{x\stackrel{~}{c}c}{x\stackrel{~}{c}+c}\right),xL(x)=\frac{d}{dx}\left(\frac{x^2\stackrel{~}{c}^2}{x\stackrel{~}{c}+c}\right)\mathrm{and}xM(x)=\frac{d}{dx}\left(\frac{x^2\stackrel{~}{c}}{x+\sigma \stackrel{~}{c}}\right)$$ ($`cc(x)`$, $`\stackrel{~}{c}\stackrel{~}{c}(x)`$ here). Finally we can write our full exact RG equation: $$\begin{array}{ccc}\hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}S& [𝐀,𝐁,𝐂]=\frac{1}{2\mathrm{\Lambda }^2}\left(\frac{1}{N}\frac{\delta }{\delta 𝒜_\mu }\frac{\delta S}{\delta 𝒜_\mu }\right)\{c^{}\}\frac{\delta \mathrm{\Sigma }_g}{\delta 𝒜_\mu }+\frac{1}{2\mathrm{\Lambda }^4}\left(\frac{ł}{N}łS\right)\{L\}ł\mathrm{\Sigma }_g\hfill & \\ & +\frac{1}{2\mathrm{\Lambda }^2}\left(\frac{1}{N}\frac{\delta }{\delta 𝐁_\mu }\frac{\delta S}{\delta 𝐁_\mu }\right)\{Kc^{}\}_𝐀\frac{\delta \mathrm{\Sigma }_g}{\delta 𝐁_\mu }+\frac{1}{2\mathrm{\Lambda }^4}\left(\frac{1}{N}\frac{\delta }{\delta 𝐂}\frac{\delta S}{\delta 𝐂}\right)\{ML\}\frac{\delta \mathrm{\Sigma }_g}{\delta 𝐂},\hfill & \\ \hfill \mathrm{where}& ł=\frac{\delta }{\delta 𝐂}+_\mu \frac{\delta }{\delta 𝒜_\mu }\mathrm{and}\mathrm{\Sigma }_g=g^2S2\widehat{S}.\hfill & (2.3)\hfill \end{array}$$ In here, the first term on the RHS is just our pure gauge field flow equation with $`A_\mu `$ replaced by $`𝒜_\mu `$ and trace by supertrace. Finiteness considerations motivated the precise form of the other terms (see sec. 9 and ref. .) By the wine $`\{Kc^{}\}_𝐀`$ we mean that in (2.2) only $`𝐀`$ is used, rather than the full peppered $`𝒜`$. In $`ł`$ we apply the $`\delta /\delta 𝒜_\mu `$ first and then $`_\mu `$, i.e. we understand this $`\delta /\delta 𝒜_\mu `$ as not differentiating the $`𝒜_\mu `$ in the $`_\mu `$ of $`_\mu \delta /\delta 𝒜_\mu `$. Note that as in ref. , by prime (as in $`c^{}`$) we mean differentiation with respect to its argument (here $`p^2/\mathrm{\Lambda }^2`$). The position space representation (2.1), the covariantization, and the resulting vertices (2.1), are all labelled by the underlying kernel \[i.e. in this case replacing the letter $`W`$ by $`c^{}`$ throughout (2.1) – (2.2)\]. (2.3) yields precisely the Feynman rules chosen earlier for the $`\beta `$ function calculation . The coupling $`g`$ is defined as in ref. via the field strength for $`A_\mu ^1`$ in the $`A^1`$ part of the action: $$S\|_{A^1}=\frac{1}{2g^2}\mathrm{tr}d^Dx\left(F_{\mu \nu }^1\right)^2+O(^3)$$ (discarding the vacuum energy). Note that we can and do impose, as in ref., the requirement of ‘quasilocality’, and thus in particular that $`S`$ has a derivative expansion to all orders. (N.B. Other than of course its supersymmetry, this is the crucial fundamental property assumed of the Wilsonian effective action in supersymmetric theories, which in turn justifies its holomorphy .) Clearly, the flow equation (2.3) is manifestly $`SU(N)\times SU(N)`$ gauge invariant. It also leaves the partition function $$𝒵=𝒟[𝐀,𝐁,𝐂]\mathrm{e}^{^{NS}}$$ invariant. To see this, note that (2.3) implies that $$\begin{array}{ccc}\hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}\mathrm{e}^{^{NS}}& =\frac{1}{2\mathrm{\Lambda }^2}\frac{\delta }{\delta 𝒜_\mu }\{c^{}\}\left(\frac{\delta \mathrm{\Sigma }_g}{\delta 𝒜_\mu }\mathrm{e}^{^{NS}}\right)\frac{ł}{2\mathrm{\Lambda }^4}\{L\}\left(\mathrm{e}^{^{NS}}ł\mathrm{\Sigma }_g\right)\hfill & (2.4)\hfill \\ & \frac{1}{2\mathrm{\Lambda }^2}\frac{\delta }{\delta 𝐁_\mu }\{Kc^{}\}_𝐀\left(\frac{\delta \mathrm{\Sigma }_g}{\delta 𝐁_\mu }\mathrm{e}^{^{NS}}\right)\frac{1}{2\mathrm{\Lambda }^4}\frac{\delta }{\delta 𝐂}\{ML\}\left(\frac{\delta \mathrm{\Sigma }_g}{\delta 𝐂}\mathrm{e}^{^{NS}}\right),\hfill & \end{array}$$ and hence is a total functional derivative. Actually, for this to be true we need the $`\delta /\delta 𝒜_\mu `$ in the leftmost $`ł`$ to act on everything in the expression and thus also on the $`𝒜_\mu `$ in $`_\mu `$, in apparent contradiction with the definition given for (2.3). However, the difference between the two definitions gives $`\frac{\delta }{\delta 𝒜_\mu }𝒜_\mu `$, which vanishes by the supertrace mechanism. \[To see this set $`𝐘=1`$ in (2.1).\] We can see indirectly that the exact RG equation (2.3) corresponds to integrating out by the arguments already given in ref. . Similarly from , we still have that $`S`$ may be expanded in traces and products of traces. We can write these as products of supertraces of the fields $`𝐀`$, $`𝐁`$ and $`𝐂`$, and if necessary from (2.1) or (2.1), embedded $`\sigma _3`$’s. Representing the supertraces by closed loops we have the same diagrammatic notation for the full RG equations as before: $$\mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}\text{}=g^2\text{}+2\text{}+\frac{g^2}{N}\text{}\frac{2}{N}\text{}$$ Fig.1. Diagrammatic representation of the flow equation. A circumflex in a circle indicates $`\widehat{S}`$. where here we have taken the (naïve ) large $`N`$ limit. Once again this implies that at most a single (super)trace survives in the effective action. The expansion of the loops in powers of the fields yields Feynman diagrams as before . Note that the closed loops may no longer necessarily be interpreted as integrals over pure gauge (super) Wilson loops because $`𝐂`$s, $`\sigma _3`$s and individual $`𝐁`$s may be inserted. However, the contributions consist of pure gauge (viz. $`𝒜`$) sections joining isolated $`𝐂`$s, $`\sigma _3`$s and $`𝐁`$s, and thus may be interpreted as integrals over pure gauge Wilson lines joining a countable number of such points. If we represent diagrammatically the insertion of a $`\sigma _3`$ into (a join in) such a superloop, by a filled arrow pointing to the insertion point, the supersowing special cases (2.1), and the supersplitting special cases (2.1), can be seen to be caused by the same local process, as shown in fig. 2. $$\frac{1}{2}\text{}\pm \frac{1}{2}\text{}$$ Fig.2. Insertion of $`\sigma _3`$s occurs in the same way, whether or not the attaching wine results in a tree or loop correction. Contributions to (2.3) containing terms of form of fig. 3 are apparently required by (2.4) (where the wine attaches at the top to $`S`$ or $`\widehat{S}`$). However, in the large $`N`$ limit these contributions vanish. The reason is as follows. Firstly the $`\{Kc^{}\}_𝐀`$, or $`\{ML\}`$, wines obviously cannot bite their own tails since they do not contain $`𝐁`$, or $`𝐂`$, respectively. Secondly the other end of the wine in fig. 3 must attach to a vertex with at least two points, because all one-point vertices vanish by charge conjugation invariance (see sec. 6). Since only one superloop can contain fields in the large $`N`$ limit, the closed superloop in fig. 3 must thus be field free. But for the two remaining possibilities $`\{c^{}\}`$ and $`\{L\}`$, this loop is formed by a ‘full’ supermatrix differential $`\delta /\delta 𝒜`$ which thus yields $`\mathrm{str}\mathrm{\hspace{0.17em}1}=0`$ by (2.1). Fig.3. A wine biting its own tail. 3. Perturbative expansion We recall the general structure of the perturbative expansion . As in ref. we keep $`D`$ general at this stage. It will be helpful to access $`D=4`$ via the limit $`D4`$ . It will also be helpful to write (2.3) as $$\mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}S=a_0[S,g^2S2\widehat{S}]+a_1[g^2S2\widehat{S}],$$ where we have expanded $`\mathrm{\Sigma }_g`$ and written the classical terms as the bilinear functional $`a_0`$ and the quantum terms as the linear functional $`a_1`$. Since this has the same form as the pure gauge case, the isolation of perturbative contributions from these equations proceeds as before . Thus we see from (3.1), that $`S1/g^2`$ at the classical level \[consistent with (2.4)\], and by iteration, using (3.1), that $`S`$ has as expected the weak coupling expansion $$S=\frac{1}{g^2}S_0+S_1+g^2S_2+\mathrm{}.$$ Substituting this expansion in (3.1) and recalling that $`g`$ will run at the quantum level, we see that the $`\beta `$ function must also take the standard form $$\beta :=\mathrm{\Lambda }\frac{g}{\mathrm{\Lambda }}=\beta _1g^3+\beta _2g^5+\mathrm{}.$$ By (2.4), $`g^2`$ has dimension $`4D`$. The $`\beta _i`$ thus carry dimensions $`(D4)i`$ and (as we will confirm) are not universal except when $`D=4`$. From (3.1) and (3.1), we obtain the loopwise expansion of (3.1): $$\begin{array}{ccc}\hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}S_0& =a_0[S_0,S_02\widehat{S}]\hfill & (3.1)\hfill \\ \hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}S_1& =2\beta _1S_02a_0[S_0\widehat{S},S_1]+a_1[S_02\widehat{S}]\hfill & (3.2)\hfill \\ \hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}S_2& =2\beta _2S_02a_0[S_0\widehat{S},S_2]a_0[S_1,S_1]+a_1[S_1],\hfill & (3.3)\hfill \end{array}$$ etc. . In a similar way we obtain equations for the weak coupling expansion of integrated operators which in this notation are identical to those of ref. . 4. Feynman rules This section sets out the nomenclature we will use for the Feynman rules and the resulting expressions as derived by expanding (2.3) in a power series in the fields. Along the way we explain, pedantically, the precise method for translating fig. 1 to equations for individual vertices. As before , this is diagrammatically represented by replacing the thick black lines of fig. 1, which represent the exact expressions without expansion in the fields, by thin lines, and a series of contributions with increasing numbers of points (blobs). The points represent individual fields and appear in all places on the composite Wilson loop with equal weight. Where necessary the points will now be labelled by the flavour of the field that attaches to it. Specializing to a single supertrace as appropriate for the large $`N`$ limit, the field expansion of the effective action, which is illustrated diagrammatically in fig. 4, takes the form: $$S=\frac{1}{s_n}d^Dx_1\mathrm{}d^Dx_nS_{a_1\mathrm{}a_n}^{\sigma ^j𝐗_1\mathrm{}𝐗_n}(x_1,\mathrm{},x_n)\mathrm{str}\sigma _3^j𝐗_1^{a_1}(x_1)\mathrm{}𝐗_n^{a_n}(x_n),$$ where the superfields $`𝐗_i^{a_i}`$ are $`𝐀^{\mu _i}`$, $`𝐁^{\mu _i}`$ or $`𝐂`$, the indices $`a_i=\mu _i`$, or null for $`𝐂`$, and $`j=0`$ or $`1`$. (We will omit the $`\sigma `$ superscript for $`j=0`$ and write it without exponent when $`j=1`$.) Only one cyclic ordering of each list $`𝐗_1\mathrm{}𝐗_n`$ appears in the sum. Furthermore, if the list $`𝐗_1\mathrm{}𝐗_n`$ is invariant under some nontrivial cyclic permutations, then $`s_n`$ is the order of the cyclic subgroup, otherwise $`s_n=1`$. (See sec. 6.1.) In the process of computing $`S`$, there can arise any number of trapped $`\sigma _3`$s via fig. 2, however in the field expansion these can always be reduced to at most one by (anti)commutation through ($`𝐁`$), $`𝐀`$ or $`𝐂`$, and if one remains it will be placed at the beginning of the supertrace (thus resulting in a marked superloop as illustrated). Fig.4. Expansion of the action into supertraces of fields, with and without a $`\sigma _3`$ insertion. Concentrating on a given vertex, but with $`j`$ so far undetermined, and turning to the composite Wilson loops on the RHS of fig. 1, we assign the vertex’s flavours ($`𝐀`$, $`𝐁`$ or $`𝐂`$) together with their associated momenta and if appropriate Lorentz indices, to the points, summing over all cyclic permutations. For each resulting configuration, the flavours of the component wine and action vertices can then be determined. The position of any embedded $`\sigma _3`$s in component action vertices then follow. Of course some cases already vanish at this stage due to the absence of the appropriate vertex, e.g. when a $`𝐂`$-point is placed on a wine. In other cases there may be more than one choice for the component vertices, in which case the corresponding Feynman diagrams are summed over. Attachments made via partial supermatrices are expanded as in fig. 2. The full set of $`\sigma _3`$s can be combined inside the Wilson loops, (anti)commuting past external fields as necessary, and eliminated via $`\sigma _3^2=1`$ and/or $`\mathrm{str}\sigma _3=2N`$, leaving at most one $`\sigma _3`$, which is moved to its canonical position as determined by (4.1). (Actually, the expansion step fig. 2, and all $`\sigma _3`$s can be ignored for classical vertices as explained in sec. 7.) Finally, applying momentum conservation and including loop momentum integrals if appropriate, we read off according to the Feynman rules set out below, the complete expressions for the component vertices, and multiply the whole by $`1/2\mathrm{\Lambda }^2`$ . Fig.5. Expansion of the wine in gauge fields. The wine expansion (2.2) appears as fig. 5. The wine component of the vertices is then given by (2.2) (with $`𝒜`$ replaced by $`𝐀`$ in the case of $`\{Kc^{}\}_𝐀`$) and thus in general, from (2.3) – (2.3), as : $$\begin{array}{c}\text{ }W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)(2\pi )^D\delta (\underset{i=1}{\overset{n}{}}p_i+\underset{j=1}{\overset{m}{}}q_j+r+s)=(4.1)\text{ }\hfill \\ \text{ }(i)^{n+m}d^Dud^Dv𝒟𝒞_{uv}_u^v𝑑x_n^{\mu _n}_u^{x_n}𝑑x_{n1}^{\mu _{n1}}\mathrm{}_u^{x_2}𝑑x_1^{\mu _1}_v^u𝑑y_m^{\nu _m}_v^{y_m}𝑑y_{m1}^{\nu _{m1}}\mathrm{}_v^{y_2}𝑑y_1^{\nu _1}\text{ }\hfill \\ \text{ }\mathrm{exp}i\left(r.u+s.v+\underset{i}{}p_i.x_i+\underset{j}{}q_j.y_j\right),\text{ }\hfill \end{array}$$ where the $`x_i`$ integration is along the curve $`𝒞_{uv}`$, and the $`y_j`$ integration along the same curve but in the opposite direction cf. fig. 6. (As in ref. , all momenta are taken to be pointing in to the vertex.) Explicit expressions up to $`O(𝒜^2)`$, for two covariantizations, are given in sec. 5. Fig.6. Feynman rule for the wine, with momentum and Lorentz labels. From (2.3), we now have several extra pieces to take into account in the full Feynman rule for the wine vertices compared to ref. : extra factors of $`1/\mathrm{\Lambda }^2`$, contributions from $`ł`$, the flavour and (if it exists) Lorentz index of the functional derivative the wine vertex actually attaches to. (Due to $`ł`$, it is no longer the case that one just contracts these Lorentz indices.) A complete notation for the wine vertex Feynman rule can therefore be given as: $$V_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m,a_1a_2}^{𝐗_1\mathrm{}𝐗_n,𝐘_1\mathrm{}𝐘_m,𝐙_1𝐙_2}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s),$$ with momenta and indices as represented in fig. 6, where the $`𝐗_i^{\mu _i}(p_i)`$ and $`𝐘_j^{\nu _j}(q_j)`$ are $`𝐀`$s or $`𝐁`$s, and the functional derivatives are with respect to $`𝐙_1^{a_1}(r)`$ and $`𝐙_2^{a_2}(s)`$. The index $`a_1`$ ($`a_2`$) is the Lorentz index $`\alpha `$ ($`\beta `$) if $`𝐙_1`$ ($`𝐙_2`$) is $`𝐀`$ or $`𝐁`$, and null if it is $`𝐂`$. This notation is convenient for listing the Feynman rules below. They may all be expressed immediately in terms of the expansion of $`L`$ with its extra pieces. For this we write \[with the momentum arguments $`(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)`$ here suppressed\]: $$W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m,\alpha \beta }=\delta _{\alpha \beta }W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}+L_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m,\alpha \beta },$$ where shorthands (2.2) and (2.2) apply (see also (2.1) and above), and thus in particular at the zero-point level $`W_{p,\alpha \beta }=\delta _{\alpha \beta }W_p+L_{p,\alpha \beta }`$. The first few $`𝐂`$-$`𝐀_\beta `$ and $`𝐀_\alpha `$-$`𝐀_\beta `$ vertices extracted from $`ł\{L\}ł`$ are: $$\begin{array}{cc}\hfill L_{p,\beta }& L_{,,\beta }(;;p,p)=ip_\beta L_p/\mathrm{\Lambda }^2\hfill \\ \hfill L_{p,\alpha \beta }& =p_\alpha p_\beta L_p/\mathrm{\Lambda }^2\hfill \\ \hfill L_{\mu ,\beta }(p;q,r)& =i\left\{r_\beta L_\mu (p;q,r)+\delta _{\mu \beta }L_q\right\}/\mathrm{\Lambda }^2\hfill \\ \hfill L_{\mu ,\alpha \beta }(p;q,r)& =\left\{r_\beta \delta _{\mu \alpha }L_rq_\alpha \delta _{\mu \beta }L_qq_\alpha r_\beta L_\mu (p;q,r)\right\}/\mathrm{\Lambda }^2\hfill \\ \hfill L_{\mu \nu ,\beta }(p,q;r,s)& =i\left\{s_\beta L_{\mu \nu }(p,q;r,s)+\delta _{\beta \nu }L_\mu (p;r,s+q)\right\}/\mathrm{\Lambda }^2\hfill \\ \hfill L_{\mu \nu ,\alpha \beta }(p,q;r,s)& =\{\delta _{\mu \alpha }\delta _{\nu \beta }L_{p+r}+s_\beta \delta _{\mu \alpha }L_\nu (q;p+r,s)\hfill \\ & r_\alpha \delta _{\nu \beta }L_\mu (p;r,s+q)r_\alpha s_\beta L_{\mu \nu }(p,q;r,s)\}/\mathrm{\Lambda }^2,\hfill \end{array}$$ and via the coincident line identities (cf. sec. 5 and ref. ): $$\begin{array}{cc}\hfill L_{\mu ,\nu ,\beta }(p;q;r,s)& =L_{\mu \nu ,\beta }(p,q;r,s)L_{\nu \mu ,\beta }(q,p;r,s)\hfill \\ \hfill L_{\mu ,\nu ,\alpha \beta }(p;q;r,s)& =L_{\mu \nu ,\alpha \beta }(p,q;r,s)L_{\nu \mu ,\alpha \beta }(q,p;r,s).\hfill \end{array}$$ From (2.3), we thus obtain for the zero-point wine Feynman rules (4.2): $$\begin{array}{cc}\hfill V_{,,\mu \nu }^{,,\mathrm{𝐀𝐀}}(;;p,p)=c_{p,\alpha \beta }^{}& ,V_{,,\alpha \beta }^{,,\mathrm{𝐁𝐁}}(;;p,p)=K_{p,\alpha \beta },\hfill \\ \hfill V_{,,}^{,,\mathrm{𝐂𝐂}}(;;p,p)=M_p/\mathrm{\Lambda }^2& ,V_{,,\beta }^{,,\mathrm{𝐂𝐀}}(;;p,p)=L_{p,\beta },\hfill \end{array}$$ and for the one-point wine Feynman rules, using (2.2): $$\begin{array}{cc}\hfill V_{\mu ,\alpha \beta }^{𝐀,\mathrm{𝐀𝐀}}(p;q,r)& =V_{\mu ,\alpha \beta }^{𝐁,\mathrm{𝐀𝐁}}(\mathrm{"})=V_{\mu ,\alpha \beta }^{𝐁,\mathrm{𝐁𝐀}}(\mathrm{"})=c_{\mu ,\alpha \beta }^{}(p;q,r)\hfill \\ \hfill V_{\mu ,\alpha \beta }^{𝐀,\mathrm{𝐁𝐁}}(p;q,r)& =K_{\mu ,\alpha \beta }(p;q,r)\hfill \\ \hfill V_{\mu ,}^{𝐀,\mathrm{𝐂𝐂}}(p;q,r)& =\frac{1}{\mathrm{\Lambda }^2}M_\mu (p;q,r)\hfill \\ \hfill V_{\mu ,\beta }^{𝐀,\mathrm{𝐂𝐀}}(p;q,r)& =V_{\mu ,\beta }^{𝐁,\mathrm{𝐂𝐁}}(\mathrm{"})=L_{\mu ,\beta }(p;q,r)\hfill \\ \hfill V_{\mu ,\alpha }^{𝐀,\mathrm{𝐀𝐂}}(p;q,r)& =V_{\mu ,\alpha }^{𝐁,\mathrm{𝐁𝐂}}(\mathrm{"})=L_{\mu ,\alpha }(p;r,q).\hfill \end{array}$$ Of course we do not list those that vanish due to the absence of an appropriate vertex. Similarly for the two-point wine Feynman rules: $$\begin{array}{cc}\hfill V_{\mu \nu ,\alpha \beta }^{\mathrm{𝐀𝐀},\mathrm{𝐀𝐀}}(p,q;r,s)& =V_{\mu \nu ,\alpha \beta }^{\mathrm{𝐁𝐀},\mathrm{𝐁𝐀}}(\mathrm{"})=V_{\mu \nu ,\alpha \beta }^{\mathrm{𝐁𝐁},\mathrm{𝐀𝐀}}(\mathrm{"})=V_{\mu \nu ,\alpha \beta }^{\mathrm{𝐀𝐁},\mathrm{𝐀𝐁}}(\mathrm{"})=c_{\mu \nu ,\alpha \beta }^{}(p,q;r,s)\hfill \\ \hfill V_{\mu \nu ,\alpha \beta }^{\mathrm{𝐀𝐀},\mathrm{𝐁𝐁}}(p,q;r,s)& =K_{\mu \nu ,\alpha \beta }(p,q;r,s)\hfill \\ \hfill V_{\mu ,\nu ,\alpha \beta }^{𝐀,𝐀,\mathrm{𝐀𝐀}}(p;q;r,s)& =V_{\mu ,\nu ,\alpha \beta }^{𝐁,𝐀,\mathrm{𝐀𝐁}}(\mathrm{"})=V_{\mu ,\nu ,\alpha \beta }^{𝐀,𝐁,\mathrm{𝐀𝐁}}(\mathrm{"})=c_{\mu ,\nu ,\alpha \beta }^{}(p;q;r,s)\hfill \\ \hfill V_{\mu \nu ,\beta }^{\mathrm{𝐀𝐀},\mathrm{𝐂𝐀}}(p,q;r,s)& =L_{\mu \nu ,\beta }(p,q;r,s)\hfill \\ \hfill V_{\mu \nu ,\alpha }^{\mathrm{𝐀𝐀},\mathrm{𝐀𝐂}}(p,q;r,s)& =L_{\nu \mu ,\alpha }(q,p;s,r)\hfill \\ \hfill V_{\mu ,\nu ,\beta }^{𝐀,𝐀,\mathrm{𝐂𝐀}}(p;q;r,s)& =L_{\mu ,\nu ,\beta }(p;q;r,s)\hfill \\ \hfill V_{\mu \nu ,}^{\mathrm{𝐀𝐀},\mathrm{𝐂𝐂}}(p,q;r,s)& =\frac{1}{\mathrm{\Lambda }^2}M_{\mu \nu }(p,q;r,s).\hfill \end{array}$$ Here also we only list those that we will need later. The expansion of $`\widehat{S}_A`$ in (2.3) is pure $`𝒜`$ and we reserve unlabelled $`\widehat{S}`$ vertices for this: $$\widehat{S}_A=\underset{n=2}{\overset{\mathrm{}}{}}\frac{1}{n}d^Dx_1\mathrm{}d^Dx_n\widehat{S}_{\mu _1\mathrm{}\mu _n}(x_1,\mathrm{},x_n)\mathrm{str}𝒜_{\mu _1}(x_1)\mathrm{}𝒜_{\mu _n}(x_n).$$ This way the same explicit expressions as in refs. apply: $$\begin{array}{cc}\hfill \widehat{S}_{\mu \nu }(p)& \widehat{S}_{\mu \nu }(p,p)=2\mathrm{\Delta }_{\mu \nu }(p)/c_p\hfill \\ \hfill \widehat{S}_{\mu \nu \lambda }(p,q,r)& =\frac{2}{c_p}(p_\lambda \delta _{\mu \nu }p_\nu \delta _{\lambda \mu })+2c_\nu ^1(q;p,r)(p_\lambda r_\mu p.r\delta _{\lambda \mu })+\mathrm{cycles}\hfill \\ \hfill \widehat{S}_{\mu \nu \lambda \sigma }(p,q,r,s)& =\frac{1}{c_{p+q}}(\delta _{\sigma \mu }\delta _{\lambda \nu }\delta _{\lambda \mu }\delta _{\nu \sigma })+2c_\nu ^1(q;p,r+s)(p_\sigma \delta _{\lambda \mu }p_\lambda \delta _{\sigma \mu })\hfill \\ & +2c_\sigma ^1(s;p,r+q)(p_\nu \delta _{\mu \lambda }p_\lambda \delta _{\mu \nu })+2c_{\nu \lambda }^1(q,r;p,s)(p_\sigma s_\mu p.s\delta _{\sigma \mu })\hfill \\ & +c_{\nu ,\sigma }^1(q;s;p,r)(p_\lambda r_\mu p.r\delta _{\lambda \mu })+\mathrm{cycles}\hfill \end{array}$$ etc. , where in the two-point vertex we set $`p_1=p_2=p`$, and introduce the transverse combination $`\mathrm{\Delta }_{\mu \nu }(p):=p^2\delta _{\mu \nu }p_\mu p_\nu `$, in the three-point vertex we add the two cyclic permutations of $`(p_\mu ,q_\nu ,r_\lambda )`$, and in the four-point vertex the three cyclic permutations of $`(p_\mu ,q_\nu ,r_\lambda ,s_\sigma )`$. Fig.7. Feynman rules for the seed vertices. The full ‘seed’ vertices that we will need, are then given as follows, cf. fig. 7. Two-point vertices (there are no one-point vertices): $$\begin{array}{cc}\hfill \widehat{S}_{\mu \nu }^{\mathrm{𝐀𝐀}}(p)& =\widehat{S}_{\mu \nu }(p)\hfill \\ \hfill \widehat{S}_{\mu \nu }^{\mathrm{𝐁𝐁}}(p)& =\widehat{S}_{\mu \nu }(p)+2\mathrm{\Lambda }^2\delta _{\mu \nu }/\stackrel{~}{c}_p\hfill \\ \hfill \widehat{S}^{\mathrm{𝐂𝐂}}(p)& =2\mathrm{\Lambda }^2p^2/\stackrel{~}{c}_p+2\sigma \mathrm{\Lambda }^4,\hfill \end{array}$$ three-point vertices: $$\begin{array}{cc}\hfill \widehat{S}_{\mu \nu \lambda }^{\mathrm{𝐀𝐀𝐀}}(p,q,r)& =\widehat{S}_{\mu \nu \lambda }(p,q,r)\hfill \\ \hfill \widehat{S}_\lambda ^{\mathrm{𝐂𝐂𝐀}}(p,q,r)& =2\mathrm{\Lambda }^2\{p_\lambda \stackrel{~}{c}_p^1q_\lambda \stackrel{~}{c}_q^1p.q\stackrel{~}{c}_\lambda ^1(r;q,p)\}\hfill \\ \hfill \widehat{S}_{\mu \nu \lambda }^{\mathrm{𝐁𝐁𝐀}}(p,q,r)& =\widehat{S}_{\mu \nu \lambda }(p,q,r)+2\mathrm{\Lambda }^2\delta _{\mu \nu }\stackrel{~}{c}_\lambda ^1(r;q,p)\hfill \\ \hfill \widehat{S}_{\mu \nu }^{\mathrm{𝐁𝐁𝐂}}(p,q,r)& =2i\mathrm{\Lambda }^2\left\{r_\nu \stackrel{~}{c}_\mu ^1(p;r,q)+r_\mu \stackrel{~}{c}_\nu ^1(q;p,r)\right\},\hfill \end{array}$$ and four-point vertices: $$\begin{array}{cc}\hfill \widehat{S}_{\mu \nu \lambda \sigma }^{\mathrm{𝐀𝐀𝐀𝐀}}(p,q,r,s)& =\widehat{S}_{\mu \nu \lambda \sigma }(p,q,r,s)\hfill \\ \hfill \widehat{S}_{\mu \nu \lambda \sigma }^{\mathrm{𝐁𝐁𝐀𝐀}}(p,q,r,s)& =\widehat{S}_{\mu \nu \lambda \sigma }(p,q,r,s)+2\mathrm{\Lambda }^2\delta _{\mu \nu }\stackrel{~}{c}_{\lambda \sigma }^1(r,s;q,p)\hfill \\ \hfill \widehat{S}_{\lambda \sigma }^{\mathrm{𝐂𝐂𝐀𝐀}}(\mathrm{"})=2\mathrm{\Lambda }^2& \left\{\delta _{\lambda \sigma }\stackrel{~}{c}_{p+s}^1+p_\lambda \stackrel{~}{c}_\sigma ^1(s;q+r,p)q_\sigma \stackrel{~}{c}_\lambda ^1(r;q,p+s)p.q\stackrel{~}{c}_\lambda ^1(r,s;q,p)\right\}\hfill \end{array}$$ 5. Explicit expressions for wine vertices The wine vertices are defined implicitly by the choice of covariantization, and (2.1) or equivalently (2.2). In this section we derive formulae for the covariantizations (2.3) and (2.3), although in this paper we will only use (2.3) for concrete calculations. 5.1. Straight Wilson Lines In the straight line case (2.3), the integration over curves in (4.1), is trivial, being replaced by its normalisation (2.3), whilst $`𝒞_{uv}`$ is just the straight line which we may parametrize as $`x_i=u+(vu)t_i`$ and $`y_i=v+(uv)t_i^{}`$. Translation invariance implies that the integrand depends on $`u`$ and $`v`$ only in the combination $`vu`$. We replace $`vu`$ with derivatives with respect to the its conjugate momentum $$p=s+\underset{i=1}{\overset{n}{}}p_it_i+\underset{j=1}{\overset{m}{}}q_j(1t_j^{}).$$ (This expression drops out by substitution in the exponential of (4.1) and is unique up to the use of overall momentum conservation.) Integrating over the overall position yields the overall momentum conserving $`\delta `$ function. Cancelling this from both sides of (4.1) we arrive at $$\begin{array}{c}\text{ }W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)=(5.1)\text{ }\hfill \\ \text{ }(1)^m_0^1𝑑t_n_0^{t_n}𝑑t_{n1}\mathrm{}_0^{t_2}𝑑t_1_0^1𝑑t_m^{}_0^{t_m^{}}𝑑t_{m1}^{}\mathrm{}_0^{t_2^{}}𝑑t_1^{}\frac{}{p^{\mu _1}}\mathrm{}\frac{}{p^{\mu _n}}\frac{}{p^{\nu _1}}\mathrm{}\frac{}{p^{\nu _m}}W_p.\text{ }\hfill \end{array}$$ As an example, the one-point vertex is given by $$W^\mu (p;r,s)=\frac{2}{\mathrm{\Lambda }^2}_0^1𝑑tW_{tp+s}^{}(tp+s)^\mu .$$ We note that these straight line vertices are completely symmetric on all Lorentz indices. 5.2. Covariantization via (2.3) We first explain why this covariantization is of the coincident line form (2.3), equivalently (4.1). In expanding (2.3), or preferably (since the supermatrices are not necessary here) $$u\{W\}v=\mathrm{tr}d^Dxu(x)W(D^2/\mathrm{\Lambda }^2)v(x),$$ in powers of the gauge field, we only need write $$\begin{array}{c}\text{ }u\{W\}v=\underset{m,n=0}{\overset{\mathrm{}}{}}d^Dxd^Dyd^Dx_1\mathrm{}d^Dx_nW_{\mu _1\mathrm{}\mu _n}(x_1,\mathrm{},x_n;x,y)\text{ }\hfill \\ \text{ }\mathrm{tr}\left[u(x)A_{\mu _1}(x_1)\mathrm{}A_{\mu _n}(x_n)v(y)\right],\text{ }\hfill \end{array}$$ where these vertices are the coefficients of the ordered product of gauge fields $`A_{\mu _1}(x_1)\mathrm{}`$ $`A_{\mu _n}(x_n)`$ as before, but moreover, each of these gauge fields is understood to act by commutation on the expression to its right. Expanding the commutators and relabelling the gauge fields that appear on the right hand side of $`v(y)`$ in terms of $`A_{\nu _i}(y_i)`$, as in (2.1), the other vertices are then given in terms of them by $$W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)=()^m\underset{interleaves}{}W_{\lambda _1\mathrm{}\lambda _{m+n}}(k_1,\mathrm{},k_{m+n};r,s),$$ where we have (trivially) transferred to momentum space, and the sum runs over all interleaves of the sequences $`p_1^{\mu _1},\mathrm{},p_n^{\mu _n}`$ and $`q_m^{\nu _m},\mathrm{},q_1^{\nu _1}`$ i.e. combined sequences $`k_1^{\lambda _1},\mathrm{},k_{m+n}^{\lambda _{m+n}}`$ in which the $`p^\mu `$s remain ordered with respect to each other, and similarly the $`q^\nu `$s remain in reverse order. These are in fact the “coincident line” identities of ref. . Now we employ a single-Wilson-line representation for the $`W`$s on the RHS of (5.2), i.e. (4.1) with $`m=0`$. (This follows from a path integral representation of the kernel in (5.2), expressing the $`A`$s in the adjoint representation.) By isolating the integrals over the coordinates conjugate to the $`q_j`$s in the right hand terms of (5.2), we readily see that they collect together, and on reversing their direction, form the $`y_j`$ integrals of the coincident line representation (4.1). Now we compute the contribution to $`W_{\mu _1\mathrm{}\mu _n}(p_1,\mathrm{},p_n;r,s)`$ from the $`(D^2)^m`$ term in the Taylor expansion of $`W(D^2/\mathrm{\Lambda }^2)`$. (This expansion must exist by quasilocality .) In this paper we only need explicit expressions to order $`A^2`$, but we derive a formula for the general vertex below, in order to analyse its large momentum behaviour. Expanding, $$\begin{array}{c}\text{ }(D^2)^m=(^2)^m+\underset{\alpha =0}{\overset{m1}{}}(^2)^\alpha (i.A+iA.+A^2)(^2)^{m1\alpha }(5.2)\text{ }\hfill \\ \text{ }\underset{\genfrac{}{}{0pt}{}{\alpha ,\beta ,\gamma 0}{\alpha +\beta +\gamma =m2}}{}(^2)^\alpha (.A+A.)(^2)^\beta (.A+A.)(^2)^\gamma +O(A^3).\text{ }\hfill \end{array}$$ We already have noted in (2.2) that the zeroth order in $`A`$ gives $$W(;r,s)=W_s,$$ and this is trivially confirmed by (5.2). Transforming the $`O(A)`$ terms to momentum space, we see that (5.2) supplies a contribution to $`W_\mu (p;r,s)`$ of the form $$(p+2s)_\mu \underset{\alpha =0}{\overset{m1}{}}r^{2\alpha }s^{2m22\alpha }=(rs)_\mu \frac{r^{2m}s^{2m}}{p.(rs)}$$ (where we have used momentum conservation) and thus resumming the Taylor expansion, $$W_\mu (p;r,s)=(rs)_\mu \frac{W_rW_s}{p.(rs)}.$$ Similarly at $`O(A^2)`$, (5.2) supplies a contribution to $`W_{\mu \nu }(p,q;r,s)`$ of form $$\delta _{\mu \nu }\underset{\alpha =0}{\overset{m1}{}}r^{2\alpha }s^{2m22\alpha }+(p+2s+2q)_\mu (q+2s)_\nu \underset{\genfrac{}{}{0pt}{}{\alpha ,\beta ,\gamma 0}{\alpha +\beta +\gamma =m2}}{}r^{2\alpha }(s+q)^{2\beta }s^{2\gamma }.$$ The latter sum may readily be evaluated e.g. by noting that it is equal to $$\frac{1}{2\pi i}\frac{dz}{z^{m1}}\frac{1}{(1r^2z)(1[s+q]^2z)(1s^2z)}$$ for a contour of infinitessimal radius encircling the origin (which we close onto the other poles), and thus after resumming the expansion of $`W`$, $$\begin{array}{c}\text{ }W_{\mu \nu }(p,q;r,s)=\delta _{\mu \nu }\frac{W_sW_r}{s^2r^2}(5.3)\text{ }\hfill \\ \text{ }(p+2r)_\mu (q+2s)_\nu \left\{\frac{W_{s+q}}{q.(q+2s)p.(p+2r)}+\frac{1}{s^2r^2}\left[\frac{W_r}{p.(p+2r)}\frac{W_s}{q.(q+2s)}\right]\right\}\text{ }\hfill \end{array}$$ The case where the gauge fields appear on either side of the wine is then given in terms of this by (5.2): $$W_{\mu ,\nu }(p;q;r,s)=W_{\mu \nu }(p,q;r,s)W_{\nu \mu }(q,p;r,s).$$ 5.3. Special momenta and covariantization (2.3) The formulae (5.3) and (5.3) are ambiguous at certain special momenta. One way to determine the correct value is simply to return to first principles and resum (5.3) and (5.3) at the special values. However it is comforting to find that these results also appear uniquely by recalling that these formulae are only valid when the total momentum flowing into the vertex vanishes (i.e. is conserved), and taking the limit as the special configuration is approached. Thus (5.3) needs care at the point $`p=0`$, i.e. $`r=s`$, however by momentum conservation we can replace $`p.(rs)`$ by $`s^2r^2`$ after which the limit $`rs`$ is trivial in (5.3), giving $$W_\mu (0;s,s)=\frac{2s_\mu }{\mathrm{\Lambda }^2}W_s^{}.$$ (This also follows from resumming (5.3) with $`p=0`$ and $`r=s`$, or from the gauge transformation relations cf. (6.2) or ref. : $$p^\mu W_\mu (p;sp,s)=W_{s+p}W_s$$ by expanding to first order in $`p`$ after which $`p^\mu `$ may be removed uniquely from both sides. The longitudinal parts of higher powers in $`p`$ are also determined uniquely by gauge invariance. See appendix A of ref. .) Note that the case $`r=s`$ ($`p=2s`$) in (5.3) trivially gives zero \[by a limit or from (5.3)\]. (5.3) needs care at the point $`p=q`$ corresponding to $`r=s`$: $$\begin{array}{ccc}\hfill W_{\mu \nu }(p,p;r,r)& =\underset{ϵ0}{lim}W_{\mu \nu }(p,pϵ;r+ϵ,r)\hfill & (5.4)\hfill \\ & =\frac{\delta _{\mu \nu }}{\mathrm{\Lambda }^2}W_r^{}+(p+2r)_\mu (p+2r)_\nu \left\{\frac{W_{r+p}W_r}{[p.(p+2r)]^2}\frac{W_r^{}}{\mathrm{\Lambda }^2p.(p+2r)}\right\}.\hfill & \end{array}$$ There are other special momentum configurations that need careful definition in (5.3), but we will need only (5.4) here. 5.4. The general-point vertex in covariantization (2.3) The order $`A^n`$ term in $`(D^2)^m`$ supplies a contribution to $`W_{\mu _1\mathrm{}\mu _n}(p_1,\mathrm{},p_n;r,s)`$ which is a sum over $`a`$ insertions of $`A^2`$ and $`b`$ insertions of $`.A+A.`$ and over all permutations of these factors, cf. (5.2). In momentum space this reads $$\underset{\genfrac{}{}{0pt}{}{a,b}{2a+b=n}}{}\underset{\mathrm{perms}}{}T_1^{i_1}\mathrm{}T_{a+b}^{i_{a+b}}\underset{\genfrac{}{}{0pt}{}{\alpha _0,\mathrm{},\alpha _{a+b}}{{\scriptscriptstyle \alpha _j}=mab}}{}\underset{k=0}{\overset{a+b}{}}P_{I_k}^{2\alpha _k},$$ where the $`i_k=`$ 1 or 2, according to whether a $`.A+A.`$ or $`A^2`$ term is taken from $`D^2`$ respectively. They yield respectively the tensors $$T_k^1=p_{I_k}^{\mu _{I_k}}+2P_{I_k}^{\mu _{I_k}}\text{and}\text{ }T_k^2=\delta ^{\mu _{I_k1}\mu _{I_k}}.$$ $`I_k=_{j=1}^ki_j`$ keeps track of the number of gauge fields accounted for, with $`I_0:=0`$ and $`I_{a+b}=n`$, and the total momentum flow (directed from $`s`$) into the $`k^{\mathrm{th}}`$ insertion is given by $$P_{I_k}=s+\underset{j=I_k+1}{\overset{n}{}}p_j.$$ These definitions are illustrated in fig. 8. Note that $`P_{I_0}=P_0=r`$ and $`P_{I_{a+b}}=P_n=s`$. Fig.8. Tensor structure and momentum flow in the general vertex. As in (5.3), the RH sum in (5.5) may be evaluated readily by writing first $$\underset{\genfrac{}{}{0pt}{}{a,b}{2a+b=n}}{}\underset{\mathrm{perms}}{}T_1^{i_1}\mathrm{}T_{a+b}^{i_{a+b}}\frac{1}{2\pi i}\frac{dz}{z^{m+1ab}}\underset{k=0}{\overset{a+b}{}}\frac{1}{1zP_{I_k}^2}$$ for a contour of infinitessimal radius encircling the origin. Closing on to the other poles and resumming the Taylor expansion of $`W`$, we thus obtain $$W_{\mu _1\mathrm{}\mu _n}(p_1,\mathrm{},p_n;r,s)=\underset{\genfrac{}{}{0pt}{}{a,b}{2a+b=n}}{}\underset{\mathrm{perms}}{}T_1^{i_1}\mathrm{}T_{a+b}^{i_{a+b}}\underset{j=0}{\overset{a+b}{}}W_{P_{I_j}}\underset{\genfrac{}{}{0pt}{}{k=0}{kj}}{\overset{a+b}{}}\frac{1}{P_{I_j}^2P_{I_k}^2}.$$ 5.5. Large momentum behaviour of covariantization (2.3) We determine the behaviour of such a vertex when a given momentum flow through the vertex becomes large. This is needed to study the finiteness properties of the exact RG and is used for the proof in sec. 9. We have already derived in appendix A of ref. , lower bounds on the divergences, by gauge invariance considerations. Simply by power counting, we can now obtain from (5.5), upper bounds on these divergences. In some cases, where no cancellations can occur or where the lower and upper bound agree, we can thus be confident that the large momentum behaviour is then precisely known. Furthermore, we can readily furnish simple expressions for the coefficients of the leading large momentum behaviour. It is helpful in this section to adopt the following convention. By $``$ we will mean equal up to corrections that decay relative to the stated term as the large momentum $`k\mathrm{}`$; by $``$ we mean also that the stated term is determined only up to a coefficient of $`O(k^0)`$. In fact this level of precision is much more than we need here, since the proof in sec. 9 only relies on the following properties: when the large momentum flow (the loop momentum $`k`$) is from end to end ($`sk`$, $`rk`$, $`k\mathrm{}`$) the behaviour of the $`n`$ point vertex is no worse than that of the zero point vertex $`W_k`$, and when the large momentum instead enters via the side (one $`p_j\pm k`$) the ‘covariantized differentiated propagators’ $`\{c^{}\}`$, $`\{K\}`$, $`\{L\}`$, $`\{M\}`$ do not diverge. Consider first the case where the large momentum flow is from end to end ($`sk`$, $`rk`$, $`k\mathrm{}`$, all the $`p_j`$ finite). In this case all the $`P_{I_j}^2k^2`$, so that $`P_{I_j}^2P_{I_k}^2k`$. Including all the $`T^1k`$ ($`T^2k^0`$ of course), one may imagine from (5.5) that we simply have $`W_{\mu _1\mathrm{}\mu _n}(p_1,\mathrm{},p_n;r,s)W_k`$, where we have used the fact that the leading behaviour apparently arises from having $`a=0`$. Whilst good enough for our purposes the ultraviolet behaviour is clearly better than this, since from (5.5), each term in the Taylor series goes as $`k^{2mn}`$ (independent of $`a`$ and $`b`$) and thus as a whole: $$W_{\mu _1\mathrm{}\mu _n}(p_1,\mathrm{},p_n;r,s)W_k/k^n\underset{\genfrac{}{}{0pt}{}{a,b}{2a+b=n}}{}2^b\underset{\mathrm{perms}}{}\stackrel{~}{T}_1^{i_1}\mathrm{}\stackrel{~}{T}_{a+b}^{i_{a+b}}(srk),$$ where the tensor sum is $`O(k^0)`$, with $`\stackrel{~}{T}^2=T^2`$, $`\stackrel{~}{T}^1=k^{\mu _{I_k}}/k`$. Furthermore we note that this expression agrees with the lower bound from ref. , so we are confident that (5.5) is correct. The problem with the analysis of (5.5) is that when many terms are involved as here, there can be many cancellations in (5.5) of the leading behaviours as determined by power counting. But of course such power counting always yields an upper bound on the large momentum behaviour. From (5.2) we readily find the large momentum behaviour for the $`n+m`$ point vertex in the case the large momentum $`srk`$ flows from end to end, in particular we see that it is $`W_k/k^{m+n}`$ in agreement with the lower bound from . For the more complex cases, the results depend on whether the kernel $`W_k`$ decays or grows for large momentum. In the growing cases, corresponding in this paper to $`W=c^1`$ or $`\stackrel{~}{c}^1`$, we can again obtain a too pessimistic upper bound from power counting (5.5). The exact result, agreeing with the lower bound from , again follows most straightforwardly from (5.5). Thus, from the maximum power of $`k`$ in (5.5), we see in particular that an $`n+m`$ vertex with large momentum $`k`$ entering and leaving from any two points is again $`W_k/k^{m+n}`$. When the kernel $`W_k`$ decays for large $`k`$ (corresponding in this paper to $`c_k^{}`$, $`K_k`$, $`L_k`$ and $`M_k`$) the leading large $`k`$ behaviour depends on the precise configuration. Thus we will show that $$\begin{array}{c}\text{ }W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)\text{ }\hfill \\ \text{ }\frac{k^{\mu _n}}{k^2}W_{\mu _1\mathrm{}\mu _{n1},\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_{n1};q_1,\mathrm{},q_m;r,s+p_n)(sp_nk)(5.5)\text{ }\hfill \\ \text{ }\left(2\frac{k^{\mu _{n1}}k^{\mu _n}}{k^4}\frac{\delta ^{\mu _{n1}\mu _n}}{k^2}\right)W_{\mu _1\mathrm{}\mu _{n2},\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_{n2};q_1\mathrm{}q_m;r,s+p_{n1}+p_n)\text{ }\hfill \\ \text{ }(sp_{n1}k)(5.6)\text{ }\hfill \\ \text{ }\frac{1}{k^{nj+1}}W_{\mu _1\mathrm{}\mu _{j1},\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_{j1};q_1,\mathrm{},q_m;r,s+\underset{i=j}{\overset{n}{}}p_i)(sp_jk)(5.7)\text{ }\hfill \end{array}$$ Closely similar identities then follow easily for $`k`$ leaving via $`q_1`$, $`q_2`$ or any $`q_i`$, from the charge conjugation invariance identity $$\begin{array}{c}\text{ }W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)=\text{ }\hfill \\ \text{ }()^{n+m}W_{\nu _m\mathrm{}\nu _1,\mu _n\mathrm{}\mu _1}(q_m\mathrm{}q_1;p_n\mathrm{}p_1;r,s),(5.8)\text{ }\hfill \end{array}$$ equivalently reversing the direction of the Wilson lines in (4.1), . For the cases where the large momentum both enters and leaves through the side, say $`p_{j+1}p_jk`$, we will show that $$W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}\left[\delta ^{\mu _j\mu _{j+1}}\frac{\mathrm{\Delta }^{\mu _j\mu _{j+1}}(k)}{k^2}\right],$$ where here the notation means to keep only those terms in $`W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}`$ which contain $`\delta ^{\mu _j\mu _{j+1}}`$ and for these make the replacement indicated (no other changes required). In particular this allows us to read off from (5.3), $$W_{\mu \nu }(k,krs;r,s)\frac{\mathrm{\Delta }_{\mu \nu }(k)}{k^2}\frac{W_sW_r}{s^2r^2}$$ and from (5.4), $$W_{\mu \nu }(k,k;p,p)\frac{\mathrm{\Delta }_{\mu \nu }(k)}{k^2\mathrm{\Lambda }^2}W_p^{}.$$ For general separation, on one side: $$W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)1/k^{v1}(p_{j+v}p_jk).$$ Clearly these conclusions hold also for the $`q`$ side (e.g. from (5.8), or exchange or Lorentz symmetry, see sec. 6 or ). If $`k`$ enters and leaves by different sides, say $`p_iq_jk`$, then by (5.2), (5.9) and (5.9), we see that the vertex behaves as $`\mathrm{\Delta }^{\mu _i\nu _j}(k)/k^2`$, irrespective of the values of $`i`$ and $`j`$ (reflecting the fact that in (4.1) these points can always get close to each other). These properties are arrived at, as follows. Setting $`sk`$, $`p_nk`$ in (5.5), all other momenta fixed and finite, the leading term arises when $`i_{a+b}=1`$, through $`T_{a+b}^1k^{\mu _n}`$. Since we can neglect the $`j=a+b`$ term, and factor out $`T_{a+b}^1`$ and the denominator $`P_{I_j}^2P_{a+b}^2k^2`$ in the $`ja+b`$ terms, we are left precisely with the expression for the $`n1`$ point vertex, giving (5.5) for $`m=0`$. Now if $`k`$ leaves through $`p_{n1}k`$, then the leading terms are furnished by $`T_{a+b}^1T_{a+b1}^1`$ and $`T_{a+b}^2`$. Factoring out the now two divergent denominators gives the $`m=0`$ cases of (5.6). Proceeding similarly with the $`k`$ leaving point further down the line, we readily find (5.7) for $`m=0`$. The large momentum behaviour of these three cases is confirmed by the fact that they saturate the lower bounds derived in ref. (and in addition the first two clearly involve no cancellations). Now using (5.2), we readily see that they hold also for $`m>0`$ (the terms on the RHS of (5.2) only contributing at leading order when no $`q`$ momenta get between the divergent pair, $`p_j`$ and $`s`$). The cases where both large momenta are on the side, follow similarly. Thus with $`p_{j+1}p_jk`$ in (5.5), we see that the leading terms come from $`T^2=\delta ^{\mu _j\mu _{j+1}}`$ in which case no $`P_{I_k}`$ diverges, and $`(T^1)^2k^{\mu _j}k^{\mu _{j+1}}`$ with $`1/k^2`$ from factoring out a divergent denominator. Since there is no opportunity for cancellation we can be confident that (5.9) is then correct, despite the fact that the gauge invariance analysis gives a lower bound at $`1/k`$ . Indeed this disagreement between upper and lower bounds is allowed precisely because the leading behaviour (5.9) is transverse in $`k`$ and thus ‘escapes’ the Ward identities. (5.9) follows from a similar analysis and agrees with gauge invariance analysis expectations after taking into account the more divergent initial case (5.9). Finally the $`m>0`$ cases are established as before, by use of (5.2). We note that it is quite straightforward to go further and establish the order of the next-to-leading terms that we have been neglecting (typically down by $`1/k`$), and even their precise form. 5.6. Large momentum behaviour of straight line vertices We finish this section with a brief remark about vertices (5.1). Broadly speaking, the straight line vertices have a comparable behaviour. However, let $`W_k`$ decay for large $`k`$, and consider as an example (5.2), where the large momentum leaves through the side: $$W^\mu (kr;r,k)=\frac{2}{\mathrm{\Lambda }^2}_0^1𝑑tW_{tk+(t1)r}^{}[tk+(t1)r]^\mu .$$ We cannot just take the leading term from the $`t`$ integrand because this leads to an integral that does not converge at $`t=0`$. Substituting $`t=\mathrm{\Lambda }\sqrt{x}/k`$, we obtain the leading behaviour: $$W^\mu (kr;r,k)=\frac{1}{k}_0^{\mathrm{}}\frac{dx}{\sqrt{x}}W^{}(\zeta ^2)\zeta ^\mu ,$$ where $`\zeta =\widehat{k}\sqrt{x}r/\mathrm{\Lambda }`$ (so $`\zeta ^2=x2r.\widehat{k}/\mathrm{\Lambda }+r^2/\mathrm{\Lambda }^2`$), with $`\widehat{k}`$ being the unit vector in direction $`k`$, and corrections to the above being $`O(1/k^2)`$. We see that as in (5.5), $`W_\mu 1/k`$, however unlike (5.5) and the other cases above, it cannot be expressed as an inverse power of $`k^2`$ with a coefficient which is analytic in its momenta (here $`k`$ and $`r`$). Although mostly a matter of taste, it is the more regular large momentum behaviour of the vertices following from (2.3), that led us to use this covariantization for the concrete calculations reported in this paper. 6. Symmetries As well as cyclic and exchange symmetry, and the symmetries of gauge invariance and charge conjugation, that are inherited and preserved from the formulation in ref. (and may be interpreted geometrically in terms of Wilson loops ), some new symmetries appear: fermion number, which is the $`U(1)`$ remainder of the original global $`U(N\|N)`$ symmetry, and an interesting $`Z_2`$ symmetry that exchanges the two $`SU(N)`$ subgroups whilst effectively changing the sign of $`g^2`$. This latter symmetry is thus a “theory space” symmetry, a symmetry of the flow equation (2.3) but not of the action $`S`$. These symmetries provide the key to understanding the formulation at a deeper level. We comment on them below, providing definitions where necessary. We also comment on reality, and include for later some comments on Poincaré invariance and dimensional assignments. 6.1. Cyclicity Some action vertices inherit symmetries from cyclicity of the supertrace (for supermatrices). We have already mentioned this in sec. 4 where these vertices were defined divided by the order of the symmetry group. Thus the vertices in (4.2) are fully cyclically symmetric: $$\widehat{S}_{\mu _1\mathrm{}\mu _n}(p_1,\mathrm{},p_n)=\widehat{S}_{\mu _2\mathrm{}\mu _n\mu _1}(p_2,\mathrm{},p_n,p_1).$$ Similarly, the vertices (4.2) appear in $`\widehat{S}`$ with a factor $`1/2`$, and the $`\mathrm{𝐀𝐀}`$ and $`\mathrm{𝐁𝐁}`$ vertices are consequently symmetric under $`\mu \nu `$, while only the pure $`𝐀`$ vertices in the selections presented in (4.2) and (4.2), have any cyclic symmetry. As mentioned below (4.1), and explained later, odd-loop contributions carry an insertion of $`\sigma _3`$. Since $`𝐀`$ and $`𝐂`$ commute with $`\sigma _3`$ but $`𝐁`$ anticommutes with $`\sigma _3`$, odd-loop vertices with $`𝐁`$ are antisymmetric under cyclic permutations that result in the original order of flavours but cycle an odd number of $`𝐁`$s. As an example, we see that the odd-loop contributions to $`S_{\mu \nu }^{\mathrm{𝐁𝐁}}(p)`$ vanish, because these must be antisymmetric under $`\mu \nu `$, but no such tensor can be constructed. 6.2. Exchange Symmetry From the comment below (2.1), we have $$W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)=W_{\nu _1\mathrm{}\nu _m,\mu _1\mathrm{}\mu _n}(q_1,\mathrm{},q_m;p_1,\mathrm{},p_n;s,r),$$ and similarly for the full vertices, $$\begin{array}{c}\text{ }V_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m,a_1a_2}^{𝐗_1\mathrm{}𝐗_n,𝐘_1\mathrm{}𝐘_m,𝐙_1𝐙_2}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)=\text{ }\hfill \\ \text{ }V_{\nu _1\mathrm{}\nu _m,\mu _1\mathrm{}\mu _n,a_2a_1}^{𝐘_1\mathrm{}𝐘_m,𝐗_1\mathrm{}𝐗_n,𝐙_2𝐙_1}(q_1,\mathrm{},q_m;q_1,\mathrm{},q_n;s,r),\text{ }\hfill \end{array}$$ as is clear from fig. 6, with again similar identities for the $`\alpha \beta `$ vertices of (4.2) – (4.2). 6.3. Poincaré invariance Note that (as usual) all vertices accompany $`\delta `$ functions over the sum of their momentum arguments. The vertices are thus meaningful only when momentum is conserved at the vertex (cf. in particular sec. 5). As in ref. we will use the fact that Lorentz invariance implies that changing the sign of all momentum arguments in any vertex, changes the sign of those with an odd number of Lorentz indices and has no effect on those with an even number. (Of course this applies to any even dimension $`D`$. In odd dimensions we need also parity, which is a symmetry realised straightforwardly here, to rule out the appearance of $`\epsilon _{\mu _1\mathrm{}\mu _D}`$.) 6.4. Dimensions There is of course a scale invariance corresponding to the naïve, or engineering, dimensions. However for a number of reasons (see secs. 2, 3 and 9), especially in general dimension $`D`$, the assignments are a little novel. Of course $`[S]=0`$ but we also have $$[g^2]=4D,[\beta _i]=(D4)i,[_i]=(D4)i+4,[\widehat{}]=4,$$ where the $`_i`$ and $`\widehat{}`$ are the Lagrangians corresponding to $`S_i`$ and $`\widehat{S}`$. And we have $$[𝐀]=[𝐁]=1\text{and}\text{ }[𝐂]=0.$$ 6.5. Gauge invariance Fig.9. Graphical representation of gauge invariance identities. These trivial Ward identities follow from the gauge invariance relations (2.1) or from Wilson line representations as in (4.1), just as they did in ref. . They apply to any pure gauge section (i.e. that came from pure $`𝒜`$ or pure $`A`$) and diagrammatically appear as in fig. 9. Thus, $$q^\nu U_{\mathrm{}a\nu b\mathrm{}}^{\mathrm{}\mathrm{𝐗𝐀𝐘}\mathrm{}}(\mathrm{},p,q,r,\mathrm{})=U_{\mathrm{}ab\mathrm{}}^{\mathrm{}\mathrm{𝐗𝐘}\mathrm{}}(\mathrm{},p,q+r,\mathrm{})U_{\mathrm{}ab\mathrm{}}^{\mathrm{}\mathrm{𝐗𝐘}\mathrm{}}(\mathrm{},p+q,r,\mathrm{}),$$ where $`U`$ is some vertex, and $`a`$ and $`b`$ are Lorentz indices or null as appropriate. Refs. give as an example the relations for pure gauge (seed) action vertices. For wine vertices where the point is at the end of the line : $$\begin{array}{c}\text{ }p_1^{\mu _1}W_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)=(6.1)\text{ }\hfill \\ \text{ }W_{\mu _2\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_1+p_2,p_3,\mathrm{},p_n;q_1,\mathrm{},q_m;r,s)\text{ }\hfill \\ \text{ }W_{\mu _2\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m}(p_2,\mathrm{},p_n;q_1,\mathrm{},q_m;r+p_1,s),\text{ }\hfill \end{array}$$ with similar identities for contraction with $`p_n^{\mu _n}`$, $`q_1^{\nu _1}`$ and $`q_m^{\nu _m}`$, as is clear from fig. 6. In particular, note the signed momentum that appears for $`L_{\mu ,\beta }`$ of (4.2): $$p^\mu L_{\mu ,\beta }(p;q,r)=L_{q,\beta }L_{r,\beta }.$$ 6.6. Charge conjugation invariance Recall that the action of charge conjugation on the gauge fields $`A_\mu ^iA_{\mu }^{i}{}_{}{}^{T}`$, corresponds to reversal of the sign of the underlying Wilson loops . This means that the action is invariant under replacing all supertraces of $`n`$ gauge fields by the signed reversed order: $$\mathrm{str}𝐀_{\mu _1}𝐀_{\mu _2}\mathrm{}𝐀_{\mu _n}()^n\mathrm{str}𝐀_{\mu _n}\mathrm{}𝐀_{\mu _2}𝐀_{\mu _1}.$$ The ‘peppering’ prescription which as we have seen amounts to replacing $`𝐀`$ by $`𝒜=𝐀+𝐁`$, means that symmetry (6.2) must extend to $`𝐁`$ fields also. From the form of $``$ in (2.3), we then see that $`𝐂`$ must also be odd under charge conjugation. However, to implement this symmetry on the fields, we need to include an extra sign for every pair of $`B`$ or $`\overline{B}`$ to compensate for the anticommutation on rearranging the order as in (6.2). Since by fermion number conservation we know there are as many $`\overline{B}`$s as $`B`$s (see below), we can incorporate this by an extra sign in the definition of the supermatrix transpose; thus we define $$𝐗^T:=\left(\begin{array}{cc}X_{}^{11}{}_{}{}^{T}& X_{}^{21}{}_{}{}^{T}\\ X_{}^{12}{}_{}{}^{T}& X_{}^{22}{}_{}{}^{T}\end{array}\right).$$ (We could instead place the sign on $`X^{12}`$, or an $`i`$ \[or $`i`$\] in front of both fermionic parts.) The action of charge conjugation symmetry on the fields is then simply summarised as $$𝒜_\mu 𝒜_\mu ^T,𝐂𝐂^T.$$ Note that the extension to $`SU(N\|N)`$ thus forces charge conjugation to be no longer $`Z_2`$, rather it closes on the $`Z_2`$ in fermion number conservation: $`𝐁𝐁`$, making a $`Z_4`$ in all. At the ‘Wilson loop’ level, the symmetry (6.2) thus holds for all the fields $`𝐀`$, $`𝐁`$ and $`𝐂`$. Of course charge conjugation (6.2) does not act on embedded $`\sigma _3`$s. Instead they maintain their relative position in the Wilson loop, i.e. flip order together with the fields in (6.2) but with $`()^n`$ counting only the fields, as is clear from (6.2). This brings us to an important difference with spontaneously broken $`SU(N\|N)`$: in the present formulation the embedded $`\sigma _3`$s transform differently from $`𝐂`$, and therefore cannot be directly regarded as arising from $`<𝒞>`$. This difference is entirely due to the ‘missing’ $`\sigma _3`$ in $`\widehat{S}_B`$ discussed below (2.3), which would replace the lone $`𝐁`$ in $``$ by $`i\sigma _3𝐁`$, and thus determine $`𝐂`$ to be even under charge conjugation. 6.7. Reality Note that of course, we also have the requirement of reality of the Euclidean action $`\widehat{S}`$ and the flow equation (equivalent, in a time-reversal invariant theory such as this, to unitarity of the Minkowksi theory . Instantons break time-reversal invariance and lead to a complex Euclidean $`S`$, but the underlying $`\widehat{S}`$ and exact RG equation must still be real. Instanton contributions will be considered in more detail elsewhere.) For the gauge fields, the constraints of reality on the vertices follow after Hermitian conjugation and the substitution (or identification) $`A_\mu ^iA_{\mu }^{i}{}_{}{}^{}`$. This extends to a change of variables on the superfields thus, $`𝒜_\mu 𝒜_\mu ^{}`$ and $`𝐂𝐂^{}`$. But note that for the same reasons as before, the transpose part must be defined as in (6.2), and thus the superfields are with this definition only “pseudo-real”: Hermitian conjugation, twice-performed, closes on $`𝐁𝐁`$. Combining reality, charge conjugation and the comments on Poincaré invariance, one readily shows that in momentum space the (wine or action) vertices associated with an odd number of $`𝐂`$ fields (and any number of $`𝐀`$s or $`𝐁`$s) are pure imaginary, whereas those associated with an even number of $`𝐂`$ fields are real. As with the other symmetries outlined in this section, this may be readily verified on the Feynman rules of sec. 4. 6.8. Fermion Number In any vertex there are always as many $`B`$s as $`\overline{B}`$s. This has to be so by the remaining $`SU(N)\times SU(N)`$ invariance (cf. sec. 2 and ref. ), but it also implies the existence of a fermion number $`U(1)`$ symmetry: $`BB\mathrm{e}^{^{i\vartheta }}`$ and $`\overline{B}\overline{B}\mathrm{e}^{^{i\vartheta }}`$ (and similarly with $`B`$ replaced by $`D`$, when not eaten). This $`U(1)`$ is generated by $`\sigma _3`$: $$𝐗\mathrm{e}^{^{i\vartheta \sigma _3/2}}𝐗\mathrm{e}^{^{i\vartheta \sigma _3/2}},$$ where $`𝐗`$ runs over the fields, and thus extends the global $`SU(N\|N)`$ to $`U(N\|N)`$, appearing in this sense in the usual way. Note however there is of course also an essential difference compared to the usual (i.e. bosonic) groups, in that this extra $`U(1)`$ acts non-trivially in the adjoint representation. (Also note that the fermion number $`Z_2`$ invariance mentioned above and whose explicit representation appears below (2.1), is just the subgroup generated by $`\vartheta =\pi `$.) 6.9. Duality While the above symmetries are interesting to explore in their own right and as we will see, are required in practical calculation, by far the most intriguing symmetry we uncovered is a $`Z_2`$ duality symmetry which exchanges the rôle of the two groups in $`SU(N)\times SU(N)`$ and at the same time, in a sense that we make explicit below, changes the sign of the squared coupling constant. At the level of the flow equation (2.3), it is implemented by $$g^2g^2,𝒜𝒜^e\text{and}\text{ }𝐂𝐂^e,$$ where $`𝐗^e:=\sigma _1𝐗\sigma _1`$ and thus for example $$𝒜^e=\left(\begin{array}{cc}A^2& \overline{B}\\ B& A^1\end{array}\right).$$ ($`\sigma _1`$ was introduced in sec. 2.) This transformation is of course not part of $`U(N\|N)`$, not the least because $`\sigma _1`$ has bosonic off-diagonal elements. Note also that $`\left(\sigma _3\right)^e=\sigma _3`$. Thus from the identity $`\mathrm{str}𝐗=\mathrm{tr}\sigma _3𝐗`$ \[or explicitly from (6.2)\] the supertrace of a string of fields is antisymmetric under (6.2). Since from (2.3), $`\widehat{S}`$ is such a single supertrace of superfields, it changes sign. Similarly the RHS of the flow equation (2.3) picks up a sign via the single supertrace in (2.2).<sup>7</sup> Note that the change of variables implied by (6.2), and similarly $`C^1C^2`$, means that the functional derivatives (2.1) are minus their duals. Changing also the sign of $`g^2`$ in $`\mathrm{\Sigma }_g`$ of (2.3), we see that the exact RG is indeed invariant. It follows that if $`S[𝐀,𝐁,𝐂](g^2)`$ is a solution, then so is $`S[𝐀^e,𝐁^e,𝐂^e](g^2)`$. If we imagine $`g`$ in (2.3) to be a fixed (i.e. independent of $`\mathrm{\Lambda }`$) expansion parameter, for example the classical or bare coupling, $`g=g_0`$ at $`\mathrm{\Lambda }=\mathrm{\Lambda }_0`$, then it is easy to see that we can take $`S`$ to be self-dual: $$S[𝐀,𝐁,𝐂](g_0^2)=S[𝐀^e,𝐁^e,𝐂^e](g_0^2).$$ Indeed this follows immediately if as will be the case, the ‘initial’ condition, $`S`$ at $`\mathrm{\Lambda }=\mathrm{\Lambda }_0`$, is taken to be self-dual. As a corollary we find from (3.1), that all even (odd) order loop corrections $`S_n`$ are even (odd) under duality. In the large $`N`$ limit, since $`S`$ is a single supertrace, an even (odd) loop $`S_n`$ must thus contain an even (odd) number of embedded $`\sigma _3`$s, and this fact can readily be confirmed explicitly by considering fig. 2. The self-duality at the very least is obscured when we come to renormalize however. We see immediately that the $`\beta `$ function (3.1) cannot be invariant under (6.2) unless all the odd-loop $`\beta _{2n+1}`$ vanish, which is not the case. From the expansion (3.2), (3.3), etc. , the non-zero $`2\beta _{2n+1}S_0`$ terms mix together terms with even and odd numbers of embedded $`\sigma _3`$s, so that the corollary above no longer holds. The underlying reason for these complications is as follows. By the above analysis, at one loop (using (2.4) and solving $`\beta `$) $$S[𝐀,0,0]=\frac{1}{2}\mathrm{str}d^Dx\left(\frac{1}{g_0^2}+2\beta _1\sigma _3\mathrm{ln}\frac{\mathrm{\Lambda }_0}{\mathrm{\Lambda }}\right)𝐅_{\mu \nu }^2+O(^3).$$ Defining $`1/g^2=1/g_0^2+2\beta _1\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda })`$ only absorbs the divergence for the $`A^1`$ part of the action. We see that in the continuum limit we are forced to introduce two renormalised couplings: $`g=g_1`$ for $`A^1`$, and $`g_2`$ for $`A^2`$. Under duality we map to a solution of (2.3), for which the renormalization condition (2.4) now insists $`g=g_2`$. Let us call the couplings for such a solution $`\stackrel{~}{g}_2=\stackrel{~}{g}`$ and $`\stackrel{~}{g}_1`$. Then by the duality of (2.3) and the initial bare action, $$\begin{array}{cc}\hfill \stackrel{~}{g}_2^2(g_0^2,\mathrm{\Lambda }_0/\mathrm{\Lambda })& =g^2(g_0^2,\mathrm{\Lambda }_0/\mathrm{\Lambda })\hfill \\ \hfill g_2^2(g_0^2,\mathrm{\Lambda }_0/\mathrm{\Lambda })& =\stackrel{~}{g}_1^2(g_0^2,\mathrm{\Lambda }_0/\mathrm{\Lambda }).\hfill \end{array}$$ We comment further in the conclusions. 7. Classical vertices without gauge fixing As we will see, even here there are surprises. Classical solutions will turn out to suffer a form of divergence, arising from integration over $`\mathrm{\Lambda }`$, which is regularised by careful choice of the $`\mathrm{\Lambda }=\mathrm{\Lambda }_0\mathrm{}`$ ‘boundary conditions’, and will require the introduction of some new ‘renormalised’ parameters . This has nothing to do with unbounded momentum integrals, and nothing to do with gauge invariance per se: it arises in the Pauli-Villars sector from the existence of positive powers of the cutoff and some freedom to add extra interactions. However, the implementation of a Pauli-Villars exact RG scheme is in itself one of the novel developments we report in this paper. The dictionary for translating the Feynman diagrams, which themselves follow from expanding (the relevant parts of) fig. 1, has already been given in sec. 4. In fact the expanded Wilson loops look identical to those in refs. , a consequence of the equality of form with that of the pure gauge case of both fig. 1 and the perturbative development in sec. 3. As already noted in sec. 2, the tree-level insertions of $`\sigma _3`$ in (2.1) serve only to ensure that $`𝐘`$ and $`𝐗`$, the remainders of the Feynman diagram on either side, are bosonic or fermionic as appropriate.<sup>8</sup> more strictly block diagonal or block off-diagonal as appropriate But these restrictions are automatically incorporated in the explicit Feynman rules, cf. sec. 4. Constructing the tree-level vertices out of them, expanding using fig. 2 and (anti)commuting the $`\sigma _3`$s together, they all combine to no overall effect. Thus we omit this step in this section. Solving the flow equations for the vertices introduces integration constants, i.e. terms independent of $`\mathrm{\Lambda }`$. Of course these must be chosen to satisfy all the symmetries of the theory, and we will signal which symmetries provide non-trivial constraints. Moreover at the classical level, since $`g`$ does not run, there is no difficulty in preserving the self-duality (6.2), which follows here because the integration constants will be chosen to be single supertraces without embedded $`\sigma _3`$s. Thus in this formulation, there are no classical vertices with an embedded $`\sigma _3`$: $$S_{a_1\mathrm{}a_n}^{0\sigma 𝐗_1\mathrm{}𝐗_n}=0.$$ The reader can find in sec. 6, the list of symmetries together with relevant comments and definitions. It is helpful also, to borrow the conclusions on ‘drifting’ from sec. 9, in particular those summarised in Lemmas 3 and 4 and Corollary 4. Thus we already know that pure-$`𝐀`$ classical vertices are constructed only out of $`\widehat{S}_A`$ vertices, such as in (4.2), and $`c^{}`$ vertices i.e. the non-$`L`$ part of $`c_{\mu _1\mathrm{}\mu _n,\nu _1\mathrm{}\nu _m,\alpha \beta }^{}`$ \[cf. (4.2)\]. In other words, the pure-$`𝐀`$ $`S_0`$ vertices are unchanged from ref. . Similarly to (4.2), let us reserve the unlabelled $`S_{\mu _1\mathrm{}\mu _n}^0`$ vertices for these, which thus have the same explicit expressions as in refs. . We also know there are no $`S_0`$ vertices with just one $`𝐂`$. Contributions that seem at variance with the above conclusions, vanish as a consequence of drifting, which itself is a consequence of the exact preservation of gauge invariance. Although these statements are readily verified, in the interests of compactness we omit the explicit computations. 7.1. Two point vertices From (3.1), we thus have : $$\mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}\text{}_{}^{}{}_{p^\mu }{}^{p^\nu }=2\underset{p^\mu }{\overset{p^\nu }{\text{}}}\underset{p^\mu }{\overset{p^\nu }{\text{}}}+(p^\mu p^\nu )$$ Fig.10. Feynman diagrams for the two-point vertex. As in refs , we now adopt the convention that the empty circle corresponds to $`S_0`$, not $`S`$ as in fig. 1, and we have noted that once again since actions’ one-point vertices vanish (for example by charge conjugation invariance) we must have at least one blob per lobe. Here however, we must also assign flavours $`𝐀`$, $`𝐁`$ or $`𝐂`$ to the two points. From (4.2) and (4.2), at first sight we appear to generate a mixed $`𝐂`$-$`𝐀`$ vertex. Actually, the required $`𝐂`$-$`𝐀`$ zero-point wine vertex anihilates the $`𝐀`$-$`𝐀`$ lobe by gauge invariance: $`p^\mu \widehat{S}_{\mu \nu }(p)=0`$. This is nothing but a perturbative verification of the consequences of ‘drifting’ as discussed above. It is consistent to set to zero all $`S_0`$ vertices for which the RHS of the flow equation vanishes, as here, and we will do so in this paper. (In fact in this case, it may be verified that the requirements of gauge invariance and Lorentz invariance already disallow a $`𝐂`$-$`𝐀`$ vertex.) We are left with: $$\begin{array}{ccc}\hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}S_{\mu \nu }^0(p)& =\frac{1}{2\mathrm{\Lambda }^2}c_p^{}[2\widehat{S}_{\mu \alpha }(p)S_{\mu \alpha }^0(p)]S_{\alpha \nu }^0(p)+(p_\mu p_\nu ),\hfill & (7.1A)\hfill \\ \hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}S_{\mu \nu }^{0\mathrm{𝐁𝐁}}(p)& =\frac{1}{2\mathrm{\Lambda }^2}[2\widehat{S}_{\mu \alpha }^{\mathrm{𝐁𝐁}}(p)S_{\mu \alpha }^{0\mathrm{𝐁𝐁}}(p)]K_{p,\alpha \beta }S_{\beta \nu }^{0\mathrm{𝐁𝐁}}(p)+(p_\mu p_\nu ),\hfill & (7.1B)\hfill \\ \hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}S^{0\mathrm{𝐂𝐂}}(p)& =\frac{1}{\mathrm{\Lambda }^4}M_p\left[2\widehat{S}^{\mathrm{𝐂𝐂}}(p)S^{0\mathrm{𝐂𝐂}}(p)\right]S^{0\mathrm{𝐂𝐂}}(p),\hfill & (7.1C)\hfill \end{array}$$ where we have used (4.2), and on $`(7.1A)`$, the drifting simplifications mentioned above. For completeness we recall how $`(7.1A)`$ is solved . By gauge invariance and dimensions, $$S_{\mu \nu }^0(p)=2\mathrm{\Delta }_{\mu \nu }(p)/f(p^2/\mathrm{\Lambda }^2).$$ From (3.1), we require $`f(0)=1`$ so as to be consistent with (2.4) in the $`g0`$ limit. Substituting (4.2), one readily finds the unique solution to be $`f=c`$. Similarly, substituting (4.2), (4.2), (4.2) and (2.3), one readily verifies that the two-point classical and seed vertices for $`𝐁`$ and $`𝐂`$ may also be taken to be equal. Thus in total: $$S_{\mu \nu }^0(p)=\widehat{S}_{\mu \nu }(p),S_{\mu \nu }^{0\mathrm{𝐁𝐁}}(p)=\widehat{S}_{\mu \nu }^{\mathrm{𝐁𝐁}}(p),S^{0\mathrm{𝐂𝐂}}(p)=\widehat{S}^{\mathrm{𝐂𝐂}}(p).$$ In the case of $`𝐁`$ and $`𝐂`$ however, these are not the most general solutions, presumably reflecting the freedom of reparametrization invariance in the non-gauge sector , but we will specialize to these equalities since they simplify the higher-point vertices as in ref. . In fact the expressions for $`K`$, $`L`$ and $`M`$ in (2.3) were determined to make these equalities possible. (See also sec. 9.) 7.2. Three-point vertices Similarly, from (3.1) and the top two lines of fig. 1, we obtain the following diagrams for the three-point vertex: $$\mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}\text{}=\mathrm{\hspace{0.17em}\hspace{0.17em}2}\text{}+2\text{}$$ Fig.11. Feynman diagrams for the three-point vertex. Here we have already simplified with (7.2), which works precisely in the same way as in ref. . We are thus again left on the RHS with terms that are already determined, allowing the differential equation to be integrated immediately. Again this simplification, provided by the equalities in (7.2), persists to all higher point $`S_0`$ vertices. The proof is identical to that in ref. . Up to cyclicity, the possible three-point flavours are $`\mathrm{𝐀𝐀𝐀}`$, $`\mathrm{𝐁𝐁𝐀}`$, $`\mathrm{𝐁𝐁𝐂}`$, $`\mathrm{𝐂𝐂𝐀}`$, and $`\mathrm{𝐂𝐂𝐂}`$. Note that odd numbers of $`𝐁`$s are ruled out by fermion number (the supertrace of such an odd number vanishes identically) and the classical $`\mathrm{𝐂𝐀𝐀}`$ vertex is ruled out by drifting (cf. discussion above). For the pure $`𝐀`$ vertex, all the $`L`$ terms vanish by drifting (again as outlined above) and we thus obtain the same flow equation as in ref. , with the same result: $$\begin{array}{cc}\hfill S_{\mu \nu \lambda }^0(p,q,r)=& _\mathrm{\Lambda }^{\mathrm{}}\frac{d\mathrm{\Lambda }_1}{\mathrm{\Lambda }_1^3}\left\{c_r^{}\widehat{S}_{\mu \nu \alpha }(p,q,r)\widehat{S}_{\alpha \lambda }(r)+c_\nu ^{}(q;p,r)\widehat{S}_{\mu \alpha }(p)\widehat{S}_{\alpha \lambda }(r)\right\}\hfill \\ & +2(r_\nu \delta _{\mu \lambda }r_\mu \delta _{\nu \lambda })+\mathrm{cycles}.\hfill \end{array}$$ As in refs. , we adopt the convention that the terms in a particular $`\mathrm{\Lambda }`$-integral are understood to have $`\mathrm{\Lambda }`$ replaced by the integration variable (here $`\mathrm{\Lambda }_1`$). The term “cycles” means that we add to the expression the two cyclic permutations of $`(p_\mu ,q_\nu ,r_\lambda )`$. Recall from ref. (see also ref. ), that the continuum limit, corresponding to the upper limit $`\mathrm{\Lambda }_1=\mathrm{}`$, trivially exists, and that the integration constant is fixed by gauge invariance to be the unique covariantization of $`\mathrm{\Delta }_{\mu \nu }`$, i.e. the usual ‘bare’ three-point vertex, as is also clear from the $`\mathrm{\Lambda }\mathrm{}`$ limit of (7.2), since from (2.4), (6.2), and by dimensions $$S_0[𝐀,0,0]=\frac{1}{2}\mathrm{str}d^Dx𝐅_{\mu \nu }^2+O(^3/\mathrm{\Lambda })$$ (or simply by restriction to $`A^1`$). In a similar way we obtain for the $`\mathrm{𝐁𝐁𝐀}`$ vertex: $$\begin{array}{cc}& S_{\mu \nu \lambda }^{0\mathrm{𝐁𝐁𝐀}}(p,q,r)=_\mathrm{\Lambda }^{\mathrm{}}\frac{d\mathrm{\Lambda }_1}{\mathrm{\Lambda }_1^3}\{\widehat{S}_{\mu \nu \alpha }^{\mathrm{𝐁𝐁𝐀}}(p,q,r)c_r^{}\widehat{S}_{\alpha \lambda }(r)+\widehat{S}_{\mu \alpha \lambda }^{\mathrm{𝐁𝐁𝐀}}(p,q,r)K_{q,\alpha \beta }\widehat{S}_{\beta \nu }^{\mathrm{𝐁𝐁}}(q)\hfill \\ & +\widehat{S}_{\alpha \nu \lambda }^{\mathrm{𝐁𝐁𝐀}}(p,q,r)K_{p,\alpha \beta }\widehat{S}_{\beta \mu }^{\mathrm{𝐁𝐁}}(p)+\widehat{S}_{\mu \alpha }^{\mathrm{𝐁𝐁}}(p)c_{\nu ,\alpha \beta }^{}(q;p,r)\widehat{S}_{\beta \lambda }(r)+\widehat{S}_{\lambda \alpha }(r)c_{\mu ,\alpha \beta }^{}(p;r,q)\widehat{S}_{\beta \nu }^{\mathrm{𝐁𝐁}}(q)\hfill \\ & +\widehat{S}_{\nu \alpha }^{\mathrm{𝐁𝐁}}(q)K_{\lambda ,\alpha \beta }(r;q,p)\widehat{S}_{\beta \mu }^{\mathrm{𝐁𝐁}}(p)\}\hfill \\ & +2\left\{(r_\nu \delta _{\mu \lambda }r_\mu \delta _{\nu \lambda })(1+\gamma ^{BBA}/2)+p_\lambda \delta _{\nu \mu }p_\nu \delta _{\lambda \mu }+q_\mu \delta _{\lambda \nu }q_\lambda \delta _{\mu \nu }+\frac{\stackrel{~}{c}_0^{}}{\stackrel{~}{c}_0^2}\delta _{\mu \nu }(qp)_\lambda \right\}\hfill \end{array}$$ where $`\gamma ^{BBA}`$ is a dimensionless real free parameter, and the integral needs interpreting with some care – as we explain below. The first term has been simplified, once again by drifting – i.e. gauge invariance \[as above $`(7.1)`$\]. In fact in this way, expanding the wine vertices by (4.2), all the terms in (7.2) containing $`L`$, either vanish or considerably simplify. We omit the details. To justify (7.2), we replace the top limit in the $`\mathrm{\Lambda }_1`$ integral by $`\mathrm{\Lambda }_0`$, in which case clearly the integration constant is identified with the ‘bare’ value of $`S_{\mu \nu \lambda }^{0\mathrm{𝐁𝐁𝐀}}(p,q,r)`$, i.e. its value at $`\mathrm{\Lambda }=\mathrm{\Lambda }_0`$. Unlike the pure-$`𝐀`$ case viz. (2.4), we do not have a renormalization condition to fix this. (Actually, after identifying all relevant and marginal directions, these could be introduced of course, but as part of the regulating structure they are not needed – nor do they simplify the calculations – at least in this paper.) However by quasilocality and dimensions, since we may discard all terms that vanish in the limit $`\mathrm{\Lambda }_0\mathrm{}`$, finding the most general integration constant reduces to looking for the usual bare-action type terms<sup>9</sup> i.e. polynomial in momenta, balanced as required by non-negative powers of $`\mathrm{\Lambda }_0`$ consistent with the symmetries of the theory. The exact preservation of $`SU(N)\times SU(N)`$ gauge invariance makes this process simple and elegant. Thus the most general integration constant is the $`\mathrm{𝐁𝐁𝐀}`$ part of $$_0\|_{\mathrm{\Lambda }=\mathrm{\Lambda }_0}=\mathrm{str}\left\{\frac{1}{2}_{\mu \nu }^2+\frac{\stackrel{~}{c}_0^{}}{\stackrel{~}{c}_0^2}𝐁_\mu ^2𝐁_\mu +i\gamma ^{BBA}𝐁_\mu 𝐅_{\mu \nu }𝐁_\nu +\mathrm{}\right\},$$ where the first two terms are fixed by gauge invariance from (7.2).<sup>10</sup> viz. gauge covariantizing (4.2), or from (2.3), or directly by expanding (4.2); here and later the solution may also be readily derived directly in momentum space, using (6.1). The ellipses refer to gauge invariant terms that do not contain the $`\mathrm{𝐁𝐁𝐀}`$ vertex (and note by (2.3), include some that diverge in the $`\mathrm{\Lambda }\mathrm{}`$ limit). Note that the last term has the right reality and charge conjugation properties. As we will verify, the universality of the continuum limit means that physical quantities will be independent of $`\gamma ^{BBA}`$. But we cannot quite just set it to zero: by dimensions we see that if the integral in (7.2) yields such a term, then it is logarithmically divergent (viz. $`_\mathrm{\Lambda }^{\mathrm{\Lambda }_0}𝑑\mathrm{\Lambda }_1/\mathrm{\Lambda }_1`$). However, also note that since (7.2) contains the most general possible such $`\mathrm{\Lambda }`$-independent terms, the $`\mathrm{\Lambda }_1`$ integral can only diverge this way. A straightforward calculation confirms that the integral in (7.2) indeed diverges, as $`16\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda })\left(r_\nu \delta _{\mu \lambda }r_\mu \delta _{\nu \lambda }\right)`$, and thus for a finite continuum limit we set the constant to $$\gamma ^{BBA}=16\mathrm{ln}(\mathrm{\Lambda }_0/\mu )+\mathrm{finite},$$ where $`\mu `$ is a finite mass scale, and $`\mathrm{\Lambda }_0\mathrm{}`$. At a more sophisticated level, we may simply impose a definite prescription for discarding the infinities in $`\mathrm{\Lambda }`$ integrals arising in finite continuum solutions such as (7.2), for example minimal subtraction of $`\mathrm{\Lambda }_0`$ divergences, safe in the knowledge that in reality these divergences are actually cancelled by opposite divergences in parameters in the most general integration constant. We could go further with this prescription, and discard these parameters, since this just amounts to choosing them to be precisely the opposing divergences. We will keep them however, as an extra test of universality, but for simplicity report from now on this more sophisticated approach. (Actually, we also checked the calculations the dumb way, as in (7.2) and above. We omit the details.) In this way, the $`\mathrm{𝐁𝐁𝐂}`$ vertex is found to be $$\begin{array}{cc}& S_{\mu \nu }^{0\mathrm{𝐁𝐁𝐂}}(p,q,r)=_\mathrm{\Lambda }^{\mathrm{}}\frac{d\mathrm{\Lambda }_1}{\mathrm{\Lambda }_1^3}\{\frac{1}{\mathrm{\Lambda }_1^2}\widehat{S}_{\mu \nu }^{\mathrm{𝐁𝐁𝐂}}(p,q,r)M_r\widehat{S}^{\mathrm{𝐂𝐂}}(r)+\widehat{S}_{\mu \nu \alpha }^{\mathrm{𝐁𝐁𝐀}}(p,q,r)L_{r,\alpha }\widehat{S}^{\mathrm{𝐂𝐂}}(r)\hfill \\ & +\widehat{S}_{\mu \alpha }^{\mathrm{𝐁𝐁𝐂}}(p,q,r)K_{q,\alpha \beta }\widehat{S}_{\beta \nu }^{\mathrm{𝐁𝐁}}(q)+\widehat{S}_{\alpha \nu }^{\mathrm{𝐁𝐁𝐂}}(p,q,r)K_{p,\alpha \beta }\widehat{S}_{\beta \mu }^{\mathrm{𝐁𝐁}}(p)\widehat{S}_{\mu \alpha }^{\mathrm{𝐁𝐁}}(p)L_{\nu ,\alpha }(q;r,p)\widehat{S}^{\mathrm{𝐂𝐂}}(r)\hfill \\ & +\widehat{S}^{\mathrm{𝐂𝐂}}(r)L_{\mu ,\beta }(p;r,q)\widehat{S}_{\beta \nu }^{\mathrm{𝐁𝐁}}(q)\}+i\gamma ^{BBC}_1\delta _{\mu \nu }(p^2q^2)+i\gamma ^{BBC}_2(p_\mu p_\nu q_\mu q_\nu ),\hfill \end{array}$$ where the two dimensionless real $`\gamma _i^{BBC}`$ parametrise the most general integration constant. Once again, the $`L`$ parts considerably simplify on using the gauge invariance relations (6.1). Note that charge conjugation invariance requires the integration constants to be odd under $`p_\mu q_\nu `$. Dimensions suggest and explicit calculation confirms that the $`\gamma _i^{BBC}`$ mop up logarithmic divergences in the $`\mathrm{\Lambda }_1`$ integral. The $`\mathrm{𝐂𝐂𝐀}`$ vertex is found to be: $$\begin{array}{cc}\hfill S_\lambda ^{0\mathrm{𝐂𝐂𝐀}}(p,q,r)& =_\mathrm{\Lambda }^{\mathrm{}}\frac{d\mathrm{\Lambda }_1}{\mathrm{\Lambda }_1^5}\{\widehat{S}_\lambda ^{\mathrm{𝐂𝐂𝐀}}(p,q,r)[M_pS^{\mathrm{𝐂𝐂}}(p)+M_qS^{\mathrm{𝐂𝐂}}(q)]\hfill \\ & +\mathrm{\Lambda }_1^2\widehat{S}_\alpha ^{\mathrm{𝐂𝐂𝐀}}(p,q,r)c_r^{}\widehat{S}_{\alpha \lambda }(r)+\widehat{S}^{\mathrm{𝐂𝐂}}(q)M_\lambda (r;q,p)\widehat{S}^{\mathrm{𝐂𝐂}}(p)\}\hfill \\ & +2\frac{\stackrel{~}{c}_0^{}}{\stackrel{~}{c}_0^2}(p^2+q^2)(qp)_\lambda +\gamma ^{CCA}(r^2p_\lambda r_\lambda r.p),\hfill \end{array}$$ where $`\gamma ^{CCA}`$ is another dimensionless real free parameter. Here the integration constant is constrained by gauge invariance, (7.2) and (2.3) to be the $`\mathrm{𝐂𝐂𝐀}`$ vertex in $$_0\|_{\mathrm{\Lambda }=\mathrm{\Lambda }_0}=\mathrm{str}_\mu 𝐂\left(\mathrm{\Lambda }_0^2\delta _{\mu \nu }+\frac{\stackrel{~}{c}_0^{}}{\stackrel{~}{c}_0^2}\delta _{\mu \nu }^2+i\gamma ^{CCA}𝐅_{\mu \nu }\right)_\nu 𝐂+\mathrm{},$$ where the ellipses do not contain $`\mathrm{𝐂𝐂𝐀}`$ vertices. The first term precisely cancels an equal and opposite quadratic divergence in the $`\mathrm{\Lambda }_1`$ integral, as it must, and thus by the prescription below (7.2), is not displayed in (7.2). The last term, which verifies the reality and charge conjugation symmetries, mops up a logarithmic divergence in the $`\mathrm{\Lambda }_1`$ integral. Finally, the three-point classical $`𝐂`$ vertex is: $$S^{0\mathrm{𝐂𝐂𝐂}}(p,q,r)=i_\mathrm{\Lambda }^{\mathrm{}}\frac{d\mathrm{\Lambda }_1}{\mathrm{\Lambda }_1^5}\left\{\left[\widehat{S}^{\mathrm{𝐂𝐂}}(p)\widehat{S}^{\mathrm{𝐂𝐂}}(q)\right]L_r\widehat{S}^{\mathrm{𝐂𝐂}}(r)+\mathrm{cycles}\right\},$$ where ‘cycles’ stands for the two cyclic permutations of $`(p,q,r)`$. Here we have ‘drifted’ $`\widehat{S}_\lambda ^{\mathrm{𝐂𝐂𝐀}}L_{r,\lambda }`$ terms by using (4.2) and (6.1). Remarkably, the integration constant must vanish. Charge conjugation and cyclic symmetry require that it be antisymmetric under exchange of any pair of momenta. In this case, any such expression can be generated by writing down Lorentz invariant terms antisymmetric under $`pq`$ and then adding the cyclic iterands. Using momentum conservation, it is then straightforward to show that all such polynomials up to dimension 4 vanish on adding the cycles. Similarly one checks that the superficially quartically divergent integral in (7.2) is actually finite, as consistency requires. 7.3. Four-point vertices As already discussed, the diagrams are the same as in ref. : $$\begin{array}{cc}\hfill \mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}\text{}=\text{}& +2\text{}+2\text{}+2\text{}\hfill \\ & +2\text{}+2\text{}+\text{}\hfill \end{array}$$ Fig.12. Feynman diagrams for the four-point vertex. We will concentrate on the vertices that we will need for the $`\beta _1`$ computation: $`\mathrm{𝐀𝐀𝐀𝐀}`$, $`\mathrm{𝐁𝐁𝐀𝐀}`$ and $`\mathrm{𝐂𝐂𝐀𝐀}`$. (All the three-point $`S_0`$ vertices are however needed, either to derive these four-point vertices or directly.) The $`\mathrm{𝐂𝐀𝐀𝐀}`$ vertex need not be considered since it vanishes by drifting. For the same reasons as given for the three-point vertex (7.2), we obtain the same four-$`𝐀`$ vertex as in ref. : $$\begin{array}{cc}\hfill S_{\mu \nu \lambda \sigma }^0(p,& q,r,s)=_\mathrm{\Lambda }^{\mathrm{}}\frac{d\mathrm{\Lambda }_1}{\mathrm{\Lambda }_1^3}\{c_{p+q}^{}(\widehat{S}_{\mu \nu \alpha }(p,q,r+s)\frac{1}{2}S_{\mu \nu \alpha }^0(p,q,r+s))S_{\alpha \lambda \sigma }^0(p+q,r,s)\hfill \\ & +\widehat{S}_{\sigma \alpha }(s)\widehat{S}_{\mu \nu \alpha }(p,q,r+s)c_\lambda ^{}(r;p+q,s)+\widehat{S}_{\lambda \alpha }(r)\widehat{S}_{\mu \nu \alpha }(p,q,r+s)c_\sigma ^{}(s;r,p+q)\hfill \\ & +\widehat{S}_{\mu \alpha }(p)\widehat{S}_{\alpha \sigma }(s)c_{\nu \lambda }^{}(q,r;p,s)+\frac{1}{2}\widehat{S}_{\mu \alpha }(p)\widehat{S}_{\alpha \lambda }(r)c_{\nu ,\sigma }^{}(q;s;p,r)\hfill \\ & +c_s^{}\widehat{S}_{\sigma \alpha }(s)\widehat{S}_{\mu \nu \lambda \alpha }(p,q,r,s)+\mathrm{cycles}\}+2\delta _{\mu \sigma }\delta _{\nu \lambda }4\delta _{\mu \lambda }\delta _{\nu \sigma }+2\delta _{\mu \nu }\delta _{\lambda \sigma },\hfill \end{array}$$ where ‘cycles’ stands for the three cyclic permutations of $`(p_\mu ,q_\nu ,r_\lambda ,s_\sigma )`$. The $`\mathrm{𝐁𝐁𝐀𝐀}`$ vertex is: $$\begin{array}{cc}& S_{\mu \nu \lambda \sigma }^{0\mathrm{𝐁𝐁𝐀𝐀}}(p,q,r,s)=_\mathrm{\Lambda }^{\mathrm{}}\frac{d\mathrm{\Lambda }_1}{\mathrm{\Lambda }_1^3}\{S_{\mu \nu \alpha }^{0\mathrm{𝐁𝐁𝐀}}(p,q,r+s)c_{p+q}^{}S_{\alpha \lambda \sigma }^0(p+q,r,s)\hfill \\ & S_{\mu \alpha \sigma }^{0\mathrm{𝐁𝐁𝐀}}(p,q+r,s)K_{p+s,\alpha \beta }S_{\beta \nu \lambda }^{0\mathrm{𝐁𝐁𝐀}}(p+s,q,r)+\widehat{S}_{\mu \nu \alpha }^{\mathrm{𝐁𝐁𝐀}}(p,q,r+s)c_{p+q}^{}S_{\alpha \lambda \sigma }^0(p+q,r,s)\hfill \\ & +S_{\mu \nu \alpha }^{0\mathrm{𝐁𝐁𝐀}}(p,q,r+s)c_{p+q}^{}\widehat{S}_{\alpha \lambda \sigma }(p+q,r,s)+\widehat{S}_{\mu \alpha \sigma }^{\mathrm{𝐁𝐁𝐀}}(p,q+r,s)K_{p+s,\alpha \beta }S_{\beta \nu \lambda }^{0\mathrm{𝐁𝐁𝐀}}(p+s,q,r)\hfill \\ & +S_{\mu \alpha \sigma }^{0\mathrm{𝐁𝐁𝐀}}(p,q+r,s)K_{p+s,\alpha \beta }\widehat{S}_{\beta \nu \lambda }^{\mathrm{𝐁𝐁𝐀}}(p+s,q,r)+\widehat{S}_{\mu \nu \lambda \alpha }^{\mathrm{𝐁𝐁𝐀𝐀}}(p,q,r,s)c_s^{}\widehat{S}_{\alpha \sigma }(s)\hfill \\ & +\widehat{S}_{\mu \nu \alpha \sigma }^{\mathrm{𝐁𝐁𝐀𝐀}}(p,q,r,s)c_r^{}\widehat{S}_{\alpha \lambda }(r)+\widehat{S}_{\mu \alpha \lambda \sigma }^{\mathrm{𝐁𝐁𝐀𝐀}}(p,q,r,s)K_{q,\alpha \beta }\widehat{S}_{\beta \nu }(q)\hfill \\ & +\widehat{S}_{\alpha \nu \lambda \sigma }^{\mathrm{𝐁𝐁𝐀𝐀}}(p,q,r,s)K_{p,\alpha \beta }\widehat{S}_{\beta \mu }(p)+\widehat{S}_{\mu \nu \alpha }^{\mathrm{𝐁𝐁𝐀}}(p,q,r+s)c_\lambda ^{}(r;p+q,s)\widehat{S}_{\alpha \sigma }(s)\hfill \\ & +\widehat{S}_{\mu \alpha \sigma }^{\mathrm{𝐁𝐁𝐀}}(p,q+r,s)c_{\nu ,\alpha \beta }^{}(q;p+s,r)\widehat{S}_{\beta \lambda }(r)+\widehat{S}_{\lambda \sigma \alpha }(r,s,p+q)c_{\mu ,\alpha \beta }^{}(p;r+s,q)\widehat{S}_{\beta \nu }^{\mathrm{𝐁𝐁}}(q)\hfill \\ & +\widehat{S}_{\alpha \nu \lambda }^{\mathrm{𝐁𝐁𝐀}}(p+s,q,r)K_{\sigma ,\alpha \beta }(s;q+r,p)\widehat{S}_{\beta \mu }^{\mathrm{𝐁𝐁}}(p)+\widehat{S}_{\mu \nu \alpha }^{\mathrm{𝐁𝐁𝐀}}(p,q,r+s)c_\sigma ^{}(s;r,p+q)\widehat{S}_{\alpha \lambda }(r)\hfill \\ & +\widehat{S}_{\mu \alpha \sigma }^{\mathrm{𝐁𝐁𝐀}}(p,q+r,s)K_{\lambda ,\beta \alpha }(r;q,p+s)\widehat{S}_{\beta \nu }^{\mathrm{𝐁𝐁}}(q)+\widehat{S}_{\lambda \sigma \alpha }(r,s,p+q)c_{\nu ,\beta \alpha }^{}(q;p,r+s)\widehat{S}_{\beta \mu }^{\mathrm{𝐁𝐁}}(p)\hfill \\ & +\widehat{S}_{\alpha \nu \lambda }^{\mathrm{𝐁𝐁𝐀}}(p+s,q,r)c_{\mu ,\beta \alpha }^{}(p;s,q+r)\widehat{S}_{\beta \sigma }(s)+\widehat{S}_{\mu \alpha }^{\mathrm{𝐁𝐁}}(p)c_{\nu \lambda ,\alpha \beta }^{}(q,r;p,s)\widehat{S}_{\beta \sigma }(s)\hfill \\ & +\widehat{S}_{\sigma \alpha }(s)K_{\mu \nu ,\alpha \beta }(p,q;s,r)\widehat{S}_{\beta \lambda }(r)+\widehat{S}_{\lambda \alpha }(r)c_{\sigma \mu ,\alpha \beta }^{}(s,p;r,q)\widehat{S}_{\beta \nu }^{\mathrm{𝐁𝐁}}(q)\hfill \\ & +\widehat{S}_{\nu \alpha }^{\mathrm{𝐁𝐁}}(q)c_{\lambda \sigma ,\alpha \beta }^{}(r,s;q,p)\widehat{S}_{\beta \mu }^{\mathrm{𝐁𝐁}}(p)+\widehat{S}_{\mu \alpha }^{\mathrm{𝐁𝐁}}(p)c_{\nu ,\sigma ,\alpha \beta }^{}(q;s;p,r)\widehat{S}_{\beta \lambda }(r)\hfill \\ & +\widehat{S}_{\nu \alpha }^{\mathrm{𝐁𝐁}}(q)c_{\lambda ,\mu ,\alpha \beta }^{}(r;p;q,s)\widehat{S}_{\beta \sigma }(s)\}\hfill \\ & +2\delta _{\mu \sigma }\delta _{\nu \lambda }4\delta _{\mu \lambda }\delta _{\nu \sigma }+2\delta _{\mu \nu }\delta _{\lambda \sigma }2\frac{\stackrel{~}{c}_0^{}}{\stackrel{~}{c}_0^2}\delta _{\mu \nu }\delta _{\lambda \sigma }+\gamma ^{BBA}(\delta _{\mu \sigma }\delta _{\nu \lambda }\delta _{\mu \lambda }\delta _{\nu \sigma }).\hfill \end{array}$$ The integrand is a straightforward translation from the diagrams, together with some cancellations of $`L`$ wine terms, in particular all $`𝐂`$-$`𝐀`$ wine terms, as a result of drifting. Further such simplifications occur only at the expense of expanding the seed vertices containing $`𝐁`$, as $`\widehat{S}_A+\widehat{S}_B`$, to isolate the pure gauge part $`\widehat{S}_A`$. Note that, as with the four-$`𝐀`$ vertex , gauge invariance (6.1) strongly constrains the overall form of the integrand and acts as a powerful consistency check. The integration constant is dimension zero and therefore cannot generate any new gauge invariant terms. Consequently the integration constant follows from expansion of (7.2), and the $`\mathrm{\Lambda }_1`$ integral diverges in just such a way as to be cancelled by the divergence in (7.2). By drifting, all the $`L`$ wine terms cancel out in the $`\mathrm{𝐂𝐂𝐀𝐀}`$ vertex: $$\begin{array}{cc}& S_{\lambda \sigma }^{0\mathrm{𝐂𝐂𝐀𝐀}}(p,q,r,s)=_\mathrm{\Lambda }^{\mathrm{}}\frac{d\mathrm{\Lambda }_1}{\mathrm{\Lambda }_1^3}\{S_\alpha ^{0\mathrm{𝐂𝐂𝐀}}(p,q,r+s)c_{p+q}^{}S_{\alpha \lambda \sigma }^0(p+q,r,s)\hfill \\ & \frac{1}{\mathrm{\Lambda }_1^2}S_\sigma ^{0\mathrm{𝐂𝐂𝐀}}(p,q+r,s)M_{p+s}S_\lambda ^{0\mathrm{𝐂𝐂𝐀}}(p+s,q,r)+\widehat{S}_\alpha ^{\mathrm{𝐂𝐂𝐀}}(p,q,r+s)c_{p+q}^{}S_{\alpha \lambda \sigma }^0(p+q,r,s)\hfill \\ & +S_\alpha ^{0\mathrm{𝐂𝐂𝐀}}(p,q,r+s)c_{p+q}^{}\widehat{S}_{\alpha \lambda \sigma }(p+q,r,s)+\frac{1}{\mathrm{\Lambda }_1^2}\widehat{S}_\sigma ^{\mathrm{𝐂𝐂𝐀}}(p,q+r,s)M_{p+s}S_\lambda ^{0\mathrm{𝐂𝐂𝐀}}(p+s,q,r)\hfill \\ & +\frac{1}{\mathrm{\Lambda }_1^2}S_\sigma ^{0\mathrm{𝐂𝐂𝐀}}(p,q+r,s)M_{p+s}\widehat{S}_\lambda ^{\mathrm{𝐂𝐂𝐀}}(p+s,q,r)+\widehat{S}_{\lambda \alpha }^{\mathrm{𝐂𝐂𝐀𝐀}}(p,q,r,s)c_s^{}\widehat{S}_{\alpha \sigma }(s)\hfill \\ & +\widehat{S}_{\alpha \sigma }^{\mathrm{𝐂𝐂𝐀𝐀}}(p,q,r,s)c_r^{}\widehat{S}_{\alpha \lambda }(r)+\frac{1}{\mathrm{\Lambda }_1^2}\widehat{S}_{\lambda \sigma }^{\mathrm{𝐂𝐂𝐀𝐀}}(p,q,r,s)\left[M_q\widehat{S}^{\mathrm{𝐂𝐂}}(q)+M_p\widehat{S}^{\mathrm{𝐂𝐂}}(p)\right]\hfill \\ & +\widehat{S}_\alpha ^{\mathrm{𝐂𝐂𝐀}}(p,q,r+s)c_\lambda ^{}(r;p+q,s)\widehat{S}_{\alpha \sigma }(s)+\frac{1}{\mathrm{\Lambda }_1^2}\widehat{S}_\lambda ^{\mathrm{𝐂𝐂𝐀}}(s+p,q,r)M_\sigma (s;q+r,p)\widehat{S}^{\mathrm{𝐂𝐂}}(p)\hfill \\ & +\widehat{S}_\alpha ^{\mathrm{𝐂𝐂𝐀}}(p,q,r+s)c_\sigma ^{}(s;r,p+q)\widehat{S}_{\alpha \lambda }(r)+\frac{1}{\mathrm{\Lambda }_1^2}\widehat{S}_\sigma ^{\mathrm{𝐂𝐂𝐀}}(p,q+r,s)M_\lambda (r;q,p+s)\widehat{S}^{\mathrm{𝐂𝐂}}(q)\hfill \\ & +\frac{1}{\mathrm{\Lambda }_1^2}\widehat{S}^{\mathrm{𝐂𝐂}}(q)M_{\lambda \sigma }(r,s;q,p)\widehat{S}^{\mathrm{𝐂𝐂}}(p)\}+2\frac{\stackrel{~}{c}_0^{}}{\stackrel{~}{c}_0^2}\{(2q+r)_\lambda (2p+s)_\sigma (p^2+q^2)\delta _{\lambda \sigma }\}\hfill \\ & +\gamma ^{CCA}\{p_\lambda q_\sigma p_\sigma q_\lambda \delta _{\lambda \sigma }p.s+s_\lambda p_\sigma \delta _{\lambda \sigma }q.r+r_\sigma q_\lambda \}+\gamma ^{CCAA}(\delta _{\lambda \sigma }r.ss_\lambda r_\sigma ).\hfill \end{array}$$ Gauge invariance again acts as a powerful consistency check on the integrand. The dimension two integration constant allows for one new gauge invariant term $`\gamma ^{CCAA}\mathrm{𝐂𝐅}_{\mu \nu }^2𝐂`$, otherwise gauge invariance requires the rest to be taken from (7.2). The quadratic divergence $`2\mathrm{\Lambda }_0^2\stackrel{~}{c}_0^1\delta _{\lambda \sigma }`$ is not displayed but cancels exactly such a divergence in the integral. Together with $`\gamma ^{CCA}`$, our final free real dimensionless parameter $`\gamma ^{CCAA}`$, mop up all logarithmic divergences in the integral. 8. The $`\beta `$ function without gauge fixing The $`\beta `$ function is determined in essentially the same way as in ref. . As for $`\widehat{S}`$ and $`S_0`$, when the vertices are not labelled by their flavours, we now mean the pure-$`𝐀`$ vertices. Using the renormalization condition (2.4), we have for the $`\mathrm{𝐀𝐀}`$ vertices $$S_{\mu \nu }(p)+S_{\mu \nu }^\sigma (p)=2/g^2\mathrm{\Delta }_{\mu \nu }(p)+O(p^3),$$ and thus by (7.2) and (4.2), $$S_{\mu \nu }(p)+S_{\mu \nu }^\sigma (p)=\frac{1}{g^2}S_{\mu \nu }^0(p)+O(p^3).$$ By (3.1) and (7.1), this implies that the $`O(p^2)`$ component of all the higher loop contributions $`S_{\mu \nu }^n(p)+S_{\mu \nu }^{n\sigma }(p)`$, must vanish. This greatly simplifies the $`O(p^2)`$ part of the $`\mathrm{𝐀𝐀}`$ vertex flow in (3.2) – (3.3), in particular reducing them to algebraic equations. Thus we see that $$a_1[S_02\widehat{S}]_{\mu \nu }^\sigma (p)=4\beta _1\mathrm{\Delta }_{\mu \nu }(p)+O(p^3),$$ where $`a_1[S_02\widehat{S}]_{\mu \nu }^\sigma (p)`$ is the $`\sigma _3\mathrm{𝐀𝐀}`$ vertex in $`a_1[S_02\widehat{S}]`$ and we have used the fact that all one-loop vertices contain a trapped $`\sigma _3`$ \[cf. discussion above (6.2)\]. This fixes $`\beta _1`$. Similarly, at $`n2`$ loops, the $`\beta _n`$ are determined by the requirement that $$a_1[S_{n1}]_{\mu \nu }(p)+a_1[S_{n1}]_{\mu \nu }^\sigma (p)=4\beta _n\mathrm{\Delta }_{\mu \nu }(p)+O(p^3).$$ And non-perturbatively from (3.1) and (3.1), $$a_1[g^2S2\widehat{S}]_{\mu \nu }(p)+a_1[g^2S2\widehat{S}]_{\mu \nu }^\sigma (p)=\frac{4}{g^3}\beta (g)\mathrm{\Delta }_{\mu \nu }(p)+O(p^3).$$ 8.1. One loop $`\beta `$ function without gauge fixing The diagrams giving the LHS of (8.1) take the same form as in ref. (apart from the fact that the large $`N`$ limit has been taken): $$\frac{2}{N}\text{}+\frac{2}{N}\text{}+\frac{2}{N}\text{},$$ Fig.13. The one-loop two-point diagrams. where, as in ref., we let the circle stand for $`\mathrm{\Sigma }_0=S_02\widehat{S}`$. Note that from (7.2) and (3.2), these diagrams also sum to $`\mathrm{\Lambda }S_{\mu \nu }^1(p)/\mathrm{\Lambda }2\beta _1\widehat{S}_{\mu \nu }(p)`$, and thus after $`\mathrm{\Lambda }`$ integration, yield the one-loop $`\mathrm{𝐀𝐀}`$ vertex. Translating as described in sec. 4, since there is no tree-level vertex with a single $`𝐂`$, we see that the wines in fig. 13 attach via two $`𝐀`$s, two $`𝐁`$s or two $`𝐂`$s. Expanding these attachments as in fig. 2, the terms that survive have a $`\sigma _3`$ pair at the top or bottom of the wine (cf. the discussion below (2.1) – the terms with either both $`\sigma _3`$ pairs, or neither $`\sigma _3`$ pair, vanish by the supertrace mechanism). Thus we obtain:<sup>11</sup> Note that the wine vertices appear reflected in fig. 13 compared to sec. 4. $$\begin{array}{cc}& a_1[\mathrm{\Sigma }_0]_{\mu \nu }^\sigma (p)=\frac{2}{\mathrm{\Lambda }^2}\frac{d^Dk}{(2\pi )^D}\{c_{k,\alpha \beta }^{}\mathrm{\Sigma }_{\alpha \beta \mu \nu }^0(k,k,p,p)\hfill \\ & +c_{\mu ,\alpha \beta }^{}(p;kp,k)\mathrm{\Sigma }_{\alpha \beta \nu }^0(pk,k,p)+c_{\mu \nu ,\alpha \beta }^{}(p,p;k,k)\mathrm{\Sigma }_{\alpha \beta }^0(k)\hfill \\ & +\frac{1}{\mathrm{\Lambda }^2}M_k\mathrm{\Sigma }_{\mu \nu }^{0\mathrm{𝐂𝐂𝐀𝐀}}(k,k,p,p)+\frac{1}{\mathrm{\Lambda }^2}M_\mu (p;kp,k)\mathrm{\Sigma }_\nu ^{0\mathrm{𝐂𝐂𝐀}}(pk,k,p)\hfill \\ & +\frac{1}{\mathrm{\Lambda }^2}M_{\mu \nu }(p,p;k,k)\mathrm{\Sigma }^{0\mathrm{𝐂𝐂}}(k)K_{k,\alpha \beta }\mathrm{\Sigma }_{\alpha \beta \mu \nu }^{0\mathrm{𝐁𝐁𝐀𝐀}}(k,k,p,p)\hfill \\ & K_{\mu ,\alpha \beta }(p;kp,k)\mathrm{\Sigma }_{\alpha \beta \nu }^{0\mathrm{𝐁𝐁𝐀}}(pk,k,p)K_{\mu \nu ,\alpha \beta }(p,p;k,k)\mathrm{\Sigma }_{\alpha \beta }^{0\mathrm{𝐁𝐁}}(k)\},\hfill \end{array}$$ where we have noted that, by Lorentz invariance, the $`p_\mu p_\nu `$ transposition just yields a factor 2 . The momentum integral is finite as required, for appropriately chosen covariantization and cutoff functions. We use (2.3) (evaluated in sec. 5) and keep $`c`$ and $`\stackrel{~}{c}`$ general except for the requirements $`c(0)=1`$ and $`\stackrel{~}{c}(0)>0`$ discussed in sec. 2, and $$c(x)x^r\text{and}\text{ }\stackrel{~}{c}(x)x^{\stackrel{~}{r}}$$ for large $`x`$. Physically, we are interested in $`D=4`$ where $`g`$ is marginal, and $`\beta _1`$ is universal. We set $$r>\stackrel{~}{r}>1.$$ This is helpful for calculation, but more stringent than necessary, since we show in the next section that physical one-loop contributions are finite for $`r>\stackrel{~}{r}>0`$. Note that although (8.1) is thus finite, it is the sum of integrals with cancelling divergences and therefore still ambiguous, a generic problem with Pauli-Villars regularisation . We will keep $`D`$ general and take the limit $`D4`$ only at the end of the calculation. This amounts to dimensional preregularisation , and allows us to discard those parts which are surface terms in any dimension $`D`$, but keep those parts which become surface terms only in $`D=4`$ dimensions; the independence of $`\beta _1`$ on the choice of cutoff functions, i.e. its universality, arises from being expressed entirely through such latter terms. For an explicit demonstration of these subtleties, see the example in ref. . Further comments can be found later in this section and in the conclusions. By (7.2), in the two-point vertices of (8.1) we have $`\mathrm{\Sigma }_0\widehat{S}`$. Applying (6.1), the first term of (8.1) becomes $$c_{k,\alpha \beta }^{}\mathrm{\Sigma }_{\alpha \beta \mu \nu }^0(k,k,p,p)c_k^{}\mathrm{\Sigma }_{\alpha \alpha \mu \nu }^0(k,k,p,p)+L_k\left\{\widehat{S}_{\mu \nu }(p)\widehat{S}_{\mu \nu }(k+p)\right\},$$ where we use the above comment, and the fact that $`L_kk_\alpha `$ vanishes by Lorentz invariance of $`k`$ integral. The expression for $`\mathrm{\Sigma }_{\alpha \alpha \mu \nu }^0(k,k,p,p)`$ derived from the last section, may be further simplified by using Lorentz invariance (in fact $`kk`$, $`\mu \nu `$ etc. ) and the coincident line identities (5.2) in a manner already described in ref. . After these manipulations, we simply substituted all the expressions derived earlier, and expanded to $`O(p^2)`$ to extract $`\beta _1`$ via (8.1). The large number of resulting terms were handled with care by algebraic computing. (We used FORM.) This calculation is described below. It seems likely that we could have proceeded more intelligently: after all, $`\beta _1`$ should be universal not only to choices of $`c`$ and $`\stackrel{~}{c}`$ but also to the choices of wine vertex and seed action $`\widehat{S}`$, i.e. covariantization (which, recall, need not even be the same in each wine). Therefore the result should fall out without the need for all these expressions to be explicit. We leave such investigations for the future. A formula for the gauge dependent part falls out simply as follows, and acts as a check on the calculation . Using (6.1) and (6.1), collecting terms using Lorentz invariance under $`kk`$, and again using (7.2), $$\begin{array}{cc}\hfill p^\mu p^\nu a_1[\mathrm{\Sigma }_0]_{\mu \nu }^\sigma (p)& =\frac{2}{\mathrm{\Lambda }^2}\frac{d^Dk}{(2\pi )^D}\{c_{k,\alpha \beta }^{}\widehat{S}_{\alpha \beta }(k)+\frac{1}{\mathrm{\Lambda }^2}M_k\widehat{S}^{\mathrm{𝐂𝐂}}(k)K_{k,\alpha \beta }\widehat{S}_{\alpha \beta }^{\mathrm{𝐁𝐁}}(k)\hfill \\ & c_{k+p,\alpha \beta }^{}\widehat{S}_{\alpha \beta }(k+p)\frac{1}{\mathrm{\Lambda }^2}M_{k+p}\widehat{S}^{\mathrm{𝐂𝐂}}(k+p)+K_{k+p,\alpha \beta }\widehat{S}_{\alpha \beta }^{\mathrm{𝐁𝐁}}(k+p)\}.\hfill \end{array}$$ By a shifting momentum $`kk+p`$ in the first three terms we see that the gauge dependent part is zero as it must be. However, expanding (8.1) to order $`p^4`$ using (4.2), and as in the ensuing calculation, leaving the radial part of the $`k`$ integral to last, we obtain $$\begin{array}{cc}\hfill a_1[\mathrm{\Sigma }_0]_{\mu \nu }^\sigma (p)& =\hfill \\ \hfill 4(D1)& \mathrm{\Omega }_D\mathrm{\Lambda }^{D4}\left\{\mathrm{\Lambda }^2\delta _{\mu \nu }_0^{\mathrm{}}𝑑x\frac{d}{dx}G_0+p_\mu p_\nu _0^{\mathrm{}}𝑑x\frac{d}{dx}G_L\right\}4\beta _1\mathrm{\Delta }_{\mu \nu }(p)+O(p^3)\hfill \\ \hfill \mathrm{where}G_0& =\frac{x^{D/2}F^{}}{D(D1)},G_L=\frac{1}{D(D+2)(D1)}\left(x^{D/2+1}F^{\prime \prime }\right)^{}\hfill \\ \hfill \mathrm{and}F& =(D1)(c^{}K)x/c(DK+xL)/\stackrel{~}{c}+(x/\stackrel{~}{c}+\sigma )M.\hfill \end{array}$$ Here $`\mathrm{\Omega }_D=2/[\mathrm{\Gamma }(D/2)(4\pi )^{D/2}]`$ is the solid angle of a $`(D1)`$-sphere divided by $`(2\pi )^D`$, and as before prime is differentiation with respect to its argument (here $`x`$). The transverse part follows from (8.1). One readily verifies, cf. (9.4), that with $`r>\stackrel{~}{r}>0`$, $`F1/x^f`$ as $`x\mathrm{}`$ with $`f>1`$, and thus the $`\delta _{\mu \nu }`$ and $`p_\mu p_\nu `$ contributions to (8.1) do indeed vanish. We sketch the main steps of the remaining calculation. After simplifying (8.1) as discussed earlier, and substituting for all $`S^0`$ and $`\widehat{S}`$ four-point vertices, we simplify the special-momenta terms by recognizing that both $`S_{a_1a_2\lambda }^{0\mathrm{𝐗𝐗𝐀}}(q,q,0)`$ and $`\widehat{S}_{a_1a_2\lambda }^{\mathrm{𝐗𝐗𝐀}}(q,q,0)`$ collapse to $`\frac{}{q^\lambda }\widehat{S}_{a_1a_2}^{\mathrm{𝐗𝐗}}(q)`$, as follows from gauge invariance.<sup>12</sup> cf. the notation of (4.1); use (6.1) on $`s^\lambda S_{\mu \nu \lambda }(q,r,s)`$ and expand for small $`s`$. After substituting for all remaining vertices, the wine expressions (5.3), (5.3) and (5.4) are substituted, together with the special-momentum case (5.4). No other special-momenta terms requiring careful limits, arise. We collected terms under the substitution $`kk`$, and expanded to overall order $`p^2`$, which is now straightforward. The angular part of the $`k`$ integration is achieved through the following equivalences under the integral $$(k.p)^2k_\mu k_\nu \frac{k^4}{D(D+2)}(p^2\delta _{\mu \nu }+2p_\mu p_\nu )\text{and}\text{ }k_\mu k_\nu \frac{k^2}{D}\delta _{\mu \nu },$$ as follows from considering Lorentz invariance (of in the first case $`k_\alpha k_\beta k_\mu k_\nu `$). The result is written as a linear combination of $`\delta _{\mu \nu }`$, $`p_\mu p_\nu `$ and $`\mathrm{\Delta }_{\mu \nu }(p)`$. As well as the radial $`k`$ integral, there are a number of $`\mathrm{\Lambda }`$ integrals to evaluate. The most involved of these are the integrals for the classical four-point vertices since these contain in turn, some inner $`\mathrm{\Lambda }`$ integrals for classical three-point vertices. Starting with these innermost $`\mathrm{\Lambda }`$ integrals, we substitute for $`K`$, $`L`$ and $`M`$ using (2.3), and expand the derivatives until they are all expressed directly in terms of differentiated cutoff functions, i.e. $`c^{(n)}(x)`$ and $`\stackrel{~}{c}^{(m)}(x)`$; note that these always appear in the numerator of the integrands. The resulting expression is very large, yet it must boil down to the one number $`\beta _1`$. The game then is to manipulate it in such a way as to successively cancel out nearly all the terms. To do this, after any differentiation or integration we express the resulting terms in a unique algebraic form: in this way any algebraic cancellations take place automatically. Denominators involve positive integer powers of $`b(x)=1/(x\stackrel{~}{c}+c)`$ and $`f(x)=1/(x+\sigma c)`$ ($`x=k^2/\mathrm{\Lambda }^2`$ being effectively the integration variable). Iteratively, every appearance of $`xb`$ is replaced by $`(1cb)/\stackrel{~}{c}`$, and every appearance of $`b/x`$ by $`(1/x\stackrel{~}{c}b)/c`$. These relations thus eliminate $`x`$ or/and $`b`$ in $`x^mb^n`$, in favour of $`c`$ and $`\stackrel{~}{c}`$. A number of other ways of bringing such terms to a canonical form also exist, but this one has the advantage that it does not exchange terms which behave well as $`x0`$ or $`\mathrm{}`$, with terms that individually do not. Similarly we tidy $`f`$ terms via iterating $`xf=1\sigma \stackrel{~}{c}f`$ and $`f/x=(1/xf)/(\sigma \stackrel{~}{c})`$. Finally terms with both $`b`$ and $`f`$ in but no explicit power of $`x`$, are brought to canonical form by iterating the substitutions $`\stackrel{~}{c}^2bf=(cbf+\stackrel{~}{c}bf)/\sigma `$ and $`cbf/\stackrel{~}{c}=\sigma \stackrel{~}{c}bf+f/\stackrel{~}{c}b`$. We now integrate by parts the inner integrals as follows. We take first the terms with a $`c^{(n)}(x)`$ or $`\stackrel{~}{c}^{(n)}(x)`$ where the derivative $`n`$ is the largest that appears. (For the inner integrals the maximum derivative is $`n=4`$; overall, it is $`n=5`$.) We then integrate by parts to systematically reduce this maximum $`n`$. If such $`\stackrel{~}{c}^{(n)}`$ terms also contain some positive power $`m`$ of $`\stackrel{~}{c}^{(n1)}`$, we eliminate $`\stackrel{~}{c}^{(n)}`$ by first writing $$\stackrel{~}{c}^{(n)}\left(\stackrel{~}{c}^{(n1)}\right)^m=\frac{1}{m+1}\frac{d}{dx}\left(\stackrel{~}{c}^{(n1)}\right)^{m+1}$$ and similarly for $`c`$. If a $`c^{(n1)}`$ multiplies the terms in (8.1), we still integrate by parts, thereby exchanging a $`\stackrel{~}{c}^{(n)}`$ for a $`c^{(n)}`$ – which we take to be the canonical order. Obstructions to further reduction of the maximum derivative $`n`$ thus occur if there are terms with a product of more than one $`n`$-derivative factor \[i.e. $`c^{(n)}`$ or $`\stackrel{~}{c}^{(n)}`$\] or if the product $`c^{(n)}\stackrel{~}{c}^{(n1)}`$ appears. Although such terms have the potential to exist they all cancel away as we iterate $`n`$ down to $`n=1`$. Incidentally, as mentioned in sec. 7, we also derived $`\beta _1`$ by keeping an explicit initial $`\mathrm{\Lambda }_0`$. Integrating by parts then often yields negative powers of $`\mathrm{\Lambda }_0`$ from the $`\mathrm{\Lambda }=\mathrm{\Lambda }_0`$ boundary (if necessary by Taylor expansion in $`1/\mathrm{\Lambda }_0`$). At this stage it is straightforward to estimate the maximum divergence of factors containing further $`\mathrm{\Lambda }`$ integrals, as a power of $`\mathrm{\Lambda }_0`$ (or as $`\mathrm{ln}\mathrm{\Lambda }_0`$), and thus we drop all terms that go overall as a negative power of $`\mathrm{\Lambda }_0`$; in fact in this way we find that at most a single power of $`1/\mathrm{\Lambda }_0^2`$ need be kept. When there are only first derivatives left the strategy for a unique simplification route is more involved because differentiating $`b`$ ($`f`$) produces $`\stackrel{~}{c}^{}`$ and $`c^{}`$ ($`\stackrel{~}{c}^{}`$). We solved this with the following algorithm. If the integrand contains a single $`c^{}`$, and otherwise only $`c`$ and $`x`$ then the final derivative can be eliminated similarly to (8.1) (introducing $`\mathrm{ln}c`$ if $`m=1`$). A similar strategy applies for $`\stackrel{~}{c}^{}`$, except that we allow a power of $`c`$ – thus exchanging $`\stackrel{~}{c}^{}`$ terms for those with $`c^{}`$, our canonical choice. Next, if there is no $`c^{}`$, one $`\stackrel{~}{c}^{}`$ and a factor of $`b^m`$ where $`m2`$ (there are such terms up to $`m=5`$) then we write: $$b^m\stackrel{~}{c}^{}=\frac{1}{m1}\frac{\stackrel{~}{c}}{c}\frac{d}{dx}b^{m1}+b^m(\stackrel{~}{c}+c^{})/c+b^{m1}\stackrel{~}{c}^{}/c(m2),$$ integrate by parts the explicit $`d/dx`$, and iterate until there are no more such terms. Next, if there is no $`c^{}`$, no $`b`$, one $`\stackrel{~}{c}^{}`$ and a factor of $`f^m`$ where $`m2`$ we write $$f^m\stackrel{~}{c}^{}=\frac{1}{\sigma }\left(f^m+\frac{1}{m1}\frac{d}{dx}f^{m1}\right)(m2),$$ integrate by parts the explicit $`d/dx`$, and iterate until there are no more such terms. After all this is done (mindful always to map the results to a unique algebraic form as above) no further simplification is possible on these integrals. For the $`\delta _{\mu \nu }`$ and $`p_\mu p_\nu `$ terms, these manipulations are sufficient to evaluate all the inner $`\mathrm{\Lambda }`$ integrals and then the outer integrals; we checked that the remaining integral, the radial $`k`$ integral, indeed takes the form (8.1) as expected. (In fact we also checked that this is so, independent of (2.3), again as expected.) For the $`\mathrm{\Delta }_{\mu \nu }(p)`$ terms however, there are quite a few inner $`\mathrm{\Lambda }`$ integrals with at most singly differentiated cutoff functions, that cannot be further simplified. At this stage we integrate by parts the outer $`\mathrm{\Lambda }`$ integral (and single $`\mathrm{\Lambda }`$ integrals), and simplify as much as possible, proceeding in a similar manner to above. (We need also to eliminate $`c^{}`$ in $`c^{}\mathrm{ln}c=(c\mathrm{ln}cc)^{}`$ etc. ) This causes some of the remaining inner $`\mathrm{\Lambda }`$ integrals to disappear by differentiation. Finally, we convert the measure in the $`k`$ integral to $$\mathrm{\Lambda }^D\mathrm{\Omega }_D_0^{\mathrm{}}𝑑xx^{D/21},$$ where $`x=k^2/\mathrm{\Lambda }^2`$, and integrate this by parts, again in the manner sketched above. In the case that we simply set $`D=4`$ from the beginning we find that all the remaining integrals are thus evaluated, converting the expression to a set of boundary terms. In this way, the remaining calculation consists of a set of limits: $`x0`$ and $`x\mathrm{}`$ from the boundaries of the radial $`k`$ integral (and first $`x_0:=k^2/\mathrm{\Lambda }_0^20`$ if the explicit initial $`\mathrm{\Lambda }_0`$ was kept). These limits may themselves involve integrals, however at this stage these integrals are easily done by substituting the appropriate asymptotic behaviour for the integrands as $`x0`$ or $`\mathrm{}`$. Indeed, just by power-counting, the vast majority of the remaining terms simply vanish. Let us give just one example. In order to be finite under (8.1), at one stage we have to collect together some separately divergent terms. Together these yield, using (8.1), a contribution to $`\beta _1`$ of $$\begin{array}{c}\text{ }\frac{1}{(4\pi )^2}\underset{y\mathrm{}}{lim}\{8yc(y)\stackrel{~}{c}(y)\left(^ydx[1/xb(x)\stackrel{~}{c}(x)]/c(x)\right)^2\text{ }\hfill \\ \text{ }[15y\stackrel{~}{c}(y)+3y^2\stackrel{~}{c}^{}(y)]^ydx[1/xb(x)\stackrel{~}{c}(x)]/c(x)\}.\text{ }\hfill \end{array}$$ In here the lower limit of the integrals corresponds to $`\mathrm{\Lambda }=\mathrm{\Lambda }_0`$ and is evaluated first. It is either explicitly finite by cancellation from other parts or the divergences are discarded in the manner already described in sec. 7. Under the ultraviolet limit $`y\mathrm{}`$, the first term vanishes and the second yields $$\frac{1}{(4\pi )^2}\left(\frac{3}{\stackrel{~}{r}1}12\right),$$ as follows from (8.1). After all the limits are evaluated in the $`D=4`$ calculation, we obtain an expression for $`\beta _1`$ which is independent of all the extra parameters $`\gamma `$ introduced in the previous section,<sup>13</sup> Ref. illustrates this by showing how the $`\left(\gamma ^{BBA}\right)^2`$ term disappears. and all the details of $`c`$ and $`\stackrel{~}{c}`$ except that there is still a dependence on the powers $`r`$ and $`\stackrel{~}{r}`$: $$\beta _1=\frac{1}{(4\pi )^2}\left(16\frac{\stackrel{~}{r}}{r}+\frac{8}{r}+\frac{1}{1\stackrel{~}{r}}r^2\stackrel{~}{r}^2+r\stackrel{~}{r}r+\frac{\stackrel{~}{r}}{2}\frac{205}{6}\right).$$ Nevertheless the dependence on $`r`$ and $`\stackrel{~}{r}`$ means that it is therefore not universal – a clear signal that we can attach no physical significance to the result. Indeed, we have already seen that the result will alter if different momentum routings are chosen (e.g. shifting $`k`$ by $`\pm p`$) in different parts of (8.1). We return to step (8.1), and leave $`D`$ general, setting $`D=4ϵ`$. Integrating by parts and expanding in small $`ϵ`$, we now obtain some extra remainder terms which are $`ϵ`$ times integrals that diverge as $`D4`$. (As in dimensional regularisation terms with any higher power of $`ϵ`$ are easily seen to vanish at one loop as $`ϵ0`$.) We do not expect it to make any difference whether $`ϵ`$ is taken positive or negative since the integral is finite for any small $`ϵ`$ providing we are allowed to shift momenta and discard the resulting surface terms. Rigorously however, we choose $`D<4`$ since then all surface terms vanish for $`D`$ small enough, after which the result may defined by analytic continuation to all $`D`$. In the present case any value of $`D<4`$ is sufficient. Throwing away those terms that clearly vanish for $`D<4`$, and restoring the factor $`N`$ absorbed by the large $`N`$ change of variables , we find : $$\begin{array}{cc}\hfill \beta _1& =\frac{N}{(4\pi )^2}\frac{ϵ}{2}[_0^{\mathrm{}}dxx^{ϵ/21}\{2+4c^26c\}+_0^{\mathrm{}}dxx^{ϵ/2}\{\frac{59}{2}\frac{\stackrel{~}{c}}{x\stackrel{~}{c}+c}\frac{59}{6}\frac{1}{x+\sigma \stackrel{~}{c}}\}\hfill \\ & +_0^{\mathrm{}}dxx^{ϵ/2}\stackrel{~}{c}_0^xdx_1\{\frac{20\stackrel{~}{c}}{(x_1\stackrel{~}{c}+c)^2}\frac{2}{\stackrel{~}{c}(x_1+\sigma \stackrel{~}{c})^2}+\frac{24\stackrel{~}{c}^{}}{\stackrel{~}{c}(x_1\stackrel{~}{c}+c)}\frac{6\stackrel{~}{c}^{}}{\stackrel{~}{c}^2(x_1+\sigma \stackrel{~}{c})}\}]\hfill \\ & \frac{11}{3}\frac{N}{(4\pi )^2}\text{as}\text{ }ϵ0,\hfill \end{array}$$ the famous one-loop result for $`SU(N)`$ Yang-Mills, completely independent of the $`\gamma `$s and $`c`$ and $`\stackrel{~}{c}`$, this time including $`r`$ and $`\stackrel{~}{r}`$. (N.B. We display in (8.1), as an example, the penultimate terms in a way closely similar to those that we computed. There is no sense however of a canonical choice here: they are easily altered by adding terms whose differences obviously vanish as $`D4^{}`$. Only the final limit is invariant.) 9. One loop finiteness For the following proof, it is helpful to constrain the form of the covariantization and to specify the general form of the asymptotic behaviour of the cutoff functions. We specialize to covariantizations that are minimal (in the sense of appendix A of ref. ), such as but not exclusively (2.3), and to cutoff functions where the scale is set by the momentum if it is much greater than $`\mathrm{\Lambda }`$ which implies here that they decay as a power for large momenta (cf. again appendix A ), i.e. are of form (8.1) for large $`x`$.<sup>14</sup> i.e. in the sense that $`x^rc(x)`$ and $`x^{\stackrel{~}{r}}\stackrel{~}{c}(x)`$ attain finite non-zero limits as $`x\mathrm{}`$. The conditions that $`c`$ and $`\stackrel{~}{c}`$ decay, and (2.3), imply $$r>0\text{,}\text{ }\stackrel{~}{r}>0\text{and}\text{ }r>\stackrel{~}{r}1.$$ In this section we prove the following theorem, which relies also on the more stringent conditions: $$\begin{array}{ccc}\hfill \mathrm{\Delta }r& >\frac{D2}{2}\hfill & (9.1a)\hfill \\ \multicolumn{3}{c}{\text{where we set }\mathrm{\Delta }r:=r\stackrel{~}{r}\text{, and}}\\ \hfill \stackrel{~}{r}& >\frac{D2}{2}.\hfill & (9.1b)\hfill \end{array}$$ ###### Theorem 1 Conditions (9.1) and $`(9.1)`$ are necessary and sufficient (at $`N=\mathrm{}`$) to ensure that the one-loop corrections in the flow equation (2.3), with any number of external $`𝐀`$ fields and no external $`𝐁`$ or $`𝐂`$ fields, are finite. Note that by replacing $`D`$ in $`(9.1)`$ by $`D+\mathrm{\Delta }`$ we can thus obtain for the one-loop momentum integral over $`k`$, any degree of convergence $`k^\mathrm{\Delta }`$ that we desire. In fact, as we discuss later, for those one-loop corrections other than the vacuum amplitude (i.e. for those corrections containing at least one external $`𝐀`$), we can relax $`(9.1)`$ to $$\mathrm{\Delta }r>\frac{D4}{2}\text{and}\text{ }\stackrel{~}{r}>\frac{D4}{2}.$$ In particular in four dimensions, $`r>\stackrel{~}{r}>0`$ is sufficient for our purposes. Clearly the cases considered are sufficient for the computation of physics to one loop. For compactness we will refer to contributions without external $`𝐁`$s or $`𝐂`$s as BC bereft. Note that the supertraces will generically still contain trapped $`\sigma _3`$s as well as external $`𝐀`$ fields. It may be possible to lift the restriction on no external $`𝐂`$s, which is made here for convenience. On the other hand we show at the end of this section that in the present formulation, there are one-loop corrections with external $`𝐁`$s that diverge, and there are divergent diagrams beyond one loop. Before launching into the proof proper, we sketch the main reasons for finiteness. This in turn helps to explain why the flow equation (2.3) takes the form it does. See also ref. . Again, it is helpful to refer to the three terms in (2.3) as $`\widehat{S}=\widehat{S}_A+\widehat{S}_B+\widehat{S}_C`$, respectively. It is also helpful to refer to the different terms in (2.3) by their covariantized kernels, i.e. as the attachment of the associated wines $`\{c^{}\}`$, $`\{L\}`$, $`\{Kc^{}\}_𝐀`$ and $`\{ML\}`$. (Thus here we will not use the fact that the covariantization is linear e.g. $`\{ML\}=\{M\}\{L\}`$ which implies the $`𝐂`$-$`𝐂`$ attachment is actually via $`\{M\}`$). Apart from ensuring that the high energy behaviour of the Feynman rules is as expected by the $`SU(N\|N)`$ properties (this is discussed in the next three paragraphs below), there are three main mechanisms for finiteness which are essentially very simple: One is the supertrace mechanism that we have already discussed in sec. 2. Another is ‘drifting’ : many potentially distastrously divergent contributions following from attaching $`L`$ in (2.3) via the $`𝒜`$ differential in $`ł`$ on or near $`\widehat{S}_A`$, disappear by gauge invariance. These two mechanisms are sufficient to ameliorate the problems caused by covariantizing the higher derivatives. The final mechanism for finiteness is simply the existence of the higher derivatives themselves. As we will see the proof of finiteness is nevertheless sufficiently involved to make a more mathematical style convenient. We expect that much simpler proofs will be possible in a manifestly local $`SU(N\|N)`$ framework. We choose the kernels (2.3) so that the kinetic terms (bilinears in the fields) in $`\widehat{S}`$ and $`S`$ can coincide as in (7.2), and thus lead to the simplifications we saw in sec. 7, mimicking ref. . This is all the more desirable in the PV sector because we want the Fermi-Bose cancellations that regularise diagrams with $`\widehat{S}_{\mu \nu }`$ in, to work also on those diagrams with $`S_{\mu \nu }`$. In fact the problem would be worse than this because when the two-point functions do not coincide, and thus the cancellations described in sec. 7.2 do not take place, the higher point functions of $`S`$ pick up integration factors that can radically alter their ultraviolet properties. With the cancellations described in sec. 7.2, diagrammatic contributions may be immediately integrated with respect to $`\mathrm{\Lambda }`$, which means in particular that their ultraviolet properties follow more or less straightforwardly from the component vertices, just as they do in the usual application of Feynman rules. (Their leading ultraviolet properties are the same as those of the integrand providing that the coefficient of the leading ultraviolet behaviour of the integrand also converges when integrated over $`\mathrm{\Lambda }`$.) At the same time however, we need the kernels themselves to have the right normalisations to lead to cancellations at high momentum just as the propagators do (as in ref. , or via the spontaneously broken $`SU(N\|N)`$ interpretation). Since the kernels are essentially the $`\mathrm{\Lambda }`$ differentials of these propagators, we need care in choosing the powers of $`\mathrm{\Lambda }`$ multiplying these propagators, equivalently the powers of $`\mathrm{\Lambda }`$ in front of the $`\widehat{S}`$ two-point functions in (2.3), or again equivalently, the natural dimensions of the fields. This requires the assignment of mass dimension one for $`𝐁`$ and the less conventional choice of dimension zero for $`𝐂`$. With these choices we have as required that the transverse $`𝐁`$’s kernel $`K`$ and $`𝐀`$’s kernel $`c^{}`$ coincide at high energy, and the longitudinal $`𝐁`$’s kernel $`L`$ and $`𝐂`$’s kernel $`M`$ coincide at high energy. Indeed from (9.1) and (2.3), $$\begin{array}{cc}\hfill K(x)& =c^{}(x)\left\{1+O\left(1/x^{1+\mathrm{\Delta }r}\right)\right\}\hfill \\ \hfill L(x)& =(x\stackrel{~}{c})^{}/x\left\{1+O\left(1/x^{1+\mathrm{\Delta }r}\right)\right\}\hfill \\ \hfill M(x)& =(x\stackrel{~}{c})^{}/x\left\{1+O\left(1/x^{1+\stackrel{~}{r}}\right)\right\}\hfill \end{array}$$ This high momentum restoration of the unbroken $`SU(N\|N)`$ clearly improves the ultraviolet behaviour of the last two terms in (2.3). Meanwhile, quantum corrections following from differentiating the full $`𝒜`$, will cancel by the supertrace mechanism (2.1), unless a $`\sigma _3`$ is trapped in between. A slightly more subtle supertrace mechanism works between $`𝐁`$ and $`𝐂`$ differentials of $``$ (reflecting the fact that in essence, $`𝐁`$ ate $`𝐃`$). As we have already mentioned, this supertrace mechanism together with the drifting property cure the problems associated with covariantizing the higher derivatives. Indeed, as we now show, the conditions $`(9.1)`$ are anyway required to ensure that there are sufficient higher derivatives to regularise, independently of the problems caused by covariantization. (Note (9.1), which is also required, has already been discussed in sec. 2.) ###### Lemma 2 $`(9.1a)`$ and $`(9.1b)`$ are necessary for ultraviolet finiteness. Fig.14. Simple one loop diagram. Proof. Consider the vacuum contribution to the simple one-loop vacuum diagram $`a_1[\widehat{S}_C]`$, as illustrated in fig. 14, where the wine is $`L`$. Clearly there is no supertrace cancellation since there is no analogous $`𝐁`$ term in (2.3). By (8.1) and (9.2), this contribution is UV finite if and only if $`(9.1b)`$ holds. Similarly, consider $`a_1[\widehat{S}_B]`$ where the kernel is $`c^{}`$ which we attach to the lone $`𝐁`$ in the $``$s. Again there is no supertrace cancellation, because there is no corresponding longitudinal $`𝐀`$ term in (2.3). By (8.1), this contribution is UV finite if and only if $`(9.1a)`$ holds. $``$$``$ Note that $`(9.1a)`$ and $`(9.1b)`$ are strictly necessary in the above examples only for the vacuum contributions. Discarding vacuum energy contributions, the first non-trivial contribution in fig. 14 must be the two-point contribution since the one-point contribution vanishes by $`\mathrm{tr}A^i=0`$ (as well as many other reasons ). By secs. 4 and 5.5 (or more generally appendix A of ref. ), wherever the two points are placed, the large momentum behaviour is improved by a power of two. Therefore in $`(9.1)`$ we can replace $`D`$ by $`D2`$, obtaining (9.2). By keeping careful track of where the external points must be, we can prove these more relaxed conditions for the general contributions, but we will not pursue it further here: the important point we wish to establish is that conditions do exist that guarantee finiteness for all physically relevant one-loop graphs. We have already shown that there are no terms of the form of fig. 3. Therefore a one loop graph is formed only by attaching wines to other wines, or lobes $`\widehat{S}`$, as shown in fig. 15. Fig.15. General form of the one loop diagram. Here we have used the fact that one of the two superloops must be field free to survive the large $`N`$ limit. Without loss of generality, we take it to be the inner Wilson loop. It follows that the inner loop also grows no tree corrections, because tree corrections have at least one lobe with only one wine attached, and these lobes must have at least one external field because there are no $`\widehat{S}`$ one-point vertices. As illustrated in fig. 15, the one loop contribution will have one or more of the following features: (a) The loop momentum $`k`$ is carried directly into and out of a lobe by wines. The lobe $`\widehat{S}`$ can have any number of external fields on the outer Wilson loop and any number of tree corrections branching off it. (b) As in (a) except that one or both wines carrying $`k`$, first attach to another wine, the rest of which is attached to some tree correction. (c) The loop momentum passes through three wines (in one of several ways) with no intermediate lobe. At first sight it is easy to construct a divergent one-loop diagram: we use the $`\widehat{S}_A`$ terms for the lobes and $`\{L\}`$ for the wines. By (9.1) or Lemma 2, the $`k^{2\stackrel{~}{r}}`$ is in general no way sufficient to soften the divergent $`k^{2r}`$ behaviour in $`\widehat{S}_A`$. However, the $`L`$ attachment must be made via the $`_\mu \delta /\delta 𝒜_\mu `$ part of $`l`$, there being no $`𝐂`$ terms in $`\widehat{S}_A`$ or the wines (where attachment may be made as in fig. 15b). Gauge invariance then implies drifting which as we will see, entirely cancels all such dangerous contributions (in BC bereft one-loop diagrams). Thus, consider an $`\{L\}`$ wine attaching to a pure $`𝒜`$ section of a superloop. The attachment takes the form $$\begin{array}{ccc}\hfill d^Dw\mathrm{str}\left[𝐗(w)_\mu \frac{\delta }{\delta 𝒜_\mu (w)}\right]\mathrm{\Phi }_{xy}& =d^Dw\mathrm{str}\left[_\mu 𝐗(w)\frac{\delta }{\delta 𝒜_\mu (w)}\right]\mathrm{\Phi }_{xy}\hfill & \\ & =i\left\{\mathrm{\Phi }_{xy}𝐗(y)𝐗(x)\mathrm{\Phi }_{xy}\right\},\hfill & (9.2)\hfill \end{array}$$ where $`𝐗(w)`$ contains $`\{L\}`$ and all that it attaches to at its other end, $`\mathrm{\Phi }_{xy}[𝒜]`$ is the pure gauge section, with ends at $`x`$ and $`y`$, and the last line follows since $`\mathrm{\Phi }_{xy}`$ transforms homogeneously under gauge transformations (2.1). The interpretation in terms of drifting, or sliding, is now clear from fig. 16. Fig.16. Drifting or sliding $`\{L\}`$. The obstructions at the ends of $`\mathrm{\Phi }_{xy}`$ are $`\sigma _3`$s, $`𝐂`$s, explicit $`𝐁`$s,<sup>15</sup> which may be regarded as arising from $`\sigma _3`$s by writing $`𝐁=\mathrm{d}_{}𝒜`$, cf. (2.1) or ‘remaindered’ $`𝐁`$s which arise in the following sense. The pure gauge section may depend only on $`𝐀`$ rather than the full supermatrix $`𝒜`$, as in $`\{Kc^{}\}_𝐀`$ of (2.3). Drifting still occurs by the first gauge transformation relation (2.1), but because $`\{L\}`$ attaches with the covariant derivative $`=𝐃i𝐁`$, we obtain remainder terms with a $`𝐁`$ at the join. Equivalently we can write the pure $`𝐀`$ section as $`\mathrm{\Phi }_{xy}[𝒜𝐁]`$, so that $`\mathrm{\Phi }_{xy}`$ expands into a $`\mathrm{\Phi }_{xy}[𝒜]`$ plus a series of smaller pure $`𝒜`$ sections with remainder $`𝐁`$s at one or both ends. If we replace the section $`\mathrm{\Phi }_{xy}[𝒜]`$ in (9.2) by a full pure gauge Wilson loop $`\phi [𝒜]`$,<sup>16</sup> Recall that we can think of these pure gauge contributions as fluctuating Wilson loops, integrated over with some suitable measure . then since this is gauge invariant, we find that attaching $`\{L\}`$ just results in zero. This is equally clear from fig. 16 by sowing the two ends of $`\mathrm{\Phi }_{xy}`$ together, i.e. by setting $`x=y`$ and taking a supertrace. $`\{L\}`$ can drift of course across a join between two sections, say $`\alpha `$ and $`\beta `$, in a composite Wilson loop (e.g. drifting off a wine and onto a lobe). We can see this either by treating the two sections together as one pure gauge section $`\mathrm{\Phi }`$ in (9.2), or more carefully as follows. We can add together the diagrams with $`\{L\}`$ attached to $`\alpha `$, and with $`\{L\}`$ attached to $`\beta `$. By (9.2) and fig. 16, the drift to the join from $`\alpha `$ cancels the drift to the join from $`\beta `$. Note however: we are making an assumption when sliding $`\{L\}`$ across a join, that the new diagram actually exists, i.e. can be constructed from (2.3). In fig. 17 we show an exception. Fig.17. A diagram with no external $`𝐁`$s where $`\{L\}`$ cannot slide across the join. We take the diagram to have no external $`𝐁`$s. The lobe is $`\widehat{S}_B`$ and the $`\{c^{}\}`$ wine attaches by differentiating the lone $`𝐁_\mu `$ in $`_\mu `$. The $`\{L\}`$ attaches to $`\{c^{}\}`$, absorbing the $`𝐁`$ that has to be there by fermion number (6.2), and then slides back to the join. The corresponding diagram with $`\{L\}`$ attached to $`\widehat{S}_B`$, but all other factors the same, does not exist however: since $`\{c^{}\}`$ and above is $`𝐁`$-free, in the corresponding diagram $`\{c^{}\}`$ must attach by differentiating an $`𝐀_\mu `$ in the same place as the lone $`𝐁_\mu `$ as illustrated in fig. 18, but there is no such $`𝐀_\mu `$ term. Fig.18. A diagram with no external $`𝐁`$s and drifted $`\{L\}`$, that does not exist. (Such a longitudinal $`𝐀_\mu `$ term is not allowed by gauge invariance.) Fortunately, as we see in this example, firstly the problem only arises where a bosonic attachment must change into a fermionic one (or vice versa) as a consequence of sliding a $`\delta /\delta 𝐁`$ across the join, and secondly the problem is always flagged by the appearance of $`\sigma _3`$s at the join expressing the fact that the attachment is only via a partial supermatrix as in fig. 2. Collecting the observations above, we thus obtain the following lemma. ###### Lemma 3 (a) $`\{L\}`$ wines drift to the ends of pure $`𝒜`$ sections of the superloop. (b) $`\{L\}`$ wines cannot attach to pure $`𝒜`$ Wilson loops. (c) $`\{L\}`$ wines drift to the ends of pure $`𝐀`$ sections up to terms involving remaindered $`𝐁`$s. (d) $`\{L\}`$ wines attach to pure $`𝐀`$ Wilson loops, in particular BC bereft tree-level contributions, only if accompanied by remaindered $`𝐁`$s. Proof. We only need note that in (d) a (BC bereft) tree-level contribution contains no $`\sigma _3`$s since they occur in pairs and thus may be commuted past $`𝐀`$s and combined together until they all disappear (cf. fig. 2and sec. 7), and $`\{L\}`$ has to attach by a bosonic differential since there are only $`𝐀`$s to differentiate, so diagrams exist for $`\{L\}`$ attaching at all points on the tree level contribution. $``$$``$ ###### Lemma 4 There are no tree-level contributions with just one $`𝐁`$ or one $`𝐂`$. Proof. We stress that this holds, whatever the number of $`𝐀`$s. That there are no one-point $`𝐁`$ diagrams is obvious, from the fermion number conservation (6.2) already mentioned. If an external $`𝐂`$ appears it must do so in $`\widehat{S}_B+\widehat{S}_C`$. But then together with any number of $`𝐀`$s, it accompanies in the $`\widehat{S}`$ lobe, either two $`𝐁`$s which at tree level cannot be removed by further attachments (again by fermion number), or another $`𝐂`$. This other $`𝐂`$ can be propagated to other similar lobes by $`\{L\}`$ or $`\{LM\}`$ in (2.3) where it may be turned into pairs of $`𝐁`$s (which again translate to external $`𝐁`$s), but it cannot disappear in any other way except apparently by attaching $`\{L\}`$ where we use the $`𝐂`$ differential in $`ł`$, and at the other end we use the $`_\mu \delta /\delta 𝒜_\mu `$ differential. However, this latter must attach to a BC bereft tree-level contribution and produce no remainder $`𝐁`$s, which is excluded by Lemma 3(d). $``$$``$ ###### Corollary 4 (a) BC bereft trees and BC bereft tree corrections are made only of lobes $`\widehat{S}_A`$ and wines $`\{c^{}\}`$. (b) In a BC bereft one-loop diagram $`\{Kc^{}\}_𝐀`$, $`\{ML\}`$ and $`\{L\}`$ must appear entirely within the loop i.e. such that the loop momentum $`k`$ goes in one end of the wine and out the other, while $`\widehat{S}_B`$ ($`\widehat{S}_C`$) must have the loop momentum enter and leave via the lone $`𝐁`$ or the $`𝐂`$ in $``$ ($`𝐂`$). Proof. (a) For a BC bereft tree or BC bereft tree correction, if there existed $`\widehat{S}_B`$ or $`\widehat{S}_C`$, then breaking open the diagram by removing the wine connection to the $``$ or $`𝐂`$ respectively, results in a tree with just one $`𝐁`$ or $`𝐂`$ in violation of Lemma 4. Similarly if the BC bereft tree (correction) contained $`\{Kc^{}\}_𝐀`$, $`\{ML\}`$ or $`\{L\}`$, then breaking open the diagram at the (tree) end of one of these wines we obtain a tree with one $`𝐁`$ and no $`𝐂`$s, one $`𝐂`$ and no $`𝐁`$s, and either one $`𝐂`$ and no $`𝐁`$s or neither $`𝐁`$s nor $`𝐂`$s, respectively. All these possibilities are excluded by Lemmas 4 and 3(d). (b) For the wines this is just a restatement of (a). For $`\widehat{S}_B`$ and $`\widehat{S}_C`$, if the loop momentum did not enter and leave in this way we would have to have an external $`𝐁`$ or $`𝐂`$ since by Lemma 4, they cannot be absorbed by attaching a tree correction. $``$$``$ Fig.19. An $`\{L\}`$ attached to $`\{c^{}\}`$ or $`\widehat{S}_A`$ which has drifted to touch a $`\{Kc^{}\}_A`$ wine. ###### Proposition 5 In a BC bereft one-loop diagram, $`\{L\}`$ can only attach either: (a) directly to the lone $`𝐁`$ or $`𝐂`$ in $``$ of $`\widehat{S}_B`$, or directly to the $`𝐂`$ in $`\widehat{S}_C`$, (b) to $`\{c^{}\}`$ as in fig. 17, or (c) to $`\{c^{}\}`$ or $`\widehat{S}_A`$ as in fig. 19. Proof. By Corollary 4(b), if $`\{L\}`$ attaches to a lobe it either does so as in (a), or attaches to $`\widehat{S}_A`$. However, it will drift off $`\widehat{S}_A`$ unless it meets an obstruction which may be regarded as being caused $`\sigma _3`$ (cf. discussion above Lemma 3). This in turn would arise only from a wine attached to $`\widehat{S}_A`$ by a partial supermatrix derivative, uniquely specifying the attachment of $`\{Kc^{}\}_𝐀`$ and hence fig. 19. Otherwise $`\{L\}`$ must attach to a wine, which by Corollary 4(b) and fig. 15, must be $`\{c^{}\}`$. If $`\{L\}`$ is not to drift off $`\{c^{}\}`$, it must meet an obstruction which can only be due to either a partial supermatrix derivative attachment, again giving uniquely $`\{Kc^{}\}_𝐀`$ as in fig. 19, or $`\{c^{}\}`$ attaching to a partial supermatrix, which uniquely specifies fig. 17. (By Corollary 4, $`\{c^{}\}`$ cannot attach to $`\{Kc^{}\}_𝐀`$ since one end of the latter would be outside the momentum loop.) $``$$``$ Note that (since $`\{L\}`$ is entirely inside the loop) if $`\{L\}`$ fails to find any obstruction then it will slide up its own tail as in fig. 3, which we saw in sec. 2 gives zero. Equivalently, if it fails to find an obstruction then there are no trapped $`\sigma _3`$s, but since $`\{L\}`$ attaches by full supermatrix differentials, the result vanishes by the supertrace mechanism. In fig. 19, by Corollary 4(b), the loop momentum enters and leaves via $`\{L\}`$ and $`\{Kc^{}\}_𝐀`$. Since these two wines meet at the same point, the pure gauge vertex above (viz. $`\widehat{S}_A`$ or $`\{c^{}\}`$) carries no loop momentum. Similarly in fig. 17, by Corollary 4(b), the loop momentum enters and leaves through $`\{L\}`$ and $`\widehat{S}_B`$, and $`\{c^{}\}`$ carries no loop momentum. We see that drifting has entirely removed the dangerous combinations of $`\{L\}`$ and $`\widehat{S}_A`$ discussed above (9.2), as promised. By fig. 15, there are at least as many wines carrying loop momentum as there are lobes carrying loop momentum. The insertion of extra wines as in fig. 15b and fig. 15c, cannot lessen the superficial<sup>17</sup> i.e. as established by power counting and not taking into account possible cancellations convergence. To be more precise, note that in the case where there is only one lobe carrying loop momentum (for example as in fig. 14), there is always at least one ‘internal’ wine, in the sense of an internal propagator, i.e. a wine carrying loop momentum $`k`$ from end to end. In a one-loop diagram with more than one $`k`$ carrying lobe, there is always at least one such internal wine between any adjacent pair of such lobes. These internal wines are therefore at least equal in number to the lobes carrying $`k`$. The extra wines, over and above these internal wines, would be only partially within the loop, but by (5.7) and (5.9), for large momentum $`k`$ they contribute at worst $`k^0`$. Consider a $`k`$-carrying lobe and all the wines on one side of the lobe that $`k`$ passes through until the next $`k`$-carrying lobe is reached (or the same lobe is reached, if there is only one $`k`$-carrying lobe in the loop). Such a combination forms a dressed internal tadpole, which we will refer to simply as a tadpole. Note that such a tadpole thus always includes at least one internal wine. Every one-loop diagram can be split up (in one of two ways) into a set of disjoint tadpoles as illustrated in fig. 20. Fig.20. An example one-loop diagram split into two tadpoles by the dashed lines. For a given tadpole, and large $`k`$, let $`k^{2p}`$ be the power contributed by the $`k`$-carrying lobe, and let $`k^{2w}`$ be the total power of $`k`$ contributed by the wines (again as established by power counting). We now use (8.1) and the results of sec. 5.5 (more generally appendix A of ), to compute upper bounds for the tadpole’s superficial degree of divergence $`2(pw)`$, for each choice of lobe and choices of internal wine. (i) Suppose that the $`k`$-carrying lobe is $`\widehat{S}_A`$, then the maximum $`p`$ that this can contribute is $`p=r+1`$. (ia) $`\{ML\}`$ cannot be an internal wine in this tadpole since it joins directly to $`\widehat{S}_B+\widehat{S}_C`$ by (2.3) and Corollary 4(b). (ib) If $`\{Kc^{}\}_𝐀`$ appears as an internal wine then by (9.2), $$pw1\mathrm{\Delta }r.$$ (ic) If $`\{L\}`$ is an internal wine in this tadpole then so is $`\{Kc^{}\}_𝐀`$ by Proposition 5, fig. 19 and Corollary 4(b). (N.B. As noted above, the vertex shown above $`\{Kc^{}\}_𝐀`$ and $`\{L\}`$ in fig. 19 is not carrying $`k`$.) Therefore (9.3) at least applies. (id) Finally $`\{c^{}\}`$ could be an internal wine, thus contributing a minimum of $`w=r+1`$, and hence $`pw0`$. (ii) Now suppose that the $`k`$-carrying lobe is $`\widehat{S}_B`$. This will contribute a maximum $`p=\stackrel{~}{r}`$ if attached to the loop by the $`𝐁`$s, or $`p=\stackrel{~}{r}+1`$ if attached to the loop via the $`𝐂`$s. (iia) If $`\{c^{}\}`$ is an internal wine in this tadpole then by Corollary 4(b) (and Lemma 4), the $`\widehat{S}_B`$ must attach via the $`𝐁`$s and $`pw1\mathrm{\Delta }r`$. (iib) Similarly if $`\{Kc^{}\}_𝐀`$ is an internal wine, then by (9.2), $`pw22\mathrm{\Delta }r`$. (iic) If $`\{ML\}`$ appears as an internal wine then this attaches to a $`𝐂`$ in $`\widehat{S}_B`$, and by (9.2), $`pw1\mathrm{min}(\stackrel{~}{r},\mathrm{\Delta }r)`$. (iid) Finally if $`\{L\}`$ is an internal wine in this tadpole then either it can attach directly to the $`𝐂`$s in which case $`pw0`$, or attach via the $`𝐁`$s in which case again $`pw0`$. (iii) Finally, suppose that the $`k`$-carrying lobe is $`\widehat{S}_𝐂`$. This gives $`p=0`$. The least convergent internal wine is $`\{L\}`$, thus we have that whatever internal wines there are, $`pw1\stackrel{~}{r}`$. Collecting these observations together we have ###### Proposition 6 In a BC bereft one-loop diagram, all tadpoles have $`pw0`$. Proof. The result follows trivially from the above observations on using the second two conditions of (9.1). $``$$``$ ###### Proposition 7 Any BC bereft one-loop diagram containing a $`\{Kc^{}\}_𝐀`$ wine, or $`\{ML\}`$ wine, or $`\widehat{S}_B`$ and an internal $`\{c^{}\}`$, or $`\widehat{S}_C`$, is finite. Proof. If $`\{Kc^{}\}_𝐀`$ appears in the diagram, it does so as an internal wine by Corollary 4(b). Splitting the diagram into tadpoles, we have by (ib), (iib) and (iii) above, and Proposition 6, that the degree of divergence of the diagram is no worse than $`D22\mathrm{\Delta }r`$, and thus the diagram is finite by $`(9.1a)`$. Similarly if $`\{ML\}`$ appears, it also must appear as an internal wine. Then by (ia), (iic), (iii) and Proposition 6, the degree of divergence of the diagram is no worse than $`D22\mathrm{min}(\stackrel{~}{r},\mathrm{\Delta }r)`$, and thus by $`(9.1)`$ the diagram is finite. If $`\widehat{S}_C`$ appears, it must do so as a $`k`$-carrying lobe by Corollary 4(b). Splitting the diagram into tadpoles, we see from Proposition 6 and (iii) above, that the diagram is finite by $`(9.1b)`$. Finally, if $`\widehat{S}_B`$ appears, it must do so as a $`k`$-carrying lobe by Corollary 4(b). If an internal $`\{c^{}\}`$ also appears in the loop then either it or another $`\{c^{}\}`$ may incorporated in a tadpole with the $`\widehat{S}_B`$, in which case the diagram is finite by Proposition 6, (iia) and $`(9.1a)`$, or else the internal $`\{c^{}\}`$s must only lie between adjacent $`k`$-carrying $`\widehat{S}_A`$s. In this case the $`\widehat{S}_B`$ can be incorporated in a tadpole with $`\{Kc^{}\}_𝐀`$, there being no other way to attach $`\widehat{S}_B`$ to the $`k`$-carrying $`\widehat{S}_A`$s, as follows by Corollary 4 and Proposition 5. But we have already shown that such a diagram is finite. $``$$``$ By Proposition 7, and Corollary 4 and Proposition 5, we are left only to prove the convergence of diagrams either (a) consisting only of $`\widehat{S}_A`$ and $`\{c^{}\}`$, or (b) where all the $`k`$-carrying wines are $`\{L\}`$s attaching to $`\widehat{S}_B`$s, either directly to the $``$s, or as in fig. 17. We complete the proof of Theorem 1, by showing that both these types vanish by the supertrace mechanism. As we will see, this is obvious for case (a). In case (b) the supertrace mechanism has been obscured by spontaneous symmetry breaking, the unitary gauge, and the ‘missing’ $`\sigma _3`$ described below (2.3). ###### Proposition 8 Diagrams of type (a) or (b) above, vanish identically. In case (a) the attachments are all made with full supermatrix differentials. There are thus no $`\sigma _3`$s generated and the internal Wilson loop in fig. 15 is $`\mathrm{str}\mathrm{\hspace{0.17em}1}=0`$. In constructing the diagram in case (b), if an $`\{L\}`$ is attached to a full $``$ at tree level (i.e. before the loop is closed) then the attachment is of the form $$\mathrm{str}d^Dw𝐗\left(\frac{\delta }{\delta 𝐂}+_\mu \frac{\delta }{\delta 𝒜_\mu }\right)\mathrm{str}d^Dx\left(𝐁_\nu +_\nu 𝐂\right)𝐘_\nu ,$$ where we have written out $`ł`$ by (2.3) and the $`_\nu `$ it attacks, using (2.3). Here $`𝐗`$ and $`𝐘_\nu `$ stand for the rest of the trees on either side. By Proposition 5, we are interested only in the case where the $`𝒜`$ differential attacks the lone $`𝐁_\nu `$. Thus using (2.1) and (2.1), the part of the above we are interested in, evaluates to $$\mathrm{str}d^Dx\left(_\nu 𝐗\mathrm{d}_{}𝐘_\nu +\mathrm{d}_+𝐗_\nu 𝐘_\nu \right)=\frac{1}{2}\mathrm{str}d^Dx\left(_\nu 𝐗\sigma _3𝐘_\nu \sigma _3\sigma _3𝐗\sigma _3_\nu 𝐘_\nu \right).$$ We see that we have exactly one $`\sigma _3`$ trapped on each side of the join as in the second term of fig. 2. As in fig. 2, the same local process takes place on supersplitting, thus if the $`\{L\}`$ attachment to $``$ closes the loop, a closely similar calculation shows that exactly one $`\sigma _3`$ gets trapped in each Wilson loop. For every occurence of fig. 17, there is a corresponding diagram with attachments as in fig. 21, where $`\{c^{}\}`$ first attaches to the covariant derivative in $`_\nu `$, i.e. to the term $`i𝐂𝒜_\nu `$ (by (2.1), (2.3) and the ordering implied in fig. 21).<sup>18</sup> Contributions fig. 17 and fig. 21 were discussed as fig. 24 of ref. . Fig.21. The partner diagram for Fig.17. Consider these attachments made at tree level, with $`\{L\}`$ and everything it attaches to at its other end being $`𝐗(w)`$, and the rest of $`\widehat{S}_B`$ and its attachments being $`𝐘_\nu (x)`$, as in (9.3), and $`\{c^{}\}`$ and all it attaches to being $`𝐙_\mu `$. By the description above Lemma 3, fig. 17 is given by $`i\mathrm{str}d^Dx\mathrm{𝐗𝐙}_\nu \mathrm{d}_{}𝐘_\nu `$. On the other hand by the description directly above, fig. 21 is given by $`i\mathrm{str}d^Dx𝐗\mathrm{d}_+\left(𝐙_\nu 𝐘_\nu \right)`$. Recognizing that, by Corollary 4, for BC bereft diagrams $`\mathrm{d}_+𝐙=𝐙`$, these contributions sum to $$i\mathrm{str}d^Dx𝐗𝐙_\nu \sigma _3𝐘_\nu \sigma _3.$$ As in (9.3), we have exactly one $`\sigma _3`$ trapped on either side of the join. Also as above, a similar calculation for the case where fig. 17 or fig. 21 closes the loop, shows that exactly one $`\sigma _3`$ gets trapped in each Wilson loop in this case. Therefore we see that in diagrams of type (b), an even number of $`\sigma _3`$s appear in the inner Wilson loop (one for every corner or join), which thus again furnishes $`\mathrm{str}\mathrm{\hspace{0.17em}1}=0`$. $``$$``$ 9.1. Divergent diagrams The above concludes the proof of finiteness of BC bereft one-loop diagrams. We now show by example where the difficulties lie in extending the ultraviolet finiteness of the present ‘unitary gauge’ formulation beyond this class, in particular we provide examples of divergent one-loop diagrams with external $`𝐁`$s, and divergent two-loop diagrams. The real difficulties appear to arise through the restriction to $`𝐀`$ of $`\{Kc^{}\}_𝐀`$ in (2.3). Thus we readily obtain one-loop divergent diagrams with remaindered $`𝐁`$s (cf. the Lemma 3 and the discussion before it) for example fig. 22. Note that since the attachment of $`L`$ to $`\{Kc^{}\}_𝐀`$ can only be made via $`𝐀`$, and the other end can only attach via $`𝐁`$ of $``$ (by fermion number conservation given that from (2.3), $`\{Kc^{}\}_𝐀`$ attaches via $`𝐁`$) there is no corresponding diagram that could cancel the divergence via the supertrace mechanism. Therefore some trapped $`\sigma _3`$s remain, after expanding as in fig. 2. Noting (5.5) and the Feynman rules, one readily shows that fig. 22 is quadratically divergent in $`D=4`$ dimensions (the integrand $`1/k^2`$ for large loop momentum $`k`$). Fig.22. Divergent diagram with a remaindered $`𝐁`$. We see that the $`𝐀`$ restriction destroys the supertrace mechanism as well as drifting. Another example is that of fig. 23. Since the $`c^{}`$ attachment has to be made via $`𝐀`$s there is no supertrace cancellation, and from the $`\widehat{S}_A`$ part of the RH lobe, we have again that the integrand $`1/k^2`$. (Note that the external $`𝐁`$s are forced via the $`𝐁`$ differential of the $`\{Kc^{}\}_𝐀`$ attachment.) Fig.23. Divergent diagram through restriction of supertrace mechanism. The problem is particularly severe beyond one loop. It is easy to give severe cases with external $`𝐁`$s, but now also BC bereft diagrams are affected, as illustrated in fig. 24. Here $`\{Kc^{}\}_𝐀`$ is attached to $`\widehat{S}_A`$ which is finite, but then a two-loop diagram is made by attaching $`L`$. $`L`$s ends will drift to the joins between $`\{Kc^{}\}_𝐀`$ and $`\widehat{S}_A`$ and cancel the diagrams where $`L`$ attaches and drifts from the other side of the join. However in the case that $`L`$ attaches via $`𝐁`$ this cancellation is imperfect because the corresponding diagram does not exist for $`L`$ attaching to $`\{Kc^{}\}`$. (Recall the arguments above Lemma 3. Equivalently we can start by noting that when $`L`$ attaches exclusively to $`\widehat{S}_A`$ the result vanishes by the supertrace mechanism.) The loop momentum routing round $`L`$ through $`\widehat{S}_A`$ gives a violently divergent integral with an integrand that diverges as $`k^{2\mathrm{\Delta }r}`$. Fig.24. Divergent BC bereft two-loop diagram. We cannot remove the $`𝐀`$ restriction however since the $`c^{}`$ in fig. 17 would then become $`K`$ and would not then cancel against fig. 21 as in Proposition 8 and ref. , resulting in quadratically divergent one-loop diagrams (integrand $`1/k^2`$) even in the BC bereft sector . Also if we removed the $`𝐀`$ restriction, $`\{Kc^{}\}`$ would bite its own tail as in fig. 3, and since it does so from (2.3) via the partial supermatrix $`𝐁`$, there is no supertrace cancellation in the inner Wilson loop, leading by (5.6), to an unregulated integrand $`1/k^2`$. (Here we note that the leading divergence would superficially arise from (5.5), but this vanishes under Lorentz invariance of the $`k`$ integral.) However this problem might need a separate solution as we explain below. Actually, tail-biting is a generic problem beyond one loop or for finite $`N`$. Recall that $`\{c^{}\}`$ did not bite its own tail because the inner Wilson loop in fig. 3 had to be field free in the large $`N`$ limit, thus giving $`\mathrm{str}\mathrm{\hspace{0.17em}1}=0`$. But at finite $`N`$, $`𝒜`$s can appear on the inner loop and thus $`\mathrm{str}\mathrm{\hspace{0.17em}1}`$ does not arise. Such a term could still be logarithmically divergent in $`D=4`$ dimensions (integrand $`1/k^4`$). At two-loops the subtleties of the large $`N`$ limit of the flow equation itself may allow such divergent tail biting diagrams even at $`N=\mathrm{}`$ as illustrated in fig. 25, providing neither field free Wilson loop vanishes by the supertrace mechanism. Fig.25. An $`N=\mathrm{}`$ divergent two-loop tail-biting diagram. For these reasons, extending these ideas to finite $`N`$ and beyond one-loop seem to require returning to more geometric covariantization (e.g. straight line wines) and regulating the wines as already described in so as to eliminate diagrams of form fig. 3. As we have seen, apart from this, divergences appear to arise only because of the restriction to $`𝐀`$ in $`\{Kc^{}\}_𝐀`$ which is necessary to ensure that the hidden supertrace mechanism described in Proposition 8 works at least in BC bereft one-loop diagrams. We expect that in a finite manifestly local $`SU(N\|N)`$ approach such problems will be absent. 10. Conclusions We have presented a methodology which for the first time allows manifestly gauge invariant continuum calculations to be performed. No gauge fixing or ghosts are required, and the full power and beauty of local gauge invariance is directly incorporated. We have demonstrated just a few of the consequences: the quantum gauge field does not renormalize – only the gauge coupling requires renormalization; possible counterterms are strongly constrained, easily and elegantly determined; the one-loop $`\beta `$ function falls out in a manifestly gauge invariant manner; the gauge invariant continuum Wilsonian effective action is for the first time precisely defined. It is important to recognize that the crucial problem solved in constructing a successful manifestly gauge invariant calculus in quantum field theory, is the passage to the continuum limit, equivalently renormalization. This is achieved by at the same time solving the long standing problem of combining gauge invariance with the exact RG. (In the exact RG, the continuum limit is almost automatic; in the gauge sector we only need to require a finite coupling constant $`g(\mathrm{\Lambda })`$ .) However, in order to write down such a gauge invariant exact RG, one must first solve the problem of finding a continuum gauge invariant physical cutoff. As we have seen from this paper and ref. , this is in a sense uniquely given by spontaneously broken $`SU(N\|N)`$ gauge theory with covariant higher derivatives. This method of regularization stands separately and is interesting and useful in its own right. It will be discussed more fully in ref. . We also comment further below. Our method should be compared with previous approaches to the exact RG in gauge theory which implement a physical cutoff at the expense of not only fixing, but breaking, the gauge invariance at least at intermediate stages , and it should be compared from this point of view with stochastic quantization (the manifestly gauge invariant form of which fails to regulate in the gauge orbit directions) and the famous loop-space, a.k.a. Migdal-Makeenko or Dyson-Schwinger, approach – as already indicated in refs. , where it is unknown how to formulate it at other than the bare level . At finite $`N`$, Mandelstam’s relations , which encode the overcomplete nature of the Wilson loop representation, also significantly complicate this approach since both equations are relations between intersecting loops. Here this overcounting has no great consequence, amounting to a harmless overparametrization of the effective action ; intersecting loops do not here carry any special significance. We stress that while it is possible and we believe useful to express the pure gauge sector and pure gauge sections geometrically in terms of Wilson loops and lines respectively, it is not necessary for our formulation. Nor need the apparent non-locality be any more than usual. Indeed if we choose $`c^1`$ and $`\stackrel{~}{c}^1`$ to be polynomials, then the seed action $`\widehat{S}`$ is just a sum over a finite number of ultralocal operators (i.e. vertices polynomial in their momenta) as would be normal for a bare action.<sup>19</sup> Despite some similarities $`\widehat{S}`$ is a renormalized action however. The effective action $`S`$ is necessarily only quasilocal (i.e. Taylor expandable in momenta) as is the case for any Wilsonian effective action. We should emphasise some subtleties in our construction. Even the inclusion of Pauli-Villars fields as part of the effective cutoff in an exact RG is novel, and has to be done carefully so that flow to lower energies amounts to lowering their masses (rather than raising them as naturally occurs for a relevant perturbation such as a mass) and thus corresponds to integrating out. We saw that at the classical level the solutions involve integrals over the effective cutoff $`\mathrm{\Lambda }`$ which diverge, requiring integration constant ‘counterterms’ or in effect a finiteness prescription for these $`\mathrm{\Lambda }`$ integrals. The Pauli-Villars contributions lead to finiteness in the momentum integrals, by subtracting separately divergent contributions. The answers are finite – but well defined only after applying and removing a preregulator. (In our case we took $`D<4`$ and only let $`D4`$ at the end of the calculation.) Only the symmetries respected in this procedure are guaranteed preserved. It is important in a regularisation such as this, that appears to generalise so straightforwardly to in particular anomalous gauge theories, that there are subtleties in its construction. For non-anomalous theories, an imperfect pre-regulator can be repaired by adding symmetry-breaking operators that vanish on removing the pre-regulator, but for anomalous theories this procedure will fail. From below (2.3), we also recall the wrong sign action for $`A_\mu ^2`$. As we noted below (2.1), and in ref. , this unphysical gauge field decouples from the physical gauge field $`A_\mu ^1`$ at energies much less than $`\mathrm{\Lambda }`$, but problems could potentially arise for $`A^1`$ if only at the non-perturbative level. At first sight the wrong-sign $`A^2`$ action leads to instability, however the fact that the kinetic term has also the wrong sign results instead in unitarity violations, whose effect in the physical sector should disappear in the limit that the overall cutoff is removed . We now understand the wrong sign $`A^2`$ action as an inevitable consequence of the underlying $`SU(N\|N)`$ gauge theory. The idea of regularising via such a gauge theory is another novel development in this paper. (For some previous applications of $`SU(N\|N)`$ gauge theory see refs.<sup>20</sup> The author thanks John Tighe for bringing these and ref. to his attention. .) Note that geometrically this corresponds to realising a superspace in the fibres of the principal bundle rather than in the base space as would correspond to the usual spacetime supersymmetry. As we saw in sec. 6.9, an intriguing property, in fact inherent to $`SU(N\|N)`$, is a duality that exchanges the rôle of $`A^1`$ and $`A^2`$ and thus also in a sense changes $`g^2`$ to $`g^2`$. The present ‘unitary gauge formulation’ closely imitates a spontaneously broken $`SU(N\|N)`$ theory but actually differs in some small details (as discussed in secs. 2 and 6.6), and it is not clear to us that this duality survives spontaneous breaking in a true local $`SU(N\|N)`$ invariant formulation. We must also be mindful that the present formulation only regulates to one loop with any number of external gauge fields. Nevertheless, it is intriguing to note that the duality symmetry, the local $`SU(N\|N)`$ and the global fermionic $`U(1)`$, do not commute with each other. In view of the intimate connection of $`g^2`$ with $`\sigma _3`$ as suggested by duality, which may be taken to imply that $`\sigma _3`$ appears with the same power as $`g^2`$, it seems that we should regard the coupling constant as the space-time independent field $`g^2\sigma _3`$. The action of local $`SU(N\|N)`$ then forces us to consider it a space-time dependent field. Perhaps these ideas hold the key to develop this duality at a deeper level. We saw however that the duality is at least obscured when we come to renormalize. To renormalize the $`A^2`$ sector we should introduce a separate coupling $`g_2`$ for $`A^2`$. (We expect this to be true also in the complete $`SU(N\|N)`$ theory.) A very intriguing scenario now arises as is clear from (6.2). The wrong sign action for $`A^2`$ is equivalent to the wrong sign $`g_2^2`$, and means that for negative $`\beta _1`$, $`g_2`$ is not asymptotically free but trivial. As $`\mathrm{\Lambda }_0\mathrm{}`$, the continuum limit for the physical Yang-Mills field $`A^1`$ is reached by $`g_00`$. But in this limit $`g_2`$ vanishes, and thus it would appear that for dynamical reasons the unphysical $`A_2`$ sector entirely decouples and loses all interactions in the continuum limit. Finally, let us remark that the formalism presented here appears to open many doors, and suggests many further avenues of exploration. We note that there is considerable freedom: the use of completeness relations for the generators, coincident lines, power law cutoff functions, or the same covariantization for the different kernels, is not necessary. Other interactions in $`\widehat{S}`$ are possible, and other ways of arranging the covariantization and regularisation may be considered along similar lines . This may help to understand more deeply holographic RG flow in the AdS/CFT correspondence and string theory . We have already touched on the issues of extending the regularisation to finite $`N`$ , of separately developing the $`SU(N\|N)`$ regularisation , of formulating a manifestly $`SU(N\|N)`$ invariant exact RG and generalisation to other groups. We do not directly compute correlators of gauge fields, the integration over modes being indirect through in effect iterating infinitessimal gauge invariant changes of variables . Nevertheless we expect that gauge invariant correlators can be computed by introducing sources for the appropriate gauge invariant operators and absorbing these sources as space-time dependent couplings. As $`\mathrm{\Lambda }0`$, all modes are integrated out and the partition function $`𝒵`$ should then be just proportional to $`\mathrm{e}^S`$. Differentiating with respect to the sources should then allow all the physics to be extracted. We should also understand if/how gauge variant operators are gauge averaged over. Other future directions include higher loop calculations such as $`\beta _2`$, investigating the Gribov problem in the continuum – perhaps through a limit as gauge fixing is removed, incorporation of matter fields in other representations; $`U(1)`$ gauge theory e.g. QED, through $`U(1\|1)`$, should be much simpler in many respects since the gauge field kernels do not need covariantizing . We have already touched on anomalies in this framework but these deserve further investigation. This formalism should allow investigations of instantons, renormalons and other controlled non-perturbative effects in a manifestly gauge invariant way. It looks possible to generalise this formalism to local coordinate invariance and thus a non-perturbative continuum framework for quantum gravity and supergravity. Of course the exact RG is tailor made to investigate Seiberg-Witten methods and super Yang-Mills theory more generally at a deeper level. And last but by no means least, we are hopeful that fully non-perturbative approximation methods can be developed to allow accurate analytic continuum calculations in realistic gauge theories, e.g. $`SU(3)`$ and QCD. Such approximation methods can involve for example the large $`N`$ (colour) limit . Many other potentially powerful approximation methods are also available within the exact RG framework . Acknowledgements The author wishes to thank Jonathon Evans, Hugh Osborn and John Tighe for perceptive comments, PPARC for financial support through an Advanced Fellowship and PPARC grant GR/K55738, and Ken Barnes for moral support. References relax T.R. Morris, Nucl. Phys. 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warning/0006/gr-qc0006003.html	ar5iv	text	# On Page’s examples challenging the entropy bound ## I Introduction In 1981 I proposed that the entropy of a complete physical system in three spatial dimensions whose total mass-energy is $`E`$, and which fits in a sphere of radius $`R`$, is necessarily bounded from above, $$S2\pi ER/\mathrm{}c,$$ (1) (henceforth I set $`\mathrm{}=c=1`$). The motivation comes from gedanken experiments in which an entropy bearing object is deposited in a black hole with a minimum of energy; a violation of the generalized second law seems to occur unless the bound applies to the object. A variant gedanken experiment in which the object is freely dropped into the hole again suggests that an entropy bound of the above form must be valid. A number of direct calculations of entropy of systems with no black hole involved have supported the validity of the bound for systems with negligible self-gravity. The entropy bound (1) is related to the holographic bound, and is often a stronger restriction on the entropy. Because of the intimate relation of information to entropy, bound (1) implies a fundamental restriction on the capacity of information storage and communication systems. This may become practically important in the distant foreseeable future (the same cannot be said of the vastly more lenient holographic information bound). Over the years Page has proposed a number of counterexamples to the entropy bound. In ref. I refuted the two from ref., while in ref. I showed Page’s proposed substitute for bound (1) to be violated. In the present paper I concentrate on Page’s recent candidates for counterexamples, and show that they do not constitute violations of the entropy bound as originally stated. It is common for critics of bound (1) to fail to include in $`E`$ the energy of some part of the system, contrary to the stipulations that the bound applies only to a complete system—one that could be dropped whole into a black hole—and that $`E`$ in the bare inequality (1) means the system’s gravitating energy (gravitating energy, of course, is immune to the well known arbitrariness of the zero of energy). Excluding part of the energy from $`E`$ empties the bound of meaning. Thus, an harmonic oscillator with sufficiently large mass (which makes the energy spacing of its levels small) manifestly violates the bound unless one includes the oscillator’s mass in $`E`$. Yet in all his new examples, Page drops from $`E`$ the energy of some part of the system. Page’s first example is a self-interacting scalar field with a symmetric double well potential, with the field vanishing at a certain spherical boundary of radius $`R`$. Page correctly points out that quantum tunnelling between the wells splits the classically degenerate ground states into a new ground state and an excited state a very small energy $`\mathrm{\Delta }E`$ above the ground state. Page hastily concludes that this violates bound (1) because the entropy of a mixed state built out of the two can reach $`\mathrm{ln}2`$, while $`2\pi R\mathrm{\Delta }E`$ can be very small in natural units ($`\mathrm{}=c=1`$). As I show in Sec. II, Page’s exclusion of the classical energy of the configuration (much larger than $`\mathrm{\Delta }E`$) is unjustified because the last contributes to the field’s gravitating energy. When $`E`$ includes it and the indispensable minimum wall energy, bound (1) is always obeyed. In all of his new examples Page excludes contributions to $`E`$ from the passive parts of the system. He quotes extensively from Schiffer’s and my paper as license for this procedure. The context is this. In the 1980’s, based on the many cases studied numerically, we had toyed with the idea that a strong form of bound (1) might apply, whereby $`S`$ and $`E`$ in inequality (1) refer solely to a collection of noninteracting massless quanta, with the contributions to $`E`$ from the cavity containing them ignored. We conjectured a theorem to this effect, which we proved for massless scalar quanta confined to a cavity of arbitrary topology by Dirichlet boundary conditions. We also sketched the proof of a generalization of that theorem for photons or massless neutrinos; we considered only the wave equation, but did not go into details about boundary conditions or constraint equations which are very specific for electromagnetism or fermion fields. Despite some efforts, we never completed these proofs. So in the intervening decade I have used bound (1) in its original formulation. Page’s new examples, which hark back to a version of the entropy bound which did not become established generically, teach us little about the true bound on entropy, and certainly do not constitute counterexamples to it. In his second example, an onion-like arrangement of a number of infinitely conducting concentric partitions separating electromagnetic fields in different states of excitation, Page establishes that some mixed states of the electromagnetic field disobey the strong bound $`S<2\pi (EE_{\mathrm{vac}})R`$, where $`EE_{\mathrm{vac}}`$ refers only to the electromagnetic energy above the vacuum, and $`R`$ to the radius of the outermost partition. But this is not a violation of Schiffer-Bekenstein theorem; the theorem (strictly proved only for scalar quanta) only claims that the ‘strong entropy bound’ is obeyed by a noninteracting field in a cavity, whereas in Page’s onion the electromagnetic field interacts throughout the volume of the sphere of radius $`R`$ with certain massive charged fields describing the conducting material of the partitions. My semiclassical discussion of the conducting partitions in Sec. III shows that even without specification of the nature of the charge carriers, the complete system does obey the original entropy bound (1). Page’s third example is a closed loop of coaxial cable of total length $`L`$ coiled up inside radius $`R`$, with $`RL`$, which confines an electromagnetic field. Page gives hand waving arguments that the lowest travelling frequencies are $`2\pi /L`$. Citing again the electromagnetic version of the Schiffer-Bekenstein theorem for license, Page estimates $`E`$ for the cable containing a mixed state based on the two lowest lying modes to be of order $`1/L`$. Since the state can have entropy of order unity, Page argues that the entropy bound is violated since $`RE1`$. If Page is correct about the low frequencies, then the Schiffer-Bekenstein theorem cannot be proved for electromagnetic fields confined to a waveguide with a not simply connected crossection. At any rate, the complete system includes the charged carriers which form the cable, and just as in the previous example, these contribute enough to $`E`$ to make bound (1) work (Sec. IV). Page revives Deutsch’s old objection to the entropy bound to the effect that weak thermal excitation of the system involves a violation. I give in Sec. V a brief argument showing why the entropy bound is always obeyed by a complete system, even at low temperature. The problem of entropy bound violations at low temperatures is evidently a red herring ! The above argument also illustrates how to deal with any mixed state which ascribes low probabilities to the high energy pure states. Page also revives the old argument that bound (1) will be violated when there is a virtually unlimited number of particle species in nature. I have earlier refuted this. In addition, as I discuss in Sec. VI, one knows today that the vacuum of quantum field theory is gravitationally unstable when many species of particles exist. Such instability makes Page’s second revived argument moot. ## II Nonlinear scalar field Page’s first example deals with a self-interacting scalar field with a multiwell potential, and considers configurations of the field which vanish at a certain boundary of radius $`R`$. For the symmetric double well potential Page correctly points out that the classically degenerate ground states of the field, each localized in one well, engender, by quantum tunnelling between the wells, a new ground state $`\psi _0`$ (energy $`E_0`$) with equal amplitude at each well and a first excited state $`\psi _1`$ a very small energy $`E_1E_0`$ above the ground state. He then hastily concludes that this violates the bound because the entropy of a mixed state built out of $`\psi _0`$ and $`\psi _1`$ can reach $`\mathrm{ln}2`$ (ground and excited states equally probable), while $`(E_1E_0)R`$ can be very small in natural units ($`\mathrm{}=c=1`$). In this interpretation Page considers the energy $`E`$ mentioned in the bound as the energy measured above the ground state. It would indeed be so if the ground state referred to a spatially unrestricted configuration, because then the bottom of the potential well would be the correct zero of energy (neglecting zero point fluctuations). But since the field is required to vanish at radius $`R`$, the energy $`E_0`$ of the described ground state is a function of $`R`$, and it makes little sense to take it as the zero of energy. For example, by expanding the system can do work ($`E/R0`$), so that its gravitating energy changes, and cannot be taken as zero for all $`R`$. The gravitating energy for the equally likely mixture of ground and excited states should be identified with $`\frac{1}{2}(E_0(R)+E_1(R))E_0(R)`$, which does go to zero as $`R\mathrm{}`$, and is, therefore, properly calibrated with respect to the global ground states of the theory.. The exponential smallness of $`E_1E_0`$, which Page pounced upon, is not very relevant, as will transpire. Since there are no solitons in $`D=1+3`$ spacetime, a finite sized field configuration \[the only interesting case—see (1)\] must be confined by a “wall”. There are three parts to the energy $`E`$ of the total system (before tunnelling is taken into account): the classical energy $`E_c`$ of the field configuration concentrated around one well but vanishing at radial coordinate $`r=R`$, the quantum correction $`E_v`$ due to the zero point fluctuations about the classical configuration, and $`E_w`$, the energy of the “wall” at $`r=R`$. As I show below, $`E_w`$ is at least of the same order as $`E_c`$, and both dominate Page’s energy. ### A Classical two-well configurations The double well potential field theory comes from the lagrangian density $$=\frac{1}{2}_\mu \varphi ^\mu \varphi \frac{1}{4}\lambda (\varphi ^2\varphi _m^2)^2.$$ (2) This gives the field equation $$_\mu ^\mu \varphi \lambda \varphi (\varphi ^2\varphi _m^2)=0.$$ (3) Every spherically symmetric configuration inside a spherical box of radius $`R`$ will thus satisfy (I use standard spherical coordinates; denotes derivative w.r.t. to $`r`$) $$r^2(r^2\varphi ^{})^{}\lambda \varphi (\varphi ^2\varphi _m^2)=0.$$ (4) Regularity requires that $`\varphi ^{}=0`$ at $`r=0`$. Page chooses $`\varphi =0`$ at $`r=R`$. The energy of such a configuration will be $$E_c=\frac{1}{2}_0^R\left[\varphi ^2+\frac{1}{2}\lambda (\varphi ^2\varphi _m^2)^2\right]r^2𝑑r.$$ (5) Since one is interested in the ground state, I require that $`\varphi `$ have its first zero at $`r=R`$. Multiplying Eq. (4) by $`r^2\varphi `$ and integrating over the box allows, after integration by parts and use of the boundary conditions, to show that $$_0^R\varphi ^2r^2𝑑r=\lambda _0^R(\varphi _m^2\varphi ^2)\varphi ^2r^2𝑑r$$ (6) whereby $$E_c=\frac{1}{4}\lambda _0^R(\varphi _m^4\varphi ^4)r^2𝑑r.$$ (7) It proves convenient to adopt a new, dimensionless, coordinate $`x\sqrt{\lambda \varphi _mr}`$ and a dimensionless scalar $`\mathrm{\Phi }\varphi /\varphi _m`$. Then Eq. (4) turns into a parameter-less equation: $$\frac{1}{x^2}\frac{d}{dx}\left(x^2\frac{d\mathrm{\Phi }}{dx}\right)+\mathrm{\Phi }(1\mathrm{\Phi }^2)=0.$$ (8) Using $`d\mathrm{\Phi }/dx=0`$ at $`x=0`$ one may integrate the equation to get $$\frac{d\mathrm{\Phi }}{dx}=\frac{1}{x^2}_0^x\mathrm{\Phi }(1\mathrm{\Phi }^2)x^2𝑑x.$$ (9) If the integration starts with $`\mathrm{\Phi }(0)>1`$, then by continuity the r.h.s. here is positive for small $`x`$, so that $`\mathrm{\Phi }`$ grows. There is thus no way for the r.h.s. to switch sign, so $`\mathrm{\Phi }(x)`$ is monotonically increasing and can never have a zero. If $`\mathrm{\Phi }(0)=1`$, it is obvious that the solution of Eq. (9) is $`\mathrm{\Phi }(x)1`$ which cannot satisfy the boundary condition at $`r=R`$. Thus the classical ground state configuration we are after requires $`\mathrm{\Phi }(0)<1`$. When $`\mathrm{\Phi }(0)<1`$ it can also be seen from Eq. (9) that $`\mathrm{\Phi }`$ is monotonically decreasing with $`x`$. For a particular $`\mathrm{\Phi }(0)`$, $`\mathrm{\Phi }(x)`$ will reach its first zero at a particular $`x`$ which I refer to as $`x_0`$. This can serve as the parameter singling out the solution in lieu of $`\mathrm{\Phi }(0)`$. One thus has a family of ground state configurations $`\mathrm{\Phi }(x,x_0)`$. Each such configuration corresponds to a box of radius $`R=x_0(\sqrt{\lambda \varphi _m})^1`$. In terms of the new variables one can write Eq. (7) as $$E_c=\frac{x_0}{4\lambda R}_0^{x_0}(1\mathrm{\Phi }^4)x^2𝑑x.$$ (10) The dependence $`E_c\lambda ^1`$ is well known from kink solutions of (3) in $`D=1+1`$, where the role of $`R^1`$ is played by the effective mass of the field. Numerical integration of Eq. (8) shows that the factor $`x_0_0^{x_0}(1\mathrm{\Phi }^4)x^2𝑑x`$ grows monotonically from 32.47 for $`\mathrm{\Phi }(0)=0(x_0=3.1416)`$ to 232.23 for $`\mathrm{\Phi }(0)=0.98(x_0=5.45)`$ to infinity as $`\mathrm{\Phi }(0)1(x_0\mathrm{})`$. Since $`E_c`$ is not exponentially small, the quantum tunnelling corrections that Page discussed are negligible, so one need only add to $`E_c`$ the zero point fluctuations energy $`E_v`$ plus the wall energy $`E_w`$ to get the full energy of the ground state, $`E_0`$. This plays the role of $`E`$ in the bound (1). I shall not bother to calculate $`E_v`$. This can be done by present techniques only for the weak coupling case $`\lambda <1`$. It is then found in other circumstances, e.g. the $`D=1+1`$ kink, that $`E_v`$ is small compared to $`E_c`$. The situation for large $`\lambda `$ (the strong coupling regime) is unclear. However, it is appropriate to recall here that the theory (2) is trivial in that it makes true mathematical sense only in the case $`\lambda =0`$. Theorists use it for $`\lambda 0`$ to obtain insights which are probably trustworthy in the small $`\lambda `$ regime, but probably not for large $`\lambda `$. I now set a lower bound on $`E_w`$. A look at Eq. (10) shows that for $`\mathrm{\Phi }(0)1`$ and so $`\mathrm{\Phi }(x)1`$), $`E_c`$ scales as $`x_0^4/RR^3`$. Numerically the exponent here only drops a little as $`\mathrm{\Phi }(0)`$ increases; for example, it is 2.86 for $`\mathrm{\Phi }(0)=0.98`$. So I take it as 3. On virtual work grounds (consider expanding $`R`$ a little bit) the $`R^3`$ dependence means the $`\varphi `$ field exerts a suction (negative pressure) of dimension $`(3E_c/4\pi R^3)`$ on the inner side of the wall. By examining the force balance on a small cap of the wall, one sees that in order for the wall to withstand the negative pressure, it must support a compression (force per unit length) $`\tau (3E_c/8\pi R^2)`$. Under this compression vibrations on the wall will propagate superluminally unless the surface energy density is at least as big as $`\tau `$ (dominant energy condition). Thus one may conclude that the wall (area $`4\pi R^2`$) must have (positive) energy $`E_w>3E_c/2`$ which adds to $`E_c`$ to give $`E>5E_c/2`$. For $`\mathrm{\Phi }(0)`$ very close to unity the coefficient is somewhat lower than $`5/2`$. However, by then $`E_cR`$ is already much larger than the corresponding quantity for $`\mathrm{\Phi }(0)1`$ (six times larger for $`\mathrm{\Phi }(0)=0.98`$). Using the value of $`E_c`$ for $`\mathrm{\Phi }(0)1`$ from the preceding argument, I thus conclude that for all physically relevant $`\mathrm{\Phi }(0)`$, $`2\pi ER>127.5\lambda ^1`$. This is certainly not exponentially small with $`R`$ as Page claimed ! True, formally it seems possible to have a violation of the bound on the entropy $`\mathrm{ln}2`$ of the 50% mixture of ground and excited states whenever $`\lambda >127.5/\mathrm{ln}2=183.95`$. However, this is the strong coupling regime. For all one knows the zero point energy may then become important and tip the scales in favor of the entropy bound. At any rate, because the theory (2) is trivial, one is more likely to be overstepping here the bounds of its applicability than to be witnessing a violation of the entropy bound at large $`\lambda `$. In summary, we have found that whenever the calculation is meaningful ($`\lambda `$ not large), the entropy bound (1) is satisfied provided $`E`$ includes all contributions to the energy. In his later defense against this observation, Page cites my paper with Schiffer as an excuse for including in $`E`$ just the excitation energy above the classical ground state. However, he neglects to point out that we ourselves restricted use of this approach to an assembly of quanta of a massless noninteracting field. I should also mention that in $`D=1+1`$ spacetime it is possible to find analytically all static classical configurations (and their energies) for the theory (2) in a box, and the distribution of energy levels is such that the entropy bound is sustained. ### B Multiwell potential Page also confronts bound (1) with a theory like (2) but with a potential having three equivalent wells. Pressumably one would like one of these centered at $`\varphi =0`$, with the other two flanking it symmetrically. Then Page’s conclusion that there are three exponentially close states (in energy) is untenable. This would require three classically degenerate configurations, which certainly exist in open space (field $`\varphi `$ fixed at one of three well bottoms). However, one is here considering a finite region of radius $`R`$ with $`\varphi =0`$ on the boundary. One exact solution is indeed $`\varphi =0`$, and it has zero energy (the zero point fluctuation energy correction will, however, depend on $`R`$). Then there are two degenerate solutions in which the field starts at $`r=0`$ in one side well and then moves to the central one with $`\varphi 0`$ as $`rR`$. By analogy with our earlier calculations, the common classical energy of these two configurations will be of $`O\left(8(\lambda R)^1\right)`$. It cannot thus be regarded as the zero of energy; this role falls to the energy of the $`\varphi =0`$ configuration. When tunnelling between wells is taken into account, one has a truly unique ground state and two excited states of classical origin split slightly in energy (plus the usual gamut of quantum excitations). The entropy of an equally weighted mixture of these states is $`\mathrm{ln}3`$. The mean energy $`E`$ is $`\frac{2}{3}(E_c+E_w)`$ of an excited state, that is $`E=O(10(\lambda R)^1)`$, so the entropy bound is satisfied, at least in the weak coupling regime where the theory makes sense. When the potential has $`n=5,7,9,\mathrm{}`$ equivalent wells with one centered at $`\varphi =0`$ and the rest disposed symmetrically about it, there will be a single zero-energy configuration ($`\varphi 0`$), and $`\frac{1}{2}(n1)`$ pairs of degenerate configurations with succesively ascending $`R`$-dependent energies. For $`n=4,6,8,\mathrm{}`$ wells there is no zero-energy configuration, but there are $`\frac{1}{2}n`$ pairs of degenerate configurations with $`R`$ dependent energies. Because of the extra energy splitting appearing here already classically, I expect, by analogy with the previous results, that the appropriate mean configuration energy (perhaps supplemented by wall energy), when multiplied by $`2\pi R`$, will bound the $`\mathrm{ln}n`$ entropy from above. ## III Electromagnetic Onion As a second counterexample to the entropy bound, Page proposes a sphere of radius $`R`$ partitioned into $`n`$ concentric shells; the partitions and the inner and outer boundaries are regarded as infinitely conducting. He points out that the lowest ($`\mathrm{}=1`$) three magnetic-type electromagnetic modes in the shell of median radius $`r`$ have frequency $`\omega 1/r`$. Since there are $`3n`$ such modes (three for each shell), Page imagines populating now one, then another and so on with a single photon of energy $`1/r`$ for the appropriate $`r`$. These one-photon states allow him to form a density matrix which, for equally weighted states, gives entropy $`\mathrm{ln}(3n)`$ and mean energy $`2/R`$ (since $`R/2`$ is the median radius of the shells if they are uniformly thick). Page concludes that bound (1) is violated because the entropy grows with $`n`$ but the mean energy does not. Page has again missed out part of the energy. The modes he needs owe their existence to the infinitely conducting partitions that confine them, each to its own shell. Were passage between shells possible, then old work already established that the entropy bound works for the electromagnetic field confined to an empty sphere (or for that matter confined to any parallelepiped. To be highly conducting, the envisaged partitions must contain a certain number of charge carriers, whose aggregated masses turn out to contribute enough to the system’s total energy $`E`$ to make it as large as required by the entropy bound (1). Ignoring the masses of the charge carriers goes against the condition that the bound applies to a complete system: the carriers are an essential component, so their gravitating energy has to included in $`E`$. I assume all partitions to have equal thickness $`d`$. One mechanism that can block the waves from crossing a partition is a high plasma frequency $`\omega _p`$ of the charge carriers in the partitions. We know that in a plasma model of a conductor with collisionless charge carriers, the electromagnetic wave vector for frequency $`\omega `$ is $`k=\omega (1\omega _p^2/\omega ^2)^{1/2}`$, so that if $`\omega <\omega _p`$, the fields do not propagate as a wave. Nevertheless they do penetrate a distance $`\delta =\omega ^1(\omega _p^2/\omega ^21)^{1/2}>\omega _p^1`$ into the plasma before their amplitudes become insignificant. In order to prevent these evanescent waves from bridging a partition, one must thus require $`\delta <d`$, i.e., $`\omega _pd>1`$. But $$\omega _p{}_{}{}^{2}=4\pi 𝒩e^2/m$$ (11) where $`𝒩`$ is the density of charge carriers of charge $`e`$ and mass $`m`$. Since $`d<R/n`$, all this gives us $`(4\pi R^2d)𝒩>mn^2d/e^2`$. Now $`4\pi R^2d`$ is the volume of material in the outermost partition. Properly accounting for the variation of partition area with its order $`i`$ in the sequence (we employ the sum $`i^2`$), tells us that for $`n1`$ the total mass-energy in charge carriers in all the partitions is $`Enm(4\pi R^2d/3)𝒩`$. Substituting our previous bound on $`(4\pi R^2d)𝒩`$, I get $`E>\frac{1}{3}n^3m^2d/e^2`$. Now a charge carrier’s Compton length has to be smaller than $`d`$, for otherwise the carriers would not be confined to the partitions; hence $`md>1`$. As a matter of principle $`e^2<1`$, because more strongly coupled electrodynamics makes structures, e.g. atoms and partitions, which are held together electrically, unstable (in our world $`e^2<10^2`$). Therefore, since $`d<R/n`$ one gets $`E>\frac{1}{3}n^4R^1`$, which strongly dominates $`2/R`$, the energy in photons that Page gets. In particular $`2\pi RE>2n^4`$, which is always much larger than the entropy in photons $`\mathrm{ln}(3n)`$. The only alternative mechanism for keeping electromagnetic waves from penetrating into a conductor is the skin effect. The skin depth is $`\delta (2\pi \omega \sigma )^{1/2}`$, where $`\sigma `$ is the conductivity. In the simple Drude model, $`\sigma =𝒩e^2(m/\tau ım\omega )^1`$, where $`\tau `$ is the slowing-down timescale for a charge carrier due to collisions, and $`ı=\sqrt{1}`$. The formula for $`\delta `$ refers to an Ohmic (real) conductivity rather than an inductive (imaginary) one. Thus one must demand that $`\omega 1/\tau `$. But then $$\delta (2\pi 𝒩e^2/m)^{1/2}.$$ (12) As before one must require $`\delta <d<R/n`$. This gives $`(4\pi R^2d)𝒩\frac{1}{2}mn^2d/e^2`$ which is just a stronger version of the lower bound on $`𝒩`$ we got before. Repeating the previous discussion verbatim shows that $`2\pi REn^4`$, which bounds Page’s $`\mathrm{ln}(3n)`$ entropy confortably. In conclusion, bound (1) is satisfied by the system photons $`+`$ charge carriers. It should be clear that the entropy of the conducting material with its many carriers, which we have been ignoring, may well dominate that in photons. Here one can fall back on the usual argument that in a random assembly of particles the entropy is of order of the number of particles, and that each particle’s Compton length is necessarily smaller than the system’s radius. From these two conditions it follows that the entropy bound applies—with room to spare—to the charge carrier system by itself. Putting all this together makes it clear that it applies to the complete onion system as well. ## IV Coaxial cable loop Page’s third example is furnished by an electromagnetic field confined to a coaxial cable of length $`L`$ which is coiled up so as to fit within a sphere of radius $`R`$, with $`RL`$, before being connected end to end to form a closed loop. Page’s entirely qualitative reasoning proceeds by analogy with a rectilinear coaxial cable with periodic boundary conditions. A rectilinear infinitely long coaxial cable has some electromagnetic modes which propagate along its axis with arbitrarily low frequency. Page notes that for the coiled-up cable, each right moving mode is accompanied by a degenerate (in frequency) left moving mode (basically this follows from time reversal invariance of Maxwell’s equations). He then argues that if the cable’s outer radius $`\varrho `$ is thin on scale $`R`$, the structure of the electromagnetic modes is little affected by the cable’s curvature. This leads him to estimate the lowest frequency $`\omega _1`$ as similar to that of the rectilinear coaxial cable with periodic boundary conditions with period $`L`$: $`\omega _12\pi L^1`$. Page then imagines a mixed electromagnetic state of energy $`EE_{\mathrm{vac}}=2\pi L^1`$ in which a single photon occupies the right- or the left-moving mode of frequency $`\omega _1`$ with equal probabilities $`\frac{1}{2}`$. The entropy of this state is $`\mathrm{ln}2`$. However, $`2\pi (EE_{\mathrm{vac}})R4\pi ^2(R/L)`$ which could be very small compared to $`\mathrm{ln}2`$. Therefore, Page exhibits this example as a violation of the Schiffer-Bekenstein ‘strong entropy bound’. As I stressed in Sec.I, the strong bound was never formalized in a theorem for electromagnetism. If Page is correct in his estimate of $`\omega _1`$, such a theorem cannot apply to a cavity with not simply connected crossection, like the coaxial cavity. (However, the myriad examples studied numerically strongly suggest that a theorem of the desired type should exist for simply connected cavities.) At any rate, in the present example the interesting question is whether the coaxial cavity plus electromagnetic field obeys the original entropy bound (1). The inner conductor of the cable—let its outer radius be $`\rho `$ and its thickness $`d`$—is an essential part of the system, for without it the lowest propagating frequency would be $`O(\varrho ^1)`$, i.e., very large on scale $`L^1`$. Now in order for the inner conductor to keep the fields out of it, as required by the whole notion of a coaxial cable, one must have either $`\rho >d>\omega _p^1`$ or $`\rho >d>\delta `$ (see Sec. III). In the case $`\rho >d>\omega _p^1`$, Eq. (11) informs us that $`𝒩\rho d>m(4\pi e^2)^1`$ where, as in Sec. III, $`m`$ denotes a charge carrier’s mass and $`𝒩`$ the carrier density. Since the volume of material in the inner conductor is $`\pi L[\rho ^2(\rho d)^2]`$, its mass energy is at least $`\pi mL𝒩\rho d`$, and thus the total mass-energy $`E`$ of cable plus field is constrained by $`E>m^2L(4e^2)^1`$. But a charge carrier has to be localized within the conductor, which requires that $`md>m\rho 1`$. Hence $`EL(4e^2\rho ^2)^1`$ so that $`2\pi ER(\pi /2e^2)(L/\rho )(R/\rho )`$. Now obviously $`R>\varrho >\rho `$ and $`LR`$ by the conditions of the problem, while $`e^2<1`$ by the condition of stability (see Sec. III), so that $`2\pi ER1`$. Thus the bound on entropy confortably bounds Page’s entropy $`\mathrm{ln}2`$. Were the mixed state instead to involve one photon in any one of the low lying modes having no crossectional nodes and wavelength along the cable of the form $`k(L/2\pi )`$ with $`k`$ an integer, one could get a bigger entropy. There are $`O(L/\rho )`$ such modes with frequency below that of the lowest lying transversally excited mode (which is obviously of order $`\rho ^1`$), so the entropy of the envisaged state is $`\mathrm{ln}(L/\rho )`$. Because $`\mathrm{ln}x<x`$ for $`x>1`$, this is obviously bounded by $`(L/\rho )(R/\rho )`$ and, therefore, by $`2\pi ER`$. In the case $`\rho >d>\delta `$, Eq. (12) gives $`𝒩\rho dm(2\pi e^2)^1`$. This is just a stronger version of the earlier lower bound on $`𝒩`$. Repeating the previous discussion shows again that $`2\pi ER`$, with $`E`$ the total energy of the system, bounds the entropy. ## V Low temperature systems Page also examined the generic example in which the system’s density matrix is diagonal and involves $`g+1`$ equally probable pure states. The entropy is $`\mathrm{ln}(g+1)`$. Obviously the greatest challenge to bound (1) is posed when $`g`$ of the states are degenerate and just a small energy $`\mathrm{\Delta }`$ above the (unique) ground state whose energy is $`ϵ_0`$. Now for any $`\mathrm{\Delta }`$ and $`g`$ the mean energy $`E`$ of the complete system is at least $`ϵ_0`$, so $`2\pi RE>2\pi Rϵ_0`$. The present system is so generic that it is not feasible to estimate $`ϵ_0`$ as I did previously. It should be clear, however, that the system’s longest possible Compton wavelength, $`ϵ_0^1`$, should lie well below $`R`$; otherwise the contention that the system fits within a definite radius $`R`$ would make no sense as it would be poorly localized on scale $`R`$. It is probably conservative to take $`Rϵ_0>3`$. Therefore, no violation of the entropy bound can occur for $`g<10^8`$. Can $`g`$ be bigger ? Now quantum mechanical systems have low degeneracies in the low lying levels. Quantum field systems have more. But even in those systems with accidental degeneracies, e.g. electromagnetic field in a cubical box, $`g<10`$. As the end of Sec.II B illustrates, it is hard to generate high degeneracies in nonlinear fields in a box. One could create a highly degenerate state artificially by lumping together states closely spaced in energy, e.g. the one-photon states using the modes with $`k=1,2,\mathrm{}`$ of the coaxial cable. But as shown in Sec.IV, the large $`L/\rho `$ required to have many of these modes defeats the attempt to keep $`2\pi RE`$ small because a long cable involves a lot of charge carriers. This illustrates the point that, unlike Page, one cannot arbitrarily legislate a high degeneracy. Rather, one must examine the energy cost paid by the passive components of the system in providing a large number of closely spaced states for the active part which can be consolidated into a formally highly degenerate one. In a related line of thought, Page revives an old challenge to the entropy bound, which is occasionally discovered anew: a system in a thermal state seems to violate the entropy bound if its inverse temperature $`\beta `$ is sufficiently high. The essence and resolution of the problem is captured by the following purely analytical treatment. Consider a system of radius $`R`$ with ground energy $`ϵ_0`$, a $`g`$-fold degenerate excited state at energy $`ϵ_1=ϵ_0+\mathrm{\Delta }`$, and higher energy states. For sufficiently large $`\beta `$ one may neglect the higher energy states in the partition function $`Z=_i\mathrm{exp}(\beta ϵ_i)`$, and so approximate it by $`\mathrm{ln}Z\beta ϵ_0+\mathrm{ln}(1+ge^{\beta \mathrm{\Delta }})`$. The mean energy is $$E=\frac{\mathrm{ln}Z}{\beta }=ϵ_0+\frac{g\mathrm{\Delta }}{e^{\beta \mathrm{\Delta }}+g}$$ (13) while the entropy takes the form $$S=\beta E+\mathrm{ln}Z=\frac{g\beta \mathrm{\Delta }}{e^{\beta \mathrm{\Delta }}+g}+\mathrm{ln}(1+ge^{\beta \mathrm{\Delta }}).$$ (14) The typical claim is“measure energies from the ground state so that $`ϵ_0=0`$; then for $`\beta >2\pi R`$ one gets $`S>2\pi RE`$ and the bound is violated”. But we have already seen in Sec. II that taking the zero of energy of a system at its ground state is not automatically justified because it may mean using as $`E`$ something distinct from the gravitating energy. And I have already discussed why $`Rϵ_0`$ should be at least of the order a few, say 3. The interesting quantity now is $$S2\pi RE=\mathrm{\Xi }(\beta \mathrm{\Delta })\frac{g(\beta \mathrm{\Delta }2\pi R\mathrm{\Delta })}{e^{\beta \mathrm{\Delta }}+g}+\mathrm{ln}(1+ge^{\beta \mathrm{\Delta }})2\pi Rϵ_0.$$ (15) The function $`\mathrm{\Xi }(y)`$ has a single maximum at $`y=2\pi R\mathrm{\Delta }`$ where $`\mathrm{\Xi }=\mathrm{ln}(1+ge^{2\pi R\mathrm{\Delta }})`$. I thus conclude that $$S<2\pi RE+[\mathrm{ln}(1+ge^{2\pi R\mathrm{\Delta }})2\pi Rϵ_0].$$ (16) For the quantity in square brackets to be nonnegative it would be necessary for $`ge^{2\pi R\mathrm{\Delta }}[e^{2\pi Rϵ_0}1]`$, i.e., $`g>10^8`$. As we saw above, this cannot be arranged. Thus the quantity in square brackets in Eq. (16) has to be negative: for sufficiently low temperature the entropy bound is upheld with room to spare. The above argument also illustrates how to deal with any mixed state which ascribes low probabilities to the high energy pure states. Early realistic numerical calculations of thermal quantum fields in boxes did reveal that, were the ground state energy to be ignored, the bound on entropy would be violated at very low temperatures, typically when $`E<10^9R^1`$ ($`R`$ enters through the “energy gap” $`\mathrm{\Delta }`$). It was also clear early that taking any reasonable ground state energy into account precludes the violation. As the temperature rises, more and more pure states are excited, and eventually $`S/E`$ peaks and begins to decrease. In this regime the entropy bound is always obeyed regardless of whether or not one includes $`ϵ_0`$ in the total energy. ## VI Proliferation of species Page also revives the old “proliferation of species” challenge to the entropy bound. Suppose there were to exist as many copies $`N`$ of a field e.g. the electromagnetic one, as one ordered. It seems as if the entropy in a box containing a fixed energy allocated to the said fields should grow with $`N`$ because the bigger $`N`$ is, the more ways there are to split up the energy. Thus eventually the entropy should surpass the entropy bound. Numerical estimates show that it would take $`N10^9`$ to do the trick. A similar picture seems to come from Eq. (16); the degeneracy factor $`g`$ should scale proportionally to $`N`$ making the factor in square brackets large, so that, it would seem, one could not use the argument based on (16) to establish that $`S<2\pi RE`$. However, as recognized in refs. , the above reasoning fails to take into account that each field species makes a contribution of zero point fluctuations energy which gets lumped in $`ϵ_0`$. If these contributions are positive, then the negative term in the square bracket eventually dominates the logarithm as $`N`$ grows, and for large $`N`$ one again recovers the entropy bound. If they are negative (which implies a Casimir suction proportional to $`N`$ on the walls which delineate the system), then the scalar field example suggests that the wall energy, which must properly be included in $`ϵ_0`$, should suffice to make the overall $`ϵ_0`$ positive. Again the entropy bound seems safe. There is an alternative view: the seeming clash between entropy bound and a large number of species merely tells us that physics is consistent only in a world with a limited number of species, such as the one that is observed. Indeed, as Brustein, Eichler and Foffa have argued, a large number of field species will make the vacuum of quantum field theory unstable against collapse into a “black hole slush”. Thus the proliferation of species argument against the entropy bound is not even physically consistent. ## VII Caveats I have stressed the robustness of the entropy bound for weakly gravitating systems. But one should recall that the bound has its limitations. These principally belong to the strongly gravitating system regime. Bound (1) does not apply in wildly dynamic situations such as those found inside black holes, and it is not guaranteed to work for large pieces of the universe (which, after all, are not complete systems). ###### Acknowledgements. This work was supported by the Hebrew University’s Intramural Research Fund. I thank Eduardo Guendelman for much information and Mordehai Milgrom for incisive critical remarks.
warning/0006/gr-qc0006010.html	ar5iv	text	# Entropy gap and time asymmetry II. ## I Introduction. In papers <sup></sup><sup></sup>It is a good occasion to correct some errata in the abstract of paper . -In line 2: it reads ”quantitative”, it must read ”qualitative”. -In line 3: it reads ”entropy”, it must read ”entropy gap”. -In line 4: it reads ”qualitative”, it must read ”quantitative”. and two of us reported a rough coincidence between the time where the minimum of the entropy gap $`\mathrm{\Delta }S=S_{act}S_{\mathrm{max}}`$ , takes place and the time where all the stars will exhaust their fuel. The rough nature of this result was explained saying that the calculation was done using a model extremely naive and simplified: a homogeneous universe Besides that the higher order terms in eq. (3) of ref. were neglected; perhaps they may be important for finite times.. In fact in the real universe nuclear reactions (that were considered as the main source of entropy) take place within the stars, that can only be properly considered in an inhomogeneous geometry. So we expend some time trying to generalize our model to these type of geometries and computing the entropy of the galaxies and stars in the condensation period, after decoupling time, in order to make the coincidence precise In fact, inhomogeneity must be considered, but in a different way, as we will see below.. It was soon clear to us that the latter entropy can be neglected with respect to the amount of entropy produced within the stars and therefore the mechanism that produced the inhomogeneity is not so important for our problem. Stars are formed when a sufficiently large mass of interstellar gas is compressed into a small enough volume, so its force of self-gravitation becomes sufficiently great to cause gravitational collapse. The instability makes the cloud to reduce its size quickly, the temperature rises and pressure forces begin to restrict the collapse until hydrostatic equilibrium is obtained. In this process, gravitational energy is converted in kinetic energy and radiation. In the classical model of star formation the time of contraction is the Kelvin-Helmholtz time: $$t_{KH}=\frac{GM^2}{RL}$$ (1) where $`M`$ is the mass, $`R`$ is the radius and $`L`$ is the luminosity of the star. For a star like the sun this time is $`3`$ $`10^7`$ years . In more recent models of star formation, the relevant time scale is the much shorter free-fall time : $$t_{ff}=(\rho G)^{1/2}=710^5\left(\frac{n}{10^4}\right)^{1/2}years$$ (2) where $`\rho `$ and $`n`$ are the cloud mass density and the cloud number density respectively. In any case, it takes less than $`1\%`$ of the lifetime $`(Gyr)`$ to the star to form. In the contraction to make the star, the energy radiated is a part of the gravitational energy and the nuclear reactions are not important. The binding gravitational energy is negligible compared with the energy liberated through nuclear reactions when the star is already formed . As a consequence we can neglect the entropy production in the phase of contraction that makes the star. So, in order to improve our rough result, we were forced to change the approach. Let us review the main formulae. The entropy gap was the conditional entropy of $`\rho (t)`$, the state of the universe at time $`t`$ with respect to the equilibrium state at that time $`\rho _{}(t)`$ : $$\mathrm{\Delta }S=tr[\rho \mathrm{log}(\rho _{}^1\rho )]$$ (3) where $$\rho (t)=\rho _{}(t)+\rho _1\mathrm{exp}(\gamma t)+O[(\gamma t)^1]$$ (4) and $$\rho _{}(\omega )=ZT^3\frac{1}{\mathrm{exp}(\frac{\omega }{T}1)}$$ (5) $`\rho _1`$ was a phenomenological coefficient constant in time, the time variation of the main irreversible process was $`\mathrm{exp}(\gamma t)`$ so $`\gamma ^1`$ was the characteristic time of this process, $`T`$ was the temperature of the universe, $`Ta^1`$ where $`a`$ is the radius or scale of the universe, and $`Z`$ is a normalization constant. Then we could approximate $`\mathrm{\Delta }S`$ as $$\mathrm{\Delta }S=CT^3\mathrm{exp}(\gamma t)\mathrm{exp}(\frac{\omega _1}{T})$$ (6) where $`\omega _1`$ was the characteristic energy where $`\rho _1`$ is peaked. The time where the minimum of $`\mathrm{\Delta }S`$ is located was: $$t_{cr}t_0\left(\frac{2}{3}\frac{\omega _1}{T_0}\frac{t_{NR}}{t_0}\right)^3$$ (7) We selected numerical values for the parameters: $`\omega _1=T_{NR},`$ the temperature of the nuclear reactions within the stars (that was used in paper as the main source of entropy), $`t_{NR}=\gamma ^1`$ the characteristic time of these nuclear reactions, $`t_0`$ the age of the universe, and $`T_0,`$ the cosmic micro-wave background temperature. $`T_{NR}`$ and $`t_{NR}`$ were chosen between the following values : $$T_{NR}=10^6\text{ }to\text{ }10^8\text{ }K$$ (8) $`t_{NR}=10^6\text{ }to\text{ }10^9\text{ }years`$ while for $`t_0`$ and $`T_0`$ we can take: $$t_0=1.5\times 10^{10}\text{ }years$$ (9) $`T_0=3K`$ In order to obtain a reasonable result we chose (with no explanation) the lower bounds for $`T_{NR}`$ and $`t_{NR\text{ }}`$ and for $`t_{cr}`$ we obtained : $$t_{cr}10^4t_010^{14}years$$ (10) concluding that the order of magnitude of $`t_{cr}`$ was a realistic one. In fact, $`10^4t_01.5\times 10^{14}years`$ after the big-bang all the stars will exhaust their fuel <sup>§</sup><sup>§</sup>§After the publication of paper , in 1996, new data about this time appear in ., so it is reasonable that this time would be of the same order than the one where the entropy gap stops its decreasing and begins to grow . But the choice of the lower bound in eqs. (8) was not explained and only the inhomogeneity was argued as above. ## II The photosphere as the main unstable system. Now we have reconsidered the problem and we conclude that, even if nuclear reactions within the stars are the main source of entropy, the parameters $`T_{NR}`$ and $`t_{NR}`$ are not the good ones to define the behavior of the term $`\mathrm{exp}(\gamma t)\rho _1`$ of equation (3)((3) of paper ), since they do not correspond to the main unstable system that we must consider In fact, to take the parameters $`T_{NR}`$ and $`t_{NR}`$ corresponds to take an homogenous gas model that fills the universe where, nuclear reactions take place in this homogenous gas, while in real universe these reactions take place within the star in a quite inhomogeneous scenario.. In fact the main production of entropy in a star is not located in its core, where the temperature is almost constant (and equal to $`T_{NR})`$, but in the photosphere where the star radiates. The energy radiated from the surface of the star is produced in the interior by fusion of light nuclei into heavier nuclei. Most stellar structures are essentially static, so the power radiated is supplied at the same rate by these exothermic nuclear reactions that take place near the center of the star . Once the star is formed, it settles into a termally stable state where all the nuclear energy is radiated at the surface and the rate of internal entropy change is extremely low . We can decompose the whole star in two branch systems, as explained in paper (or in section VII of paper ), where a chain of branch systems was introduced. We have two branch systems to study: the core and the photosphere. The core gives energy to the photosphere and in turn the photosphere diffuses this energy to the surroundings of the star, namely in the bath of microwave radiation at temperature $`T_{0\text{ .}}`$In this way, we have two sources of entropy production: the radiation of energy at the surface of the star and the change of composition inside the star (as time passes we have more helium and less hydrogen). Since the core of a star is near thermodynamic equilibrium, we neglect the second and we concentrate on the first: the radiation from the surface of the star (related with the difference between the star and the background temperatures). So the temperature of the photosphere and not the one of the core must be introduced in our formula. This is also the case for the lifetime. We must take the lifetime of the photosphere not the one of the nuclear reactions. Thus it is better to consider the photosphere as the unstable system that defines the term $`\mathrm{exp}(\gamma t)\rho _1`$ of equation (4). So we must change $`T_{NR}`$ and $`t_{NR}`$ by $`T_P`$, the temperature of the photosphere and $`t_S`$ the characteristic lifetime of the star respectively. Then we must change eq. (7) to: $$t_{cr}t_0\left(\frac{2}{3}\frac{T_P}{T_0}\frac{t_S}{t_0}\right)^3$$ (11) Considering that the mean mass of stars is $`0.64M_{}`$ with surface temperature $`4.610^3K`$ and lifetime $`3.810^{10}years`$ (see appendix) we obtain: $$t_{cr}10^{10}t_010^{20}years$$ (12) but now the computation was not done using an arbitrary choice of the lower bound in some data, but using the first meaningful figure In paper , using the photosphere data, but just orders of magnitude, we have again obtained (10). in all the data <sup>\</sup><sup>\</sup>\We have obtained a larger result than (10), but remember that the latter was obtained just by taking the lower bound in eqs. (8). If we would choose the mean values in these equations we would obtain a larger result than the one of eq. (12), without any significance.. ## III The stelliferous and the degenerated eras. Let us now compare this result with those of paper , where the future history of the universe is analyzed. $`t_{cr}10^{10}t_0=10^{20}years`$ is placed after the end of the ”stelliferous era” ($`10^6<t<10^{14}years)`$, where most of the energy generated in the universe arises from nuclear processes in the conventional star evolution, and at the beginning of the ”degenerated era” ($`10^{15}<t<10^{37}years)`$, where most of the (baryonic) mass of the universe is locked up in degenerated stellar objects: brown dwarfs, white dwarfs, and neutron stars. In this era energy is generated through proton decay and particle annihilation. So $`t_{cr}10^{20}years`$ is not a bad place for the minimum of the entropy gap, since it is bigger than the end of conventional star formation ($`10^{14}years)`$, it is also bigger than the typical time of star formation via brown dwarfs collision ($`10^{16}years)`$, and it is of the order of the time corresponding to stellar evaporation from galaxies ($`10^{19}years)`$. Namely a time where we can consider that the main mechanisms of formation of stars are ended while the beginning of the evaporation process has just started. This is the best place for $`t_{cr}`$, that must be the frontier between the growing order period (formation of structures with decaying entropy ) and the diminishing order period (decaying of the structures with growing entropy). Tacking into account the great uncertainty of all cosmological data we can say that $`t_{cr}`$ is located in the edge between the stelliferous normal era, where we are living and which is dominated by the formation of structures, and the future degenerated era, full of, by now strange, objects, where the growing of disorder will begin. So even if eq. (11) is the result of a very simple model it gives a very reasonable value. ## IV Conclusion. By choosing a more realistic model as the main source of entropy production, the photosphere of the stars, we have obtained a reasonable value for the time where the minimum of the entropy gap is reached. This is the frontier of the formation and destruction of structure periods and therefore one of the most important moments of the future history of the Universe. ## V Acknowledgments. We wish to thank Omar Benvenutto for fruitful discussions. This work was partially supported by grants Nos. CI1\-CT94-0004 of the European Community, PID-0150 and PEI-0126-97 of CONICET (National Research Council of Argentina), EX-198 of the Buenos Aires University and 12217/1 of Fundación Antorchas and the British Council. ## VI Appendix The stars have masses in the interval ,: $$0.1m=\frac{M_{}}{M_{}}100$$ (13) where $`M_{}`$ is the mass of the star and $`M_{}`$ is the mass of the Sun. The mean mass of stars can be calculated using the initial mass function (IMF) $`\xi (m)`$ defined by : $$dN=\xi (m)dm$$ (14) where $`dN`$ is the number of stars with masses between $`m`$ and $`m+dm`$ We will use as IMF : $$\frac{dN}{d(\mathrm{ln}m)}=\psi (\mathrm{ln}m)=\mathrm{exp}\left[A\frac{1}{2<\sigma >^2}\mathrm{ln}^2\left(\frac{m}{m_c}\right)\right]$$ (15) where $`<\sigma >1.57`$ , $`m_c0.1`$ and $`A`$ is a constant that sets the overall normalization of the distribution (which is not important because it cancels in the calculation of the mean mass of stars). This form of the IMF is consistent with observations (see ).As $$\xi (m)=\psi (\mathrm{ln}m)\frac{d(\mathrm{ln}m)}{dm}=\frac{1}{m}\psi (\mathrm{ln}m)=\frac{1}{m}\mathrm{exp}\left[A\frac{1}{2<\sigma >^2}\mathrm{ln}^2\left(\frac{m}{m_c}\right)\right]$$ (16) then the mean mass of stars is: $$<m>=\frac{\underset{0.1}{\overset{100}{}}m\xi (m)𝑑m}{\underset{0.1}{\overset{100}{}}\xi (m)𝑑m}=0.64$$ (17) Let us calculate the effective surface temperature of a star with mass $`m=0.64`$. The Luminosity of a star is given by the Stefan- Boltzmann law: $$L=4\pi R^2\sigma T_p^4$$ (18) where $`R`$ is the radius of the star, $`\sigma `$ is the Stefan- Boltzmann constant and $`T_p`$ is the effective temperature of the surface of the star. Then $$\frac{L}{L_{}}=\left(\frac{R}{R_{}}\right)^2\left(\frac{T_p}{T_{}}\right)^4$$ (19) where $`L_{}`$ , $`R_{}`$ ,$`T_{}`$ are, respectively, the luminosity, the radius and the surface temperature of the Sun. The radius and luminosity for stars in the main sequence, as a function of mass are (see ): $$\frac{R}{R_{}}\left(\frac{M_{}}{M_{}}\right)^\beta =m^\beta $$ (20) $$\frac{L}{L_{}}\left(\frac{M_{}}{M_{}}\right)^\eta =m^\eta $$ (21) with $`\beta 0.6`$ for stars of low mass and $`\eta 3.2`$ .So we have $$T_PT_{}\left(\frac{M_{}}{M_{}}\right)^{\frac{\eta 2\beta }{4}}T_{}m^{0.5}$$ (22) Using that $`T_{}5780K`$ and $`m=0.64`$ we obtain $$T_P4.610^3K$$ (23) The lifetime of a star can be calculated by the equation(see ): $$t_s=10^{10}\left(\frac{M_{}}{M_{}}\right)^\alpha years=10^{10}m^\alpha years$$ (24) where $`\alpha 34`$ for stars of low mass. Taking $`\alpha =3`$ we have: $$t_s3.810^{10}years$$ (25)
warning/0006/nucl-th0006021.html	ar5iv	text	# Anomalous Multiplicity Fluctuations from Phase Transitions in Heavy Ion Collisions ## Abstract Event-by-event fluctuations and correlations between particles produced in relativistic nuclear collisions are studied. The fluctuations in positive, negative, total and net charge are closely related through correlations. In the event of a phase transitions to a quark-gluon plasma, fluctuations in total and net charge can be enhanced and reduced respectively which, however, is very sensitive to the acceptance and centrality. If the colliding system experiences strong density fluctuations due, e.g., to droplet formation in a first-order phase transition, all fluctuations can be enhanced substantially. The importance of fluctuations and correlations is exemplified by event-by-event measurement of the multiplicities of $`J/\mathrm{\Psi }`$’s and charged particles since these observables should anti-correlate in the presence of co-mover or anomalous absorption. PACS numbers: 25.75+r, 24.85.+p, 25.70.Mn, 24.60.Ky, 24.10.-k Keywords: Relativistic heavy-ion collisions; Fluctuations; Quark-gluon plasma Event-by-event fluctuations have been measured at the SPS and, with the higher multiplicities of RHIC and LHC, will become an important tool for studying the anomalous fluctuations and correlations that might remain following the phase transition to a quark-gluon plasma (QGP). Studies of event-by-event fluctuations at SPS energies do not indicate the presence of new physics . It has, however, been proposed that large multiplicity fluctuations can arise from density fluctuations or droplets created by a first-order phase transition . Recently, it has also been suggested that fluctuations in net charge should be suppressed in a QGP . Here, we shall consider multiplicity fluctuations in some generality to see how they are affected by conservation of total charge and strangeness, to understand the correlations between various measured fluctuations, and to show how their measurement can reveal interesting details of the collision process. Multiplicity fluctuations between various kinds of particles can be strongly correlated. As an example, consider the multiplicities of positive and negative pions, $`N_+`$ and $`N_{}`$, in a rapidity interval $`\mathrm{\Delta }y`$ for any relativistic heavy-ion experiment. (Similar analyses can be performed for any two kinds of distinguishable particles.) We define the fluctuation in any multiplicity $`N`$ as $`\omega _N{\displaystyle \frac{N^2N^2}{N}}.`$ (1) Empirically, the fluctuations of Eq. (1) are typically of order unity in relativistic nuclear collisions which is consistent with the expectations of Poisson statistics. The net positive charge from the protons in the colliding nuclei is much smaller than the total charge produced in an ultrarelativistic heavy-ion collision. For example, $`N_+`$ exceeds $`N_{}`$ by only $`15`$% at in Pb$`+`$Pb collisions at SPS energies. The fluctuations in the number of positive and negative (or neutral) pions are also very similar, $`\omega _{N_+}\omega _N_{}`$. Charged particle fluctuations have been estimated in thermal as well as participant nucleon models including effects of resonances, acceptance, and impact parameter fluctuations. By varying the acceptance and centrality, the degree of thermalization can actually be determined empirically . Detailed analysis indicates that the fluctuations in central Pb+Pb collisions at the SPS are thermal whereas peripheral collisions are a superposition of pp fluctuations . The fluctuations in the total ($`N_{ch}=N_++N_{}`$) and net ($`Q=N_+N_{}`$) charge are defined as $`{\displaystyle \frac{(N_+\pm N_{})^2N_+\pm N_{}^2}{N_++N_{}}}=`$ (2) $`{\displaystyle \frac{N_+}{N_{ch}}}\omega _{N_+}+{\displaystyle \frac{N_{}}{N_{ch}}}\omega _N_{}\pm C,`$ (3) where the correlation is given by $`C={\displaystyle \frac{N_+N_{}N_+N_{}}{N_{ch}/2}}.`$ (4) Fluctuations in positive, negative, total and net charge can be combined to yield both the intrinsic fluctuations in the numbers of $`N_\pm `$ and the correlations in their production as well as a consistency check. These quantities can change as a consequence of thermalization and a possible phase transition. In practice, $`\omega _{N_+}\omega _N_{}`$, so that the fluctuation in total charge simplifies to $`\omega _{N_{ch}}`$ $``$ $`{\displaystyle \frac{N_{ch}^2N_{ch}^2}{N_{ch}}}=\omega _{N_+}+C,`$ (5) and that for the net charge becomes $`\omega _Q`$ $``$ $`{\displaystyle \frac{Q^2Q^2}{N_{ch}}}=\omega _{N_+}C.`$ (6) The fluctuation in net charge is related to the fluctuation in the ratio of positive to negative particles, $`\omega _QN_+/N_{}N_{ch}\omega _{N_{}/N_+}/4`$ plus volume (or impact parameter) fluctuations . The virtue of this expression is that volume fluctuations can in principle be extracted empirically. Alternatively one can vary the centrality bin size or the acceptance . In the following we shall assume that such “trivial” volume fluctuations have been removed. The analysis has so far been general and Eqs. (3-6) apply to any kind of distinguishable particles, e.g. positive and negative particles, pions, kaons, baryons, etc. - irrespective of what phase the system may be in, or whether it is thermal or not. In the following, we shall consider thermal equilibrium, which seems to apply to central collisions between relativistic nuclei, in order to reveal possible effects on fluctuations of the presence of a quark-gluon plasma. Bosons/fermions have thermal fluctuations, $`\omega _N=1\pm n_p^2/n_p`$ where $`n_p=(\mathrm{exp}(ϵ_p/T)1)^1`$ is the boson/fermion distribution function, which are slightly larger/smaller than those of Poisson statistics for a Boltzmann distribution. Massless bosons, e.g. gluons, have $`\omega _B=\xi (2)/\xi (3)=1.37`$ and massless fermions, e.g. quarks, have $`\omega _F=2\xi (2)/3\xi (3)0.91`$ independent of temperature. Massive bosons have smaller fluctuations with, for example, $`\omega _\pi =1.11`$ and $`\omega _\rho =1.01`$ when $`T=m_\pi `$. In a thermal hadron gas (HG) as created in relativistic in nuclear collisions, pions can be produced either directly or through the decay of heavier resonances, $`\rho ,\omega ,\mathrm{}`$. The resulting fluctuation in the measured number of pions is $`\omega _{N_+}=\omega _N_{}=f_\pi \omega _\pi +f_\rho \omega _\rho +f_\omega \omega _\omega +\mathrm{}.,`$ (7) where $`f_r`$ is the fraction of measured pions produced from the decay of resonance $`r`$, and $`_rf_r=1`$. These mechanisms are assumed to be independent, which is valid in a thermal system. The heavier resonances such as $`\rho ^0,\omega ,\mathrm{}`$ decay into pairs of $`\pi ^+\pi ^{}`$ and thus lead to a correlation $`C^{HG}={\displaystyle \frac{1}{3}}f_\rho +f_\omega +\mathrm{}..`$ (8) Resonances reduce the fluctuations in net charge in a HG to $`\omega _Q=0.70`$ . In addition, total charge neutrality reduces fluctuations in net charge when the acceptance is large and thus increases correlations as will be discussed below. A phase transition to the QGP can alter both fluctuations and correlations in the production of charged pions. To the extent that these effects are not eliminated by subsequent thermalization of the HG, they may remain as observable remnants of the QGP phase. As shown in Refs. , net charge fluctuations in a plasma of u, d quarks and gluons are reduced partly due to the intrinsically smaller quark charge and partly due to correlations from gluons $`\omega _Q={\displaystyle \frac{N_q}{N_{ch}}}\omega _F{\displaystyle \frac{1}{N_f}}{\displaystyle \underset{f=u,d,\mathrm{}}{\overset{N_f}{}}}q_f^2,`$ (9) where $`N_f`$ is the number of quark flavors, $`q_f`$ their charges, and $`N_q`$ the number of quarks. The total number of charged particles (but not the net charge) can be altered by the ultimate hadronization of the QGP. This effect can be estimated by equating the entropy of all pions to the entropy of the quarks and gluons. Since 2/3 of all pions are charged and since the entropy per fermion is 7/6 times the entropy per boson in a QGP $`N_{ch}{\displaystyle \frac{2}{3}}(N_g+{\displaystyle \frac{7}{6}}N_q),`$ (10) where the number of gluons is $`N_g=(16/9N_f)N_q`$. Inserting this result in (9), we see that the resulting fluctuations are $`\omega _Q=0.18`$ in a two-flavor QGP (and $`\omega _Q=0.12`$ for three flavors). As pointed out in , lattice results give $`\omega _Q0.25`$. However, according to a substantial fraction of the pions are decay products from the HG, and the entropy of the HG exceed that of a pion gas by a factor $`1.751.8`$. As described in the net charge fluctuations should be increased by this factor in the QGP, i.e. $`\omega _Q0.33`$ in a two-flavor QGP, whereas it remains similar in the HG, $`\omega _Q0.6`$. However, if the high density phase is initially dominated by gluons with quarks produced only at a later stage of the expansion by gluon fusion, the production of positively and negatively charged quarks will be strongly correlated on sufficiently small rapidity scales. If, for example, the entropy density increases by an order of magnitude in going from a HG to QGP without additional net charge production, fluctuations in net charge will be reduced significantly, $`\omega _Q^{QGP}{\displaystyle \frac{s_{HG}}{s_{QGP}}}\omega _Q^{HG}.`$ (11) The resulting fluctuation in net charge is necessarily smaller than that from thermal quark production as given by Eq. (9). A similar phenomenon occurs in string models where particle production by string breaking and $`q\overline{q}`$ pair production results in flavor and charge correlations on a small rapidity scale . The recently measured charged particle density at midrapidity in central nuclear collisions at RHIC is only $`3040`$% larger than $`pp`$ scaled up by the nuclear mass as was also found at SPS energies . If entropy and multiplicities are proportional, the net and total charged particle fluctuations should be the same as at SPS according to Eq. (11) unless anomalous non-thermal fluctuations occur as will be discussed below. The strangeness fluctuation in kaons $`K^\pm `$ might seem less interesting at first sight since strangeness is not suppressed in the QGP: The strangeness per kaon is unity, and the total number of kaons is equal to the number of strange quarks. However, if strange quarks are produced at a late stage in the expansion of a fluid initially dominated by gluons, the net strangeness will again be greatly reduced on sufficiently small rapidity scale. Consequently, fluctuations in net/total strangeness would be reduced/enhanced. The baryon number fluctuations have been estimated in a thermal model . It is, however, not known how the annihilation of baryons and antibaryons in the hadronic phase affect these results. If only charged particles are detected, but not $`K^0`$, $`\overline{K}^0`$, neutrons and antineutrons, the fluctuations have smaller correlations as compared to the total and net strangeness or baryon number. Total charge conservation is important when the acceptance $`\mathrm{\Delta }y`$ is a non-negligible fraction of the total rapidity. It reduces the fluctuations in the net charge as calculated within the canonical ensemble, Eqs. (7,8-10). If the total positive charge (which is exactly equal to the total negative charge plus the incoming nuclear charges) is independently distributed according to the single particle distributions, the resulting fluctuations within the acceptance are $`\omega _Q`$ $`=`$ $`1f_{acc},`$ (12) $`\omega _{N_{ch}}`$ $`=`$ $`1f_{acc}+2f_{acc}\omega _{N_+}`$ (13) where $`f_{acc}=N_{tot}^1_{\mathrm{\Delta }y}(dN_{ch}/dy)𝑑y`$ is the acceptance fraction or the probability that a charged particle falls into the acceptance $`\mathrm{\Delta }y`$. Since charged particle rapidity distributions are peaked near midrapidity, charge conservation effectively kills fluctuations in the net charge even when $`\mathrm{\Delta }y`$ is substantially smaller than the laboratory rapidity, $`y_{lab}6`$ (11) at SPS (RHIC) energies. Total charge conservation also has the effect of increasing $`\omega _{ch}`$ towards $`2\omega _{N_+}`$ according to Eqs. (5) and (13). Similar effects can be seen in photon fluctuations when photons are produced in pairs through $`\rho ^02\gamma `$. In the WA98 experiment, $`\omega _\gamma 2`$ is found after acceptance corrections and eliminating volume fluctuations . When the acceptance $`\mathrm{\Delta }y`$ is too small, particles from a thermal ensemble can diffuse in and out of the acceptance during hadronization and freezeout . Furthermore, correlations due to resonance production will disappear when the average separation in rapidity between decay products exceeds the acceptance. Each of these effects tend towards Poisson statistics when $`\mathrm{\Delta }y\stackrel{<}{}\delta y`$, where $`\delta y`$ denotes the rapidity interval that particles diffuse during hadronization, freezeout and decay. If $`\omega _Q`$ is the canonical thermal fluctuation of Eqs. (8,10), the resulting fluctuation after correcting for both $`\delta y`$ and total charge conservation is approximately $`\omega _Q^{exp}\left({\displaystyle \frac{\mathrm{\Delta }y}{\mathrm{\Delta }y+2\delta y}}\omega _Q+{\displaystyle \frac{2\delta y}{\mathrm{\Delta }y+2\delta y}}\right)(1f_{acc}).`$ (14) Here, the factor $`(1f_{acc})`$ is due to total charge conservation as in (12). The remainder is fluctuations from two sources: a thermal one with fluctuations $`\omega _Q`$ and a random one with Poisson fluctuations, each weighted with the fraction of the charged particles they contribute. The resulting fluctuations in total and net charge are shown in Fig. 1 assuming $`\omega _{N_+}=\omega _\pi 1.1`$ and $`\delta y=0.5`$. As mentioned above, $`f_{acc}`$ and $`\mathrm{\Delta }y`$ are related by the measured charge particle rapidity distributions . Preliminary NA49 data agree well with the net and total fluctuations in a HG ($`C=0.4`$) from Eqs. (14) and (10). Residual volume fluctuations are significant for $`\omega _{N_{ch}}`$ and have been estimated and subtracted. The curves apply to RHIC energies as well after scaling $`\delta y`$ with $`\mathrm{\Delta }y`$. The net charge fluctuations in a thermal HG corrected for finite acceptance and diffusion are slightly below the value without any intrinsic correlations given by Eq. (12). This reduction is due to the correlations in the HG that leads to $`\omega _Q<1`$. In high energy $`pp`$ collisions there are stronger rapidity correlations between unlike than like charged particles leading to similar magnitude for $`\omega _Q`$ <sup></sup><sup></sup>*The correlations are related to the two-body density distributions, e.g., $`N_+N_{}=\rho _+^{(2)}(y_1,y_2)𝑑y_1𝑑y_2`$, where the integral extends over $`y_1,y_2\mathrm{\Delta }y`$.. Therefore the net charge fluctuations does not vary by much going from peripheral $`pp`$-like high energy nuclear collisions to central collisions that are more likely to produce a thermal hadronic gas. The fluctuations in total charge are, however, very different because the total charge fluctuations in $`pp`$ collisions increase dramatically with collision energy, $`\omega _{N_{ch}}^{RHIC}6`$ and $`\omega _{N_{ch}}^{LHC}20`$ as compared to $`\omega _{N_{ch}}^{SPS}2.0`$ . Peripheral collisions will therefore be very different from central ones and the centrality dependence should be studied carefully to assess the degree of thermalization before anomalous fluctuations due to phase transitions can be determined . Large non-thermal fluctuations can arise as a consequence of density fluctuations due, e.g. to droplet formation in first-order phase transitions. These could lead to large fluctuations in multiplicities of charged particles and therefore also in the total and net charge. The above estimates for the fluctuations were of order unity. They implicitly assumed a uniform expanding system. Consider a scenario where the total multiplicity within the acceptance arises from a normal hadronic background component ($`N_{HG}`$) and from a second component ($`N_{QGP}`$) that has undergone a transition: $`N=N_{HG}+N_{QGP}.`$ (15) Its average is $`N=N_{HG}+N_{QGP}`$. Assuming that the multiplicity of each of these components is statistically independent, the multiplicity fluctuation becomes $`\omega _N=\omega _{HG}+(\omega _{QGP}\omega _{HG}){\displaystyle \frac{N_{QGP}}{N}}.`$ (16) Here, $`\omega _{HG}`$ is the standard fluctuation in hadronic matter $`\omega _{HG}1`$. The fluctuations due to the component that had experienced a phase transition, $`\omega _{QGP}`$, depend on the type and order of the transition, the speed with which the collision zone goes through the transition, the degree of equilibrium, the subsequent hadronization process, the number of rescatterings between hadronization and freezeout, etc. If thermal and chemical equilibration eliminate all signs of the transition, $`\omega _{QGP}\omega _{HG}`$. At the other extreme, the droplet scenario could produce $`\omega _{QGP}N10^210^3`$ if most hadrons arrive from a droplet so that either all or none fall into the acceptance . This is a promising signal worth looking for. Since droplets or density fluctuations are expected to be charge neutral, net charge fluctuations should vanish $`\omega _Q0`$ whereas $`\omega _{ch}2\omega _{N_+}2\omega _{QGP}`$. General correlators between all particle species should be measured event-by-event, e.g., the ratios $`{\displaystyle \frac{N_i/N_j}{N_i/N_j}}1+{\displaystyle \frac{\omega _{N_j}}{N_j}}{\displaystyle \frac{N_iN_jN_iN_j}{N_iN_j}},`$ (17) where $`N_{i,j}`$ are the multiplities in acceptances $`i`$ and $`j`$ of any particle type. (We assume that $`N_j`$ is so large that it never vanishes.) In the presence of droplets, $`N_i`$ and $`N_j`$ would be strongly correlated in nearby rapidity intervals and at all azimuthal angles. As another example, consider correlations between multiplicities of charmonium particles $`N_\psi `$, $`\psi =J/\mathrm{\Psi },\psi ^{},\chi ,..`$ and charged particle multiplicity ($`N`$) in a given rapidity interval $`\mathrm{\Delta }y`$. If a $`\psi `$ is absorbed on co-movers or anomalously suppressed by QGP, one would expect anti-correlations because the number of co-movers and QGP should scale with the multiplicity $`N`$. By contrast, direct Glauber absorption should not depend on the multiplicity of particles in $`\mathrm{\Delta }y`$ for a given centrality since it is the result of collisions with participating nucleons in Glauber trajectories along the beamline. To quantify this anti-correlation, we model the absorption of $`\psi `$s by simple Glauber theory $`{\displaystyle \frac{\overline{N}_\psi }{N_\psi ^0}}=e^{\sigma _{c\psi }\rho _cl}e^{\gamma N/N},`$ (18) where $`\overline{N}_\psi `$ is the average number of $`\psi `$’s for given charge particle multiplicity $`N`$, and $`N_\psi ^0`$ is the number of $`\psi `$s before co-mover or anomalous absorption sets in but after direct Glauber absorption on participant nucleons. In Glauber theory, the exponent is the product of the absorption cross section $`(\sigma _{c\psi }`$), the absorber density $`(\rho _c)`$, and the average path length $`(l)`$ traversed in matter. The density, therefore also the exponent, is proportional to the multiplicity $`N`$ with coefficient $`\gamma =d\mathrm{log}N_\psi /d\mathrm{log}N`$. In a simple co-mover absorption model for a system with longitudinal Bjorken scaling, $`\gamma `$ can be calculated to be approximately $`\gamma {\displaystyle \underset{c}{}}{\displaystyle \frac{dN_c}{dy}}{\displaystyle \frac{v_{c\psi }\sigma _{c\psi }}{4\pi R^2}}\mathrm{log}(R/\tau _0),`$ (19) where $`dN_c/dy`$ is the co-mover rapidity density, $`\sigma _{c\psi }`$ the absorption cross section, $`v_{c\psi }`$ the relative velocity, $`R`$ the transverse size of the overlap zone, and $`\tau _0`$ the formation time. Co-Mover absorption reduces the number of $`\psi `$ by a factor $`e^\gamma `$ where $`\gamma `$ increases with centrality up to $`\gamma 1`$ for typical parameters employed in co-mover absorption models . Since fluctuations in the exponent are small, $`\gamma \sqrt{\omega _N/N}1`$, the anti-correlation is $`{\displaystyle \frac{NN_\psi NN_\psi }{N_\psi }}=\gamma \omega _N.`$ (20) It is negative and proportional to the amount of co-mover and anomalous absorption. It vanishes when the absorption is independent of multiplicity ($`\gamma =0`$). The rapidity interval should not be less than the typical relative rapidities between co-movers and the $`\psi `$ or the rapidity range of a droplet. Even in central heavy-ion collisions $`\psi `$’s are rarely produced and so $`N_\psi =1`$ or 0. In both cases one should measure $`dN_{ch}/dy`$ and average over events with and without a $`\psi `$ separately. If comovers absorb the $`\psi `$ we expect that $`dN_{ch}/dy`$ is slightly smaller for the events with a $`\psi `$ than without, leading to the negative correlation in Eq. (20). A quantitative assessment of the average suppression of $`\psi `$s due to co-mover absorption versus direct $`\psi `$-nucleon absorption has been debated ever since the first measurements of $`J/\mathrm{\Psi }`$ suppression. The anticorrelations of Eq. (20) directly quantify the amount of co-mover or anomalous absorption and can therefore be exploited to distinguish between these and direct Glauber absorption mechanisms. In that respect it is similar to the elliptic flow parameter for $`\psi `$ . For a sample of $`N_{events}^\psi `$, the statistical uncertainty in $`\gamma `$ as determined by Eq. (20) is $`\sqrt{N/N_{events}^\psi }`$. If we take a rapidity bin $`\mathrm{\Delta }y1`$ and consider central heavy ion collisions, $`N`$ will range from $`10^210^3`$ in going from SPS to RHIC energies. A sample of $`10^410^5`$ $`\psi `$s would be sufficient to determine $`\gamma `$ with an accuracy of $`\pm 0.1`$. The analysis of the kaon to pion ratio by NA49 obtains such an accuracy by comparing to a mixed event analysis which removes systematic errors . The impact parameter fluctuations and correlations may not cancel exactly for the $`\psi /N_{ch}`$ ratio, as they do for the $`\pi ^+/\pi ^{}`$ ratio , because their production mechanisms differ. It is hard for the $`\psi `$ and soft for most of the charged particles and thus their multiplicities scale approximately with the number of binary NN collisions and the number of participants respectively. The number of NN collisions increase more rapidly with centrality and nuclear mass number $`(A^{4/3}`$. As a result the impact parameter correlations between the $`\psi `$ and $`N`$ in Eq. (20) will be slightly larger than the impact parameter fluctuations implicit in the second term in Eq. (17). The difference is a finite fraction of the total impact parameter fluctuations and will be of similar magnitude but opposite sign as correlations from co-mover absorption. The net impact parameter fluctuations will, however, depend on centrality and decrease as bin-size of the centrality cut decreases which should make it possible to separate it from other correlations. It would reveal independent information on the soft vs. hard production mechanisms. In summary, we have given a detailed analysis of total and net charge fluctuations and correlations including total charge conservation and diffusion effects and how they depend on the acceptance. The correlations may increase if a QGP is formed resulting in a reduction/enhancement of net/total charge by up to an order of magnitude depending on the model. It is important to measure fluctuations and correlations for various acceptances as well as versus centrality and/or beam energy. At SPS energies the fluctuations in net charge actually increase slightly with centrality due to thermalization which is opposite to the predicted decrease in net charge fluctuations if a QGP is formed in central heavy-ion collisions. At RHIC and LHC energies the correlations vs. centrality will be much larger because the total charge fluctuations in $`pp`$ collisions are $`\omega _{N_{ch}}^{RHIC}6`$ and $`\omega _{N_{ch}}^{LHC}20`$ as compared to $`\omega _{N_{ch}}^{SPS}2.0`$ . Peripheral collisions will therefore be very different from central ones and the centrality dependence should be studied carefully to assess the degree of thermalization before anomalous fluctuations due to phase transitions can be determined . It is important to measure all multiplicity fluctuations and correlations and understand the physical effects discussed above before a possible small increase in correlations can be attributed to the formation of a QGP. We are grateful for discussions with S. Voloshin (NA49) and T. Nayak (WA98) and for showing us preliminary data.
warning/0006/hep-lat0006013.html	ar5iv	text	# Introduction ## Introduction Supersymmetry is thought to be a crucial ingredient in any theory which attempts to unify the separate interactions contained in the standard model of particle physics. Since low energy physics is manifestly not supersymmetric it is necessary that this symmetry be broken at some energy scale. Issues of spontaneous symmetry breaking have proven difficult to address in perturbation theory and hence one is motivated to have some non-perturbative method for investigating such theories. The lattice furnishes such a framework. Unfortunately, supersymmetry being a spacetime symmetry is explicitly broken by the discretization procedure and it is highly non-trivial problem to show that it is recovered in the continuum limit. One manifestation of this problem is the usual doubling problem of lattice fermions - the naive fermion action in $`D`$ dimensions possesses not one but $`2^D`$ continuum-like modes. These extra modes persist in the continuum limit and yield an immediate conflict with supersymmetry requiring as it does an equality between boson and fermion degrees of freedom. As has been noted by several authors it is possible to circumvent this problem in a free theory by the addition of a simple Wilson mass term to the fermion action. This removes the doubles and leads to a supersymmetric free theory in the continuum limit. However such a procedure fails when interactions are introduced. Instead we shall show that the use of a non-standard lattice action allows the quantum continuum limit of such an interacting theory to admit continuum supersymmetry . Indeed, we will show that this action has an exact lattice supersymmetry in the absence of interactions (similar to that proposed in ) and very small symmetry breaking effects at non-zero interaction coupling. Finally, we write down an action for the interacting theory which is supersymmetric for all lattice spacings. ## Model The model we will study contains a real scalar field $`x`$ and two independent real fermionic fields $`\psi `$ and $`\overline{\psi }`$ defined on a one-dimensional lattice of $`L`$ sites with periodic boundary conditions imposed on both scalar and fermion fields. $$S=\underset{ij}{}\frac{1}{2}\left[x_iD_{ij}^2x_j+\overline{\psi }_i\left(D_{ij}+P_{ij}^{}\right)\psi _j\right]+\underset{i}{}\frac{1}{2}P_i^2$$ (1) The quantity $`P_i`$ is defined as $$P_i=\underset{j}{}K_{ij}x_j+gx_i^3$$ and its derivative is then just $$P_{ij}^{}=K_{ij}+3gx_i^2\delta _{ij}$$ The choice of a cubic interaction term in $`P(x)`$ guarantees unbroken supersymmetry in the continuum. The matrix $`D_{ij}`$ is the symmetric difference operator $$D_{ij}=\frac{1}{2}\left[\delta _{j,i+1}\delta _{j,i1}\right]$$ and $`K_{ij}`$ is the Wilson mass matrix $$K_{ij}=m\delta _{ij}\frac{r}{2}\left(\delta _{i,j+1}+\delta _{i,j1}2\delta _{ij}\right)$$ We work in dimensionless lattice units in which $`m=m_{\mathrm{phys}}a`$, $`g=g_{\mathrm{phys}}a^2`$ and $`x=a^{\frac{1}{2}}x_{\mathrm{phys}}`$. Notice that the boson operator $`D^2`$ is not the usual lattice Laplacian $`\mathrm{}=D_+D_{}`$ but contains a double corresponding to the extra zero in $`D=\frac{1}{2}(D_++D_{})`$. However, the boson action contains now a Wilson mass term and so this extra state, like its fermionic counterpart, decouples in the continuum limit. With the further choice $`r=1`$ the fermion matrix $`M=D+P^{}`$ is almost lower triangular and its determinant can be shown to be $$\mathrm{det}(M)=\underset{i=1}{\overset{L}{}}\left(1+m+3gx_i^2\right)1$$ which is positive definite for $`g>0`$ and $`m>0`$. This fact will be utilized in our numerical algorithm. Furthermore, the choice $`r=1`$ removes the doubles completely in the free theory and renders the fermion correlators simple exponentials. ## Simulation In order to simulate the fermionic sector we first replace the fermion field by a bosonic pseudofermion field $`\varphi `$ whose action is just $$\underset{ij}{}\frac{1}{2}\varphi _i\left(M^TM\right)_{ij}^1\varphi _j$$ This is an exact representation of the original fermion effective action provided the determinant of the fermion matrix is positive definite. The resultant (non-local) action $`S(x,\varphi )`$ can now be simulated using the Hybrid Monte Carlo (HMC) algorithm . In the HMC scheme momentum fields $`(p,\pi )`$ conjugate to $`(x,\varphi )`$ are added and a Hamiltonian $`H`$ defined which is just the sum of the original action plus additional terms depending on the momenta $`H=S+\mathrm{\Delta }S`$. $$\mathrm{\Delta }S=\underset{i}{}\frac{1}{2}\left(p_i^2+\pi _i^2\right)$$ On integrating out the momenta it is clear that this partition function is (up to a constant) identical to the original one. The augmented system $`(x,p,\varphi ,\pi )`$ is now naturally associated with some classical dynamics depending on an auxiliary time variable $`t`$ $`{\displaystyle \frac{x}{t}}`$ $`=`$ $`p`$ $`{\displaystyle \frac{p}{t}}`$ $`=`$ $`{\displaystyle \frac{H}{x}}`$ $`{\displaystyle \frac{\varphi }{t}}`$ $`=`$ $`\pi `$ $`{\displaystyle \frac{\pi }{t}}`$ $`=`$ $`{\displaystyle \frac{H}{\varphi }}`$ If we introduce a finite time step $`\mathrm{\Delta }t`$ we may simulate this classical evolution and produce a sequence of configurations $`(x(t),\varphi (t))`$. If $`\mathrm{\Delta }t=0`$ then $`H`$ would be conserved along such a trajectory. In practice $`\mathrm{\Delta }t`$ is finite and $`H`$ is not exactly conserved. However a finite length of such an approximate trajectory can still be used as a global move on the fields $`(x,\varphi )`$ which may then be subject to a Metropolis step based on $`\mathrm{\Delta }H`$. Provided the classical dynamics is reversible and care is taken to ensure ergodicity the resulting move satisfies detailed balance and hence this dynamics will provide a simulation of the original partition function. The reversibility criterion can be satisfied by using a leapfrog integration scheme and ergodicity is taken care of by drawing new momenta from a Gaussian distribution after each such trajectory. If we introduce bosonic and pseudofermionic forces $$F_i(t)=\frac{H(t)}{x_i}$$ (2) and $$_i(t)=\frac{H(t)}{\varphi _i},$$ (3) the resultant evolution equations look like $`x_i(t+\mathrm{\Delta }t)`$ $`=`$ $`x_i(t)+\mathrm{\Delta }tp_i(t){\displaystyle \frac{(\mathrm{\Delta }t)^2}{2}}F_i(t)`$ $`\varphi _i(t+\mathrm{\Delta }t)`$ $`=`$ $`\varphi _i(t)+\mathrm{\Delta }t\pi _i(t){\displaystyle \frac{(\mathrm{\Delta }t)^2}{2}}_i(t)`$ $`p_i(t+\mathrm{\Delta }t)`$ $`=`$ $`p_i(t)+{\displaystyle \frac{\mathrm{\Delta }t}{2}}\left(F_i(t)+F_i(t+\mathrm{\Delta }t)\right)`$ $`\pi _i(t+\mathrm{\Delta }t)`$ $`=`$ $`\pi _i(t)+{\displaystyle \frac{\mathrm{\Delta }t}{2}}\left(_i(t)+_i(t+\mathrm{\Delta }t)\right)`$ (4) The force terms are then given in terms of a vector $`s_i`$ which is a solution of the (sparse) linear problem $$\left(M^TM\right)_{ij}s_j=\varphi _i$$ $`{\displaystyle \frac{H}{x_i}}`$ $`=`$ $`D_{ij}^2x_j+P_{ij}^{}P_j+6gx_is_iM_{ij}s_j`$ $`{\displaystyle \frac{H}{\varphi _i}}`$ $`=`$ $`s_i`$ In order to reduce the effects of critical slowing down we have chosen to perform this update in momentum space using FFTs and a momentum dependent time step. Thus, for example, the lattice field $`x_m,(m=0\mathrm{}L1`$) can be expanded as $$x_m=\sqrt{\frac{1}{L}}\underset{n}{}\overline{x}_ne^{i\frac{2\pi mn}{L}}$$ where $`\overline{x}_n`$ is the Fourier amplitude with wavenumber $`\frac{2\pi }{L}n`$ ($`n=0\mathrm{}L1`$) and the Fourier amplitudes $`\overline{x}_n`$ are updated using equations 4 with $`\mathrm{\Delta }t=\mathrm{\Delta }t(n)`$. For the boson field $`x`$ we use $`\mathrm{\Delta }t=\tau _B`$ where $$\tau _B(n)=ϵ\frac{\left(m_{\mathrm{eff}}+2r\right)}{\sqrt{\mathrm{sin}^2\frac{2\pi n}{L}+\left(m_{\mathrm{eff}}+2r\mathrm{sin}^2\frac{2\pi n}{2L}\right)^2}}$$ For the pseudofermion field update we use the inverse function $`\mathrm{\Delta }t=\tau _F=\frac{1}{\tau _B}`$. With these choices (and $`m_{\mathrm{eff}}=ma`$) it is simple to show that the $`g=0`$ theory suffers no critical slowing down \- all modes are updated at the same rate independent of their wavelength. By setting $`m_{\mathrm{eff}}`$ at the approximate position of the massgap in the interacting case we have found very substantial reductions in the autocorrelation time for the two point functions of the theory. In practice we set $`ϵ0.1`$ and the number of leapfrog integrations per trajectory at $`N_{\mathrm{leap}}=10`$. ## Correlation functions We have measured the following correlators $$G_{ij}^B=x_ix_j$$ and $$G_{ij}^F=s_jM_{ik}s_k$$ It can be shown that the latter is simply an estimator for the original fermion correlator $`<\overline{\psi }_i\psi _j>`$. In figure 1 we show first the results of a simulation of this model for $`g_{\mathrm{phys}}=0.0`$ and $`m_{\mathrm{phys}}=10.0`$. The data set consists of $`10^6`$ Fourier accelerated HMC trajectories. The plot shows both boson and fermion massgaps, extracted from a simple exponential fit to the correlators over the first $`L/4`$ timeslices, as a function of the lattice spacing $`a=1/L`$. Notice that boson and fermion masses while receiving large O(a) systematic errors (due to the Wilson term) are degenerate within statistical errors. We see furthermore that as $`\mathrm{a}0`$ the common massgap approaches the correct continuum value. As we shall see later the free action has an exact supersymmetry at finite lattice spacing which is responsible for the boson/fermion degeneracy. We have also examined the massgaps at non-zero coupling. Figure 2 shows the same plot for $`m_{\mathrm{phys}}=10.0`$ and $`g_{\mathrm{phys}}=100.0`$. The massgaps are also listed in Table 1. The data set consists of $`10^6`$ trajectories again using lattice sizes $`L=16256`$. The effective dimensionless expansion parameter is $`g/m^2`$ so this corresponds to a regime of strong coupling. Remarkably, the boson and fermion masses are again degenerate within statistical errors O($`0.5`$%) and flow as $`\mathrm{a}0`$ to the correct continuum limit (the latter can be computed easily using Hamiltonian methods and yields $`m_{\mathrm{cont}}=16.87`$)). That this result is nontrivial can be seen when we compare it to the result of a ‘naive’ discretization of the continuum action using $`\mathrm{}`$ in place of $`D^2`$ and the usual Wilson action for the fermions – figure 3. In this case the mass plot looks very different. At large lattice spacing the extracted massgaps differ widely – the fermion having O(a) errors while the boson is much smaller (it varies as O($`\mathrm{a}^2`$) at $`g=0`$). Initially they appear to approach each other as $`\mathrm{a}0`$ but the two curves depart for fine lattice spacing and do not approach the correct continuum limit - the quantum continuum limit is not supersymmetric. Thus naive discretizations of the continuum action will break supersymmetry irreversibly even in theories such as quantum mechanics which have no divergences. At minimum it would be necessary to tune parameters to obtain a supersymmetric continuum limit. In comparison the numerical results of figure 2 indicate that supersymmetry breaking effects, if present, are very small. We examine this more carefully next. ## Supersymmetry Motivated by the form of the continuum supersymmetry transformations for this model consider the following two lattice transformations $`\delta _1x_i`$ $`=`$ $`\overline{\psi _i}\xi `$ $`\delta _1\psi _i`$ $`=`$ $`\left(D_{ij}x_jP_i\right)\xi `$ $`\delta _1\overline{\psi _i}`$ $`=`$ $`0`$ (5) and $`\delta _2x_i`$ $`=`$ $`\psi _i\overline{\xi }`$ $`\delta _2\psi _i`$ $`=`$ $`0`$ $`\delta _2\overline{\psi _i}`$ $`=`$ $`\left(D_{ij}x_j+P_i\right)\overline{\xi }`$ (6) where $`\xi `$ and $`\overline{\xi }`$ are independent anti-commuting parameters. The existence of two such symmetries reflects the $`N=2`$ character of the continuum supersymmetry. If we perform the variation corresponding to the first of these (eqn. 5) we find $$\delta _1S=\underset{i}{}\overline{\psi }_i\xi \left(P_{ij}^{}D_{jk}x_kD_{ij}P_j\right)$$ (7) The expression corresponding to the second transformation eqn. 6 is similar $$\delta _2S=\underset{i}{}\overline{\xi }\psi _i\left(P_{ij}^{}D_{jk}x_kD_{ij}P_j\right)$$ (8) In the continuum limit $`a=0`$ the difference operators become derivatives and the term inside the brackets is zero - this is the statement of continuum supersymmetry. Notice, for $`g=0`$ and $`a0`$ this term is still zero - the classical free lattice action is also supersymmetric. However this term is non-zero for finite spacing and non-zero interaction coupling - the classical lattice action breaks supersymmetry. Since we use the symmetric difference operator this breaking will be $`\mathrm{O}(gL^2)`$. From the point of view of a continuum limit such a breaking would not be important - since the theory contains no divergences all non-supersymmetric terms induced in the quantum effective action will have couplings that vanish as $`a0`$. Indeed, as was shown explicitly in , the two-dimensional Wess-Zumino model has a supersymmetric continuum limit when regulated in this way. For a lattice of size $`L=16`$ and $`g_{\mathrm{phys}}=100.0`$ we would expect symmetry breaking terms to be suppressed by a factor of $`g_{\mathrm{phys}}/L^4=0.002`$. This is consistent with what we see in the massgaps. To verify these conclusions we have studied the approximate Ward identities which follow from the lattice transformations eqn. 5 and eqn. 6. ## Ward identities The Ward identities corresponding to these approximate symmetries can be derived in the usual way. First consider the partition function with external sources $$Z(J,\theta ,\overline{\theta })=DxD\psi D\overline{\psi }e^{S+{\scriptscriptstyle J}.x+\overline{\theta }.\psi +\theta .\overline{\psi }}$$ Perform a lattice supersymmetry transformation, for example, eqn. 5. The action $`S`$ varies as in eqn. 7, and the integration measure is invariant while the source terms vary. Since the partition function does not change (the transformation can be viewed as a change of variables) we find $$0=\delta Z=DxD\psi D\overline{\psi }e^S(J.\delta x+\overline{\theta }.\delta \psi +\theta .\delta \overline{\psi }+\alpha .\overline{\psi })$$ (9) where $`\alpha _i=\left(P_{ij}^{}D_{jk}x_kD_{ij}P_j\right)`$ Furthermore, any number of derivatives with respect to the sources evaluated for zero sources will also vanish. For the first supersymmetry equation 5 this yields a set of identities connecting different correlation functions. The first non-trivial example is $$\overline{\psi _i}\psi _j+(D_{jk}x_kP_j)x_i+\alpha _k\overline{\psi }_kx_i\psi _j=0$$ The last term represents the symmetry breaking term which may be rewritten as $`<(M^T)_{jk}^1\alpha _kx_i>`$. Since $`\alpha _k`$ is a vector with random elements each of mean zero and suffering fluctuations O($`g/L^2)`$) we might expect that this term contributes rather a small correction to the Ward identity. In this spirit we will neglect it at this point and see if the predictions are substantiated by the results of the simulation. It is important to notice that this correction to the naive Ward identity is finite (quantum mechanics) and multiplied by $`1/L^2`$ and consequently the continuum limit is guaranteed to possess supersymmetry. The same will be true for any approximate Ward identity we care to construct. A second Ward identity may be derived corresponding to the second (approximate) supersymmetry equation 6 we find $$\psi _i\overline{\psi _j}+(D_{jk}x_k+P_j)(x)x_i=0$$ More conveniently we can add and subtract these equations to yield the relations $$G_{ij}^FG_{ji}^F=2D_{jk}x_kx_i$$ (10) and $$G_{ij}^F+G_{ji}^F=2P_jx_i$$ (11) Translation invariance on the lattice implies $`G_{ij}^F=G^F(t)`$ and $`G_{ji}^F=G^F(Lt)`$ where $`t=(ji)`$. We can check that these two Ward identities are satisfied numerically by forming the two (distance dependent) quantities $`w_1`$ and $`w_2`$ which are defined by $$w_1\left(t\right)=\frac{G_{ij}^F+G_{ji}^F2x_iP_j}{2x_iP_j}$$ (12) $$w_2\left(t\right)=\frac{G_{ij}^FG_{ji}^F+2x_iD_{jk}x_k}{2x_iD_{jk}x_k}$$ (13) These are shown in figure 4 and figure 5 for a lattice of size $`L=16`$ at $`g_{\mathrm{phys}}=0.0`$. In this case we expect the symmetry to be exact and indeed we see that the Ward identities are satisfied within statistical accuracy. Figures 6 and 7 show plots of $`w_1`$ and $`w_2`$ for a lattice of size $`L=16`$ at $`g_{\mathrm{phys}}=100.0`$. It is clear that within our statistical error (on the order of a few percent for these quantities) we are again not sensitive to the SUSY breaking terms and the continuum Ward identities are satisfied. ## Discussion and Conclusions We have performed a numerical study of the lattice supersymmetric anharmonic oscillator computed using path integrals. This is essentially a one-dimensional version of the Wess-Zumino model. We have utilized a lattice discretization which preserves two exact supersymmetries in the free theory. We are able to show that the interacting theory flows to a supersymmetric fixed point in the zero lattice spacing limit without fine tuning. This is to be contrasted with naive discretizations of the continuum action which require fine tuning to recover supersymmetry in the continuum limit. Furthermore, we have estimated the magnitude of supersymmetry breaking at O($`g/L^2`$) which is typically smaller than one percent even at strong coupling and for coarse lattices. Thus the lattice simulations are, in practice, very close to the supersymmetric fixed point. We have checked the first two non-trivial Ward identities following from this (approximate) invariance. Our numerical results place an upper bound on the magnitude of symmetry breaking corrections which is consistent with this estimate. It is tempting to try to interpret the numerical results as evidence of an exact lattice supersymmetry even in the presence of interactions. Using the antisymmetry of the derivative operator it is easy to show that the symmetry breaking term 8 can be rewritten $$\delta _2S=\delta _2x_iD_{ij}P_j$$ Thus $`S_{\mathrm{new}}=Sx_iD_{ij}P_j`$ will be exactly invariant under the second lattice supersymmetry transformation even for non zero interaction. Notice that the presence of an extra minus sign prevents this new action from having a second invariance corresponding to the first supersymmetry 7. This invariant action can be rewritten in the form $$S_{\mathrm{new}}=\underset{ij}{}\frac{1}{2}\left(D_{ij}x_j+P_i\right)\left(D_{ij}x_j+P_i\right)+\underset{ij}{}\overline{\psi }_i\left(D_{ij}+P_{ij}^{}\right)\psi _j$$ This allows us to identify the Nicolai map for the model . The latter is the non-trivial transformation $`x_i\xi _i`$ which maps the boson action to a free field form and whose Jacobian simultaneously cancels the fermion determinant. Here we see it explicitly $$\xi _i=D_{ij}x_j+P_i$$ It has previously been pointed out that the identification of such a map may be used to help find lattice supersymmetric actions and . This quantum mechanics model furnishes a concrete example - the lattice action which admits the Nicolai map is invariant under a transformation which interchanges bosonic and fermionic degrees of freedom. In the continuum limit this lattice action approaches its continuum counterpart and the transformation reduces to a continuum supersymmetry transformation. The presence of one exact supersymmetry is already enough to guarantee vanishing vacuum energy and boson/fermion mass degeneracy for the lattice theory. It would be interesting to extend these calculations to the two-dimensional Wess-Zumino model and verify non-perturbatively the results derived perturbatively in . ## Acknowledgements Simon Catterall was supported in part by DOE grant DE-FG02-85ER40237. We would like to thank Yoshio Kikukawa for drawing our attention to refs and .
warning/0006/astro-ph0006082.html	ar5iv	text	# A Model of Metallicity Evolution in the Early Universe ## 1 Introduction In a recent paper Prochaska & Wolfe (2000) reported new data on \[Fe/H\] for damped Ly$`\alpha `$ systems at high redshifts. They also summarized available data with redshifts from $`z1.5`$ to 4.5 (e.g., Lu et al. 1996; Lu, Sargent, & Barlow 1997; Prochaska & Wolfe 1999). Their observations are in general accord with previous studies and expand the database at $`z>3`$ substantially. Prochaska & Wolfe (2000) emphasized that there is a wide spread in \[Fe/H\] at a given $`z`$ and no damped Ly$`\alpha `$ system has \[Fe/H\] $`<2.7`$ (see also Lu et al. 1996, 1997). They also pointed out that there is only a minimal growth of \[Fe/H\] from $`z4.5`$ to 1.5. In a theoretical study of Fe and $`r`$-process abundances in very metal-poor stars in the Galaxy by Wasserburg & Qian (2000, hereafter WQ), it was proposed that the first stars formed after the Big Bang were very massive ($`100M_{}`$) and promptly enriched the interstellar medium (ISM) to \[Fe/H\] $`3`$, at which metallicity formation of normal stars (with masses $`1`$$`60M_{}`$) took over. Subsequent Fe enrichment was provided by a subset of Type II supernovae. The interpretations of WQ were based on observations of metal-poor stars in the Galaxy by Gratton & Sneden (1994), McWilliam et al. (1995), McWilliam (1998), and Sneden et al. (1996, 1998). The apparent agreement between the lower bound on \[Fe/H\] of damped Ly$`\alpha `$ systems and the critical metallicity \[Fe/H\] $`3`$ deduced by WQ for transition from formation of very massive stars to normal stars suggests that this Fe enrichment model deserves further study. Here we present a phenomenological model for the range of \[Fe/H\] at a given $`z`$ for damped Ly$`\alpha `$ systems based on the model of WQ. It is assumed that formation of normal stars started in a damped Ly$`\alpha `$ system (i.e., a protogalactic system turned on) at a time $`t^{}`$ after the Big Bang. The assembly of \[Fe/H\] at a given $`z`$ is interpreted as a sampling of $`t^{}`$ ranging from 0 to the age of the universe at $`z`$, $`t(z)`$. The value $`t^{}0`$ corresponds to the upper bound on \[Fe/H\]. Furthermore, the distribution of $`t^{}`$ at a given $`z`$ indicates the rate of turn-on of protogalaxies prior to $`t(z)`$. The data near $`z=2.2`$ suggest that this rate was initially very low and slowly reached a maximum at $`3`$ Gyr after the Big Bang. We describe the Fe enrichment model in detail in §2 and apply it to explain the data on \[Fe/H\] of damped Ly$`\alpha `$ systems in §3. Discussion and conclusions are given in §4. ## 2 Fe Enrichment and Abundances in Metal-Poor Stars The Fe enrichment model of WQ was developed to explain the relation between abundances of Fe and $`r`$-process elements ($`r`$-elements) in metal-poor stars in the Galaxy. Meteoritic data on the inventory of radioactive <sup>129</sup>I and <sup>182</sup>Hf in the early solar system require at least two distinct Type II supernova sources for the $`r`$-process (Wasserburg, Busso, & Gallino 1996; Qian, Vogel, & Wasserburg 1998; Qian & Wasserburg 2000). These are the high-frequency H events responsible for heavy $`r`$-elements with mass numbers $`A>130`$ (e.g., Ba and Eu) and the low-frequency L events responsible for light $`r`$-elements with $`A130`$ (e.g., Ag). The recurrence timescales for the H and L events are $`\mathrm{\Delta }_\mathrm{H}10^7`$ yr and $`\mathrm{\Delta }_\mathrm{L}10^8`$ yr, as required by replenishment of the appropriate radioactive nuclei in a standard mixing mass ($``$ the size of a molecular cloud) for a supernova. Additional evidence in support of different sources for the heavy and light $`r`$-elements has been found by Sneden et al. (2000). The observed wide dispersion in abundances of the heavy $`r`$-elements such as Ba and Eu over a narrow range of \[Fe/H\] $`3`$ to $`2.8`$ (McWilliam et al. 1995; McWilliam 1998; Sneden et al. 1996, 1998) led WQ to conclude that the H events cannot produce a significant amount of Fe. In contrast, there is a correlation between abundances of Fe and the heavy $`r`$-elements at \[Fe/H\] $`2.5`$ (Gratton & Sneden 1994; see also McWilliam et al. 1995). As Type Ia supernovae would occur only at metallicities much higher than \[Fe/H\] $`=2.5`$, Fe enrichment of very metal-poor stars with \[Fe/H\] $`2.5`$ must be provided by the L events. Consequently, the L events are responsible for the part of Fe contributed by Type II supernovae in general. There are $`100`$ L events during the time of $`10^{10}\mathrm{yr}`$ prior to solar system formation. To provide $`1/3`$ of the solar Fe inventory by these events, each L event must enrich a standard mixing mass with \[Fe/H\]$`{}_{\mathrm{L}}{}^{}2.5`$. The near absence of the heavy $`r`$-elements in stars with \[Fe/H\] $`4`$ to $`3`$ (McWilliam et al. 1995; McWilliam 1998) and the sharp increase in the abundances of such elements at \[Fe/H\] $`3`$ to $`2.8`$ led WQ to conclude that a source other than Type II supernovae must exist to produce Fe (along with other elements such as C, N, O, Mg, and Si) at \[Fe/H\] $`3`$. This source was attributed by WQ to very massive stars (with masses $`100M_{}`$) that first formed after the Big Bang. They further argued that formation of normal stars could not occur until \[Fe/H\] $`3`$ was reached. Presumably, this critical metallicity corresponds to conditions in the ISM that permit sufficient cooling and fragmentation to occur in collapsing gas clouds. A recent study by Bromm, Coppi, & Larson (1999) suggests that the very first stars were rather massive. The products of very massive stars formed from Big Bang debris have been discussed earlier by Ezer & Cameron (1971). However, nucleosynthesis in such stars remains to be tested with adequate stellar models. ## 3 \[Fe/H\] of Damped Ly$`\alpha `$ Systems To discuss \[Fe/H\] of a damped Ly$`\alpha `$ system at a given $`z`$, we assume the following history for its evolution: (1) at time $`t_1`$ after the Big Bang, matter consisting of Big Bang debris was isolated to form a system; (2) prompt enrichment by the first very massive stars ended at time $`t_2`$, resulting in \[Fe/H\] $`3`$ in the ISM; and (3) formation of normal stars began at time $`t^{}`$, with the L events providing further Fe enrichment to an average ISM at regular intervals. The above times are related as $`t_1<t_2<t^{}`$. As $`t^{}t(z)`$, it is necessary to define the value of $`t(z)`$ that is used. We take the redshift $`z`$ to correspond to a time $$t(z)\frac{2}{3}H_0^1\mathrm{\Omega }_\mathrm{m}^{1/2}(1+z)^{3/2}$$ (1) after the Big Bang. We take the Hubble constant $`H_0=65\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and the matter contribution to the critical density $`\mathrm{\Omega }_\mathrm{m}=0.3`$ \[equation (1) gives essentially the exact result for $`t(z)`$ at $`z1.5`$ for a flat universe with $`\mathrm{\Omega }_\mathrm{m}=0.3`$ and a cosmological constant\]. By our assumption, for $`t(z)t_2`$ the metallicity of a damped Ly$`\alpha `$ system is at least \[Fe/H\] $`3`$. Prochaska & Wolfe (2000) reported that no system at $`z1.5`$ to 4.5 has \[Fe/H\] $`<2.7`$. This means that the time required to provide prompt Fe enrichment must be less than $`t(z4.5)1.4`$ Gyr. The lower bound on \[Fe/H\] due to prompt enrichment is shown as a band between $`3`$ and $`2.8`$ in Figure 1 together with the data summarized in Prochaska & Wolfe (2000). In general, for $`t(z)t^{}`$ the metallicity of a specific damped Ly$`\alpha `$ system at $`z`$ is given by $$(\mathrm{Fe}/\mathrm{H})=(\mathrm{Fe}/\mathrm{H})_\mathrm{p}+(\mathrm{Fe}/\mathrm{H})_\mathrm{L}\frac{t(z)t^{}}{\mathrm{\Delta }_\mathrm{L}},$$ (2) where the number ratio in round brackets (Fe/H) is related to the standard square bracket notation by \[Fe/H\] $`=\mathrm{log}(\mathrm{Fe}/\mathrm{H})\mathrm{log}(\mathrm{Fe}/\mathrm{H})_{}`$, with (Fe/H)<sub>p</sub> corresponding to the prompt enrichment and (Fe/H)<sub>L</sub> to the contribution from a single L event to a standard reference mass of hydrogen. We take \[Fe/H\]$`{}_{\mathrm{p}}{}^{}=2.8`$, \[Fe/H\]$`{}_{\mathrm{L}}{}^{}=2.5`$, and $`\mathrm{\Delta }_\mathrm{L}=10^8`$ yr, the same parameters used by WQ. As recognized by Lu et al. (1996), the scatter in \[Fe/H\] at a given $`z`$ for damped Ly$`\alpha `$ systems might result from their different formation histories. Equation (2) explicitly states that the range of \[Fe/H\] at $`z`$ is caused by the different start times $`t^{}`$ for normal star formation in individual systems (see Figure 2a). The values for \[Fe/H\] corresponding to $`t^{}=1`$ and 1.5 Gyr are shown in Figure 1. Damped Ly$`\alpha `$ systems that just turned on at $`t^{}t(z)`$ would have \[Fe/H\] $``$ \[Fe/H\]<sub>p</sub>. An upper bound on \[Fe/H\] exists as damped Ly$`\alpha `$ systems that turned on at $`t^{}0`$ would have the longest history of normal star formation, and hence, the highest \[Fe/H\]. This bound is insensitive to the choice of \[Fe/H\]<sub>p</sub>. It can be seen from Figure 1 that almost all the data lie below the upper bound (usually well below this bound). Even the three exceptions are close to this bound (see §4). At a fixed $`z`$, equation (2) with equation (1) can be used to determine the start time $`t^{}`$ for normal star formation in a damped Ly$`\alpha `$ system from its \[Fe/H\]. In turn, a histogram of the number of systems at $`z`$ within a given \[Fe/H\] interval determines the probability $`p(t^{},t(z))dt^{}`$ for normal star formation to start in the interval between $`t^{}`$ and $`t^{}+dt^{}`$ after the Big Bang. Knowing the probability distribution $`p(t^{},t(z))`$ over $`0<t^{}t(z)`$ at $`z`$, we expect that the probability distribution at $`z^{}>z`$ can be obtained by discarding the part of $`p(t^{},t(z))`$ for $`t^{}>t(z^{})`$ and renormalizing the remaining part over $`0<t^{}t(z^{})`$ (see Figure 2b). For the case of a simple power-law distribution $`p(t^{},t(z))=[(\alpha +1)/t(z)][t^{}/t(z)]^\alpha `$. In this case, the average start time for normal star formation in damped Ly$`\alpha `$ systems at $`z`$ is $`t^{}(z)=[(\alpha +1)/(\alpha +2)]t(z)`$. The data in Figure 1 show a relatively high concentration in the interval $`2.0z2.4`$. The average start time in this interval is $`t^{}(z=2.2)2.5`$ Gyr. The frequency of occurrences of $`t^{}`$ calculated from the data in this interval is shown as a histogram in Figure 3. It can be seen that the frequency of occurrences is low for $`t^{}0`$ (close to the Big Bang) and slowly increases to a maximum at $`t^{}3`$ Gyr. Assuming a power-law distribution for $`t^{}`$, we obtain $`\alpha 3`$, for which $`t^{}(z)0.8t(z)`$ (see the corresponding curve for \[Fe/H\] shown in Figure 1). Values of $`t^{}`$ for all the data are shown in Figure 4. The clustering of $`t^{}`$ close to $`t(z)`$ indicates that typically there is a long delay between the Big Bang and the start time $`t^{}`$ for normal star formation in damped Ly$`\alpha `$ systems. ## 4 Discussion and Conclusions We consider the dominant cause of dispersion of \[Fe/H\] at a given $`z`$ for damped Ly$`\alpha `$ systems to be the variation in the start time $`(t^{})`$ after the Big Bang for normal star formation in different protogalaxies. The bounds on \[Fe/H\] from our model appear to closely define the observed ranges. The average $`t^{}`$ for damped Ly$`\alpha `$ systems at $`z`$ is $`t^{}(z)0.8t(z)`$. The rate of turn-on of protogalaxies was initially very low and slowly increased to a maximum at $`3`$ Gyr after the Big Bang. It is not possible to identify a turnover in this rate without more data at $`z<2`$. We suggest that the approach outlined here is a method for dating the start time for normal star formation in protogalaxies. As the formation of quasars is closely related to star formation, we consider that the histogram shown in Figure 3 may offer some insights into the rate of formation of quasars. For example, if the rate of quasar formation $`R_\mathrm{Q}`$ is proportional to the turn-on rate of protogalaxies, Figure 3 suggests that $`R_\mathrm{Q}(t^{})^3(1+z)^{4.5}`$ at $`z2`$. This gives a decrease by a factor of 2.7 in $`R_\mathrm{Q}`$ from $`z=3`$ to 4, consistent with the drop in quasar comoving space density at these redshifts (Schmidt, Schneider, & Gunn 1995). The three data points that lie above the upper bound in Figure 1 could be explained if the enrichment rate $`\beta _\mathrm{L}(\mathrm{Fe}/\mathrm{H})_\mathrm{L}/\mathrm{\Delta }_\mathrm{L}`$ of L events were increased by a factor $`2`$. It is possible that $`\beta _\mathrm{L}`$ has a spread of a factor of 2 or 3. Alternatively, we can consider the enrichment rate as a steep function of time, with $`\beta _\mathrm{L}(t)`$ starting quite high and then settling down to the value proposed by WQ. In this case (Fe/H) $`=`$ (Fe/H)$`{}_{\mathrm{p}}{}^{}+_t^{}^{t(z)}\beta _\mathrm{L}(t^{})dt^{}`$ and the straight line evolution in Figure 2a would be replaced by a curve. We note that Fe enrichment by Type Ia supernovae would not be significant for most damped Ly$`\alpha `$ systems as such enrichment appears to be significant only at \[Fe/H\] $`>1`$ in the Galaxy (Timmes, Woosley, & Weaver 1995). A more difficult matter is the timescale for condensing protogalactic globs during the expansion of the universe and the timing sequence outlined above. The range in $`t^{}`$ required is large and indicates that the time required to form protogalaxies and condense most of the baryonic matter into stars is comparable to the age of the universe at $`z2`$ ($`3.5`$ Gyr after the Big Bang) or possibly even longer. This is in conflict with ab initio models that report almost complete condensation of dark matter and possibly cloud or protogalaxy formation at $`z10`$–20 (e.g., Kamionkowski, Spergel, & Sugiyama 1994). The occurrence of \[Fe/H\] $`2.6`$ at $`z2`$ implies that there are regions where formation of normal stars did not begin until $`3.5`$ Gyr after the Big Bang. The occurrence of \[Fe/H\] $`2.6`$ to $`2.4`$ in the range of $`z2.0`$ to 4.2 indicates that a large fraction of the baryonic matter has not been condensed into protogalaxies and stars over the corresponding extended time range. If damped Ly$`\alpha `$ systems are random samples of the original medium, this indicates that the reservoir of original uncondensed baryonic material may not have been seriously diminished over $`3.5`$ Gyr. Indeed, it is possible that the bulk of this baryonic matter is dispersed and has not condensed today. Measurements of damped Ly$`\alpha `$ systems at lower redshifts (but at times before Type Ia supernovae contribute Fe) would provide a test. We have no means of establishing a priori whether the intrinsic rate of protogalaxy formation decreases with time (possibly due to decrease in global density) without appreciably depleting the reservoir of baryonic matter. A variant of the above scenario is possible which would not be in conflict with the models that suggest almost complete condensation of baryonic matter by $`z10`$. The modified scenario would be that most baryonic matter in the potential well created by non-baryonic dark matter is rapidly collected into Ur-protogalaxies, that formation of very massive stars from Big Bang debris is rapid \[$`t_2t(z)`$\], and that the explosion of these very massive stars destroys the Ur-protogalaxies and redistributes matter into the general medium providing a uniform source of Fe (and other elements such as C, N, O, Mg, and Si). The material in this second generation medium is then the source for much slower aggregation and formation of protogalaxies and stars. The explosion of very massive stars is likely to be very energetic. The potential wells formed by dark matter are considered to have a typical escape velocity of $`400`$ km s<sup>-1</sup>. The velocity of debris ejected from explosion of very massive stars is almost certainly much larger than this value. The inter-protogalactic medium would then be supplied with hot matter containing Fe, Mg, O, Si, and C. A hint in favor of this is found in the recent detection of O VI quasar absorption systems at low redshifts (Tripp, Savage, & Jenkins 2000). Knowledge of the explosion dynamics and nucleosynthetic products of very massive stars formed from Big Bang debris is fundamental to further progress. We want to thank Roger Blandford and Marc Kamionkowski for their ongoing interest in this work and for generously giving their time and providing, as usual, deep cosmic insights and education. Sir Maarten Schmidt has been kind enough to expose us to the details of quasar distribution. This work was supported in part by the Department of Energy under grant DE-FG02-87ER40328 to Y.-Z. Q. and by NASA under grant NAG5-4083 to G. J. W., Caltech Division Contribution No. 8728(1059).
warning/0006/astro-ph0006436.html	ar5iv	text	# CMB Observables and Their Cosmological Implications ## 1. Introduction With the data from the BOOMERanG (de Bernardis et al. 2000) and MAXIMA (Hanany et al. 2000) experiments, the promise of measuring cosmological parameters from the power spectrum of anisotropies in the cosmic microwave background (CMB) has come substantially closer to being fulfilled. Together they determine the location of the first peak precisely and constrain the amplitude of the power at the expected position of the second peak. The MAXIMA experiment also limits the power around the expected rise to the third peak. These observations strongly constrain cosmological parameters as has been shown through likelihood analyses in multidimensional parameter space with a variety of prior assumptions (Lange et al. 2000; Balbi et al. 2000; Tegmark & Zaldarriaga 2000b; Bridle et al. 2000). While these analyses are complete in and of themselves, the high dimension of the parameter space makes it difficult to understand what characteristics of the observations or prior assumptions are driving the constraints. For instance, it has been claimed that the BOOMERanG data favor closed universes (White et al. 2000; Lange et al. 2000) and high baryon density (Hu 2000; Tegmark & Zaldarriaga 2000b), but the role of priors, notably from the Hubble constant and big-bang nucleosynthesis is less clear. Indeed, whether CMB constraints agree with those from other cosmological observations serves as a fundamental test of the underlying adiabatic cold dark matter (CDM) model of structure formation. In this paper we show that most of the information in the power spectrum from these two data sets can be compressed into four observables. The correlation among cosmological parameters can be understood by studying their effects on the four observables. They can also be used to search for solutions outside the standard model space (e.g. Peebles et al. 2000; Bouchet et al. 2000). As an instructive application of this approach, we consider the space of flat adiabatic CDM models. Approximate flatness is clearly favored by both BOOMERanG and MAXIMA (de Bernardis et al. 2000; Hanany et al. 2000) as well as previous data, notably from the TOCO experiment (Miller et al. 1999), as shown by previous analyses (Lineweaver 1998; Efstathiou et al. 1999; Tegmark & Zaldarriaga 2000a). Our main objective in this application is to clarify the constraints derived from the CMB observations using the likelihood analyses and understand how they might change as the data evolves. Then, with the aid of a few external constraints we map out the allowed region in the plane of the matter density ($`\mathrm{\Omega }_m`$) versus the Hubble constant ($`H_0`$: we use $`h`$ to denote the Hubble constant $`H_0=100h`$ km s<sup>-1</sup>Mpc<sup>-1</sup>). The external constraints which we employ include (i) the rich cluster abundance at $`z0`$, (ii) the cluster baryon fraction, (iii) the baryon abundance from big bang nucleosynthesis (BBN), and (iv) the minimum age of the universe. We also discuss their consistency with other constraints, such as direct determinations of $`H_0`$, $`\mathrm{\Omega }_m`$, and the luminosity distance to high redshift supernovae. All errors we quote in this paper are at 67% confidence, but we consider all constraints at a 95% confidence level. In §2 we start with a statistical analysis of the CMB data. We introduce the four observables and discuss their cosmological implications. In §3, we place constraints on the $`(\mathrm{\Omega }_m,h)`$ plane and discuss consistency checks. In §4, we identify opportunities for future consistency checks and arenas for future confrontations with data. We conclude in §5. The appendix presents convenient formulae that quantify the cosmological parameter dependence of our four characteristic observables in adiabatic CDM models. ## 2. CMB Observables ### 2.1. Statistical Tests With the present precision of the BOOMERanG and MAXIMA observations (see Fig. 1), it is appropriate to characterize the power spectrum with four numbers: the position of the first peak $`\mathrm{}_1`$, the height of the first peak relative to the power at $`\mathrm{}=10`$ $$H_1\left(\frac{\mathrm{\Delta }T_\mathrm{}_1}{\mathrm{\Delta }T_{10}}\right)^2,$$ (1) the height of the second peak relative to the first $$H_2\left(\frac{\mathrm{\Delta }T_\mathrm{}_2}{\mathrm{\Delta }T_\mathrm{}_1}\right)^2,$$ (2) and the height of the third peak relative to the first $$H_3\left(\frac{\mathrm{\Delta }T_\mathrm{}_3}{\mathrm{\Delta }T_\mathrm{}_1}\right)^2,$$ (3) where $`(\mathrm{\Delta }T_{\mathrm{}})^2\mathrm{}(\mathrm{}+1)C_{\mathrm{}}/2\pi `$ with $`C_{\mathrm{}}`$ the power spectrum of the multipole moments of the temperature field. Note that the locations of the second and third peaks are set by their harmonic relation to the first peak \[see Appendix, eq. (A7)\] and so $`H_2`$ and $`H_3`$ are well-defined even in the absence of clear detections of the secondary peaks. One could imagine two different approaches towards measuring these four numbers. We could extract them using some form of parametrized fit such as a parabolic fit to the data (Knox & Page 2000; de Bernardis et al. 2000). Alternately, we could use template CDM models as calculated by CMBFAST (Seljak & Zaldarriaga 1996), and label them by the values of the four observables. We can measure $`\chi ^2`$ for these CDM models and interpret them as constraints in the four observables. Both of these methods give similar results. We chose the second one because it is more stable to changes in the $`\mathrm{}`$ ranges taken to correspond to each peak and incorporates the correct shape of the power spectra for CDM-like models. To determine the position of the first peak, we take the data that fall between $`75<\mathrm{}<375`$ and carry out a $`\chi ^2`$ fitting using a flat model template ($`\mathrm{\Omega }_\mathrm{m}+\mathrm{\Omega }_\mathrm{\Lambda }=1`$ here and below unless otherwise stated) with varying $`h`$ and $`\mathrm{\Omega }_m`$ at the fixed baryon density $`\mathrm{\Omega }_bh^2=0.02`$ and tilt parameter $`n=1`$. We include calibration errors, 10% for BOOMERanG and 4% for MAXIMA. Figure 2 shows $`\mathrm{\Delta }\chi ^2`$ as a function of $`\mathrm{}_1`$ for the BOOMERanG data alone and for the combination of BOOMERanG and MAXIMA. The figure implies $`\mathrm{}_1`$ $`=`$ $`200\pm 8(\mathrm{BOOM}),`$ $`\mathrm{}_1`$ $`=`$ $`206\pm 6(\mathrm{BOOM}+\mathrm{MAX}).`$ (4) Other choices of $`\mathrm{\Omega }_bh^2`$ and $`n`$ for the template parameters modify slightly the value of $`\chi ^2`$ but not $`\mathrm{\Delta }\chi ^2`$ or the allowed region for $`\mathrm{}_1`$. It is noteworthy that adding in the MAXIMA data steepens $`\mathrm{\Delta }\chi ^2`$ on both sides of the minimum despite the preference for $`\mathrm{}_1220`$ in the MAXIMA data alone (Hanany et al. 2000). The fact that both data sets consistently indicate a sharp fall in power at $`\mathrm{}>220`$ increases the confidence level at which a high $`\mathrm{}_1`$ can be rejected. The $`H_1`$ statistic depends both on the acoustic physics that determines the first peak and other processes relevant at $`\mathrm{}10`$. The shape of the template is therefore more susceptible to model parameters. We choose to vary $`n`$ which changes $`H_1`$ and also allow variations in $`\mathrm{\Omega }_\mathrm{\Lambda }`$ so that the position of the first peak can be properly adjusted. The other parameters were chosen to be $`\mathrm{\Omega }_bh^2=0.03`$ and $`\mathrm{\Omega }_mh^2=0.2`$. Using the BOOMERanG and MAXIMA data for $`75<\mathrm{}<375`$ in conjunction with the COBE data, we find $$H_1=7.6\pm 1.4.$$ (5) Other template choices can modify the constraint slightly but the errors are dominated by the COBE $`7\%`$ cosmic variance errors (Bunn & White 1997) and MAXIMA $`4\%`$ calibration errors on the temperature fluctuations. For $`H_2`$, we take the data for $`75<\mathrm{}<600`$ and consider templates from models with varying $`\mathrm{\Omega }_bh^2`$, $`n`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$, where the last parameter is included to ensure that the models reproduce the position of the first peak. In this case, $`\chi ^2`$ as a function of $`H_2`$ minimized over $`\mathrm{\Omega }_\mathrm{\Lambda }`$ exhibits some scatter due to information that is not contained in the ratio of the peak heights (see Fig. 3). Nonetheless, the steep dependence of $`\mathrm{\Delta }\chi ^2`$ on $`H_2`$ indicates that this statistic is robustly constrained against the variation of the template. Taking the outer envelope of $`\mathrm{\Delta }\chi ^2`$, we obtain $`H_2`$ $`=`$ $`0.37\pm 0.04(\mathrm{BOOM}),`$ $`H_2`$ $`=`$ $`0.38\pm 0.04(\mathrm{BOOM}+\mathrm{MAX}).`$ (6) $`H_3`$ is only weakly constrained by the two highest $`\mathrm{}`$ points ($`600<\mathrm{}<800`$) from MAXIMA in conjunction with the first peak data ($`75<\mathrm{}<375`$) from both experiments. We consider templates from models with varying $`\mathrm{\Omega }_mh^2`$, $`n`$, $`\mathrm{\Omega }_bh^2`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$. The latter two parameters are included to ensure that the position of the first peak and the depth of the first trough can be modeled. Minimizing $`\chi ^2`$ over these two parameters, we plot $`\mathrm{\Delta }\chi ^2`$ as a function of $`H_3`$ (see Fig. 4) to obtain the bound $$H_3=0.43\pm 0.07(\mathrm{BOOM}+\mathrm{MAX}).$$ (7) Note that the constraints on $`H_3`$ employ a template-based extrapolation: points on the rise to the third peak are used to infer its height. ### 2.2. Cosmological Implications The values for the four observables we obtained above can be used to derive and understand constraints on cosmological parameters from the experiments. The position of the first peak as measured by $`\mathrm{}_1`$ is determined by the ratio of the comoving angular diameter distance to the last scattering epoch and the sound horizon at that epoch (Hu & Sugiyama 1995). Therefore it is a parameter that depends only on geometry and sound wave dynamics \[see Appendix, eq. (A)\] through $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }`$, $`h`$, $`\mathrm{\Omega }_m`$, $`w`$, and $`\mathrm{\Omega }_bh^2`$, in decreasing order of importance. The equation of state parameter $`w=p_\mathrm{\Lambda }/\rho _\mathrm{\Lambda }`$ with $`p_\mathrm{\Lambda }`$ and $`\rho _\mathrm{\Lambda }`$ the pressure and energy density of the vacuum $`(w=1)`$ or negative-pressure energy; we use $`\mathrm{\Lambda }`$ to refer to either option. The effect of tilt is small \[see eq. (A8)\] $$\frac{\mathrm{\Delta }\mathrm{}_1}{\mathrm{}_1}0.17(n1)$$ (8) and so we neglect it when considering models with $`n1`$. Figure 5 displays the constraints in the $`\mathrm{\Omega }_m\mathrm{\Omega }_\mathrm{\Lambda }`$ plane. Notice that the confidence region is determined not by uncertainties in the measured value of $`\mathrm{}_1`$, but rather the prior assumptions about the acceptable range in $`h`$, $`\mathrm{\Omega }_bh^2`$ and $`w`$ (Lange et al. 2000; Tegmark & Zaldarriaga 2000b). Given the broad consistency of the data with flat models ($`\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_m=1`$) with $`w=1`$, we will hereafter restrict ourselves to this class of models unless otherwise stated. To better understand the dependence on the Hubble constant, we plot in Figure 6 contours of constant $`\mathrm{}_1`$ in the ($`\mathrm{\Omega }_m`$,$`h`$) plane for $`\mathrm{\Omega }_bh^2=0.019`$ and $`\mathrm{\Omega }_bh^2=0.025`$. A higher baryon abundance decreases the sound horizon at last scattering and pushes up the contours in the direction of higher $`\mathrm{\Omega }_m`$, $`h`$. The 95% limit, $`\mathrm{}_1<218`$ from equation (4), excludes the lower left region shown in Figure 6. This limit derived from the curve for $`\mathrm{\Omega }_bh^2=0.019`$ is robust in the sense that this baryon abundance represents the mimimum value allowed by the CMB, as we shall see later. This constraint is summarized approximately by $`\mathrm{\Omega }_mh^{3.8}>0.079`$. The dependence differs from the familiar combination $`\mathrm{\Omega }_mh^2`$ since in a flat universe part of the effect of lowering $`\mathrm{\Omega }_m`$ is compensated for by the raising of $`\mathrm{\Omega }_\mathrm{\Lambda }`$. The lower limit on $`\mathrm{\Omega }_mh^2`$ or equivalently on $`\mathrm{\Omega }_m`$ found in Tegmark & Zaldarriaga (2000b) reflects the prior upper limit on $`h`$, e.g. for $`h<0.8`$, $`\mathrm{\Omega }_mh^2>0.12`$ or $`\mathrm{\Omega }_m>0.18`$. The 95% confidence lower limit of the combined fit is $`\mathrm{}_1>194`$, but it does not give a robust limit in this plane in the absence of an upper limit on $`\mathrm{\Omega }_bh^2`$. The lower limit is also less robust in the sense that it comes mainly from the MAXIMA result whereas both experiments agree on the upper limit in the sense that the addition of MAXIMA only serves to enhance the confidence at which we may regect larger values. For a flat geometry, the position of the first peak is strongly correlated with the age of the universe. The correlation is accidental since $`\mathrm{}_1`$ is the ratio of the conformal ages $`𝑑t/a`$ at last scattering and the present that enters, not simply the physical age today. Nonetheless, Fig. 7 shows that the correlation is tight across the ($`\mathrm{\Omega }_m`$, $`h`$) values of interest. The upper envelope corresponds to the lowest $`\mathrm{\Omega }_m`$ $`(=0.1)`$ and implies $`\left({\displaystyle \frac{t_0}{1\mathrm{G}\mathrm{y}\mathrm{r}}}\right)`$ $``$ $`19.70.155(250\mathrm{}_1)(0.68)^{1+w}`$ (9) $`\left({\displaystyle \frac{\mathrm{\Omega }_bh^2}{0.019}}\right)^{1.7},`$ where we have included a weak scaling with $`w`$. The observations imply $`t_0<1314`$ Gyr if $`\mathrm{\Omega }_bh^20.019`$. While this is a weak constraint given the current observational uncertainties, notice that the central value of the BOOMERanG results $`\mathrm{}_1200`$ would imply $`t_0911`$Gyrs. The $`H_2`$ statistic, the ratio of the heights of the second peak to the first, mainly depends on the tilt parameter and the baryon abundance. This combination is insensitive to reionization, the presence of tensor modes or any effects that are confined to the lowest multipoles. The remaining sensitivity is to $`\mathrm{\Omega }_mh^2`$ and is modest in a flat universe due to the cancellation of two effects \[see Appendix eq. (A19)\]. Figure 8 shows contours of $`H_2=`$const. in the ($`n,\mathrm{\Omega }_bh^2`$) plane for $`\mathrm{\Omega }_mh^2=0.15`$ and $`0.20`$. The result from the previous subsection, $`H_2<0.46`$ at a 95% confidence level is shown with shades, giving a constraint $$\mathrm{\Omega }_bh^2>0.029\left(\frac{\mathrm{\Omega }_mh^2}{0.15}\right)^{0.58}(n1)+0.024.$$ (10) Note that the constraint at $`n=1`$ is approximately independent of $`\mathrm{\Omega }_mh^2`$. This limit agrees well with the those from the detailed likelihood analysis of Tegmark & Zaldarriaga (2000b, c.f. their Fig. 4) under the same assumptions for $`\mathrm{\Omega }_mh^2`$ and supports the claim that $`H_2`$ captures most of the information from the data on these parameters. An upper limit from $`H_2>0.32`$ also exists but is weaker than conservative constraints from nucleosynthesis as we discuss in the next section. We can derive a limit on $`n`$ from the indicator $`H_1`$. The cosmological parameter dependence of $`H_1`$ is more complicated than the other two we discussed above. Fortunately, most complications tend to decrease $`H_1`$ by adding large angle anisotropies. The lower limit on $`n`$ from the lower limit on $`H_1`$ in the absence of e.g. reionization or tensor modes is therefore conservative. The upper limit on $`n`$ is very weak unless one excludes the possibility of tensor modes as a prior assumption (c.f. Balbi et al. ). We search for the minimum $`n`$ that gives $`H_1`$ larger than the 2$`\sigma `$ lower limit $`H_1>4.8`$ along the parameter space that maximizes $`H_2`$. This gives a conservative lower bound of $`n>0.85`$. This bound is to good approximation independent of $`\mathrm{\Omega }_mh^2`$ for $`\mathrm{\Omega }_mh^20.15`$. Below this value, the bound tightens marginally due to the integrated Sachs-Wolfe (ISW) effect on COBE scales, but in such a way so as to maintain the bound $`\mathrm{\Omega }_bh^2>0.019`$ when combined with the constraint from $`H_2`$, the inequality (10). The analysis of Tegmark & Zaldarriaga (2000b) yields $`n>0.87`$ in the same parameter and data space indicating that not much information is lost in our much cruder parameterization. The $`H_3`$ statistic depends more strongly on $`\mathrm{\Omega }_mh^2`$ and $`n`$ since the baryons affect the height of the third and first peak similarly. In Fig. 9, we show the constraint in the ($`n`$,$`\mathrm{\Omega }_mh^2`$) plane with the baryon density fixed by requiring $`H_2=0.38`$. When combined with the constraint on the tilt $`n>0.85`$, we obtain $`\mathrm{\Omega }_mh^2<0.42`$. In Fig. 1 (dashed lines), we compare a model designed to have acceptable values for $`\mathrm{}_1`$, $`H_1`$, $`H_2`$ and $`H_3`$ with the power spectrum data from COBE, BOOMERanG and MAXIMA. This gives $`\chi ^2=27.2`$ for 30 data points and compares well with the best fit model of Tegmark & Zaldarriaga (2000b) with their “inflation prior” where $`\chi ^2=26.7`$. We summarize the constraints from the CMB as: $`\mathrm{\Omega }_mh^{3.8}`$ $`>0.079`$ (or $`t_0<14`$ Gyr); $`n>0.85`$; $`\mathrm{\Omega }_bh^2>0.019`$; $`\mathrm{\Omega }_mh^2<0.42`$. ## 3. External Constraints In this section, we combine the constraints from the CMB with those from four other observations: the light element abundances as interpreted by BBN theory, the present-day cluster abundance, the cluster baryon fraction and age of the universe. We then translate these constraints onto the ($`\mathrm{\Omega }_m`$,$`h`$) plane and discuss consistency checks. ### 3.1. Nucleosynthesis The first external constraint we consider is that on baryon abundance from primordial nucleosynthesis (see Copi et al. 1995; Olive et al. 1999; Tytler et al. 2000 for recent reviews). Olive et al. (1999) give a high baryon density option of $`0.015\mathrm{\Omega }_bh^20.023`$, and a low baryon density option $`0.004\mathrm{\Omega }_bh^20.010`$ as a $`2\sigma `$ range. A low baryon density is indicated by the traditional low value of helium abundance ($`Y_p=0.234\pm 0.003`$) (Olive et al. 1999), and agrees with a literal interpretation of the lithium abundance. There are also two Lyman limit systems which taken at face value point to high deuterium abundance (Songaila et al 1994; Carswell et al. 1994; Burles et al. 1999a; Tytler et al. 1999) and implies a low baryon density. Our lower limit on $`\mathrm{\Omega }_bh^2>0.019`$ from $`H_2`$ and $`H_1`$ is strongly inconsistent with the low baryon abundance option. In fact our limit is only marginally consistent with even the high baryon option if we take the latest determination of the deuterium abundance at face value and treat the individual errors on the systems as statistical: D/H$`=(3.4\pm 0.25)\times 10^5`$ from 3 Lyman limit systems (Kirkman et al. 2000) which implies $`\mathrm{\Omega }_bh^2=0.019\pm 0.0012`$ (Tytler et al. 2000; Burles et al. 1999b). In this paper, we provisionally accept the 2$`\sigma `$ limit of Olive et al. (1999). The issue, however, is clearly a matter of systematic errors, and we discuss in what follows where they could appear, trying to find a conservative upper limit. Since it is possible that the measured D/H abundance is high due to contamination by $`H`$, we consider the firm lower limit on the D/H abundance from interstellar clouds. The earlier UV data (McCullough 1992) show a variation of D/H from 1.2 to $`2.5\times 10^5`$. This variation is confirmed by modern high resolution spectrographs. The clouds studied are still few in number and range from D/H=(1.5$`\pm 0.1)\times 10^5`$ (Linsky 1998; Linsky et al. 1995) to $`0.7\times 10^5`$ (Jenkins et al. 1999). This variation is reasonable since the clouds are contaminated by heavy elements, indicating significant astration effects. Therefore, we take the upper value as the observational D/H abundance, and take the minimum astration effect (factor 1.5) from model calculations (see Tosi 1996; Olive et al. 1999, modified for a 10 Gyr disk age) to infer the lower limit on the primordial deuterium abundance. We take $`2\times 10^5`$ as a conservative lower limit on D/H. This value agrees with the pre-solar system deuterium abundance inferred from <sup>3</sup>He (Gloeckler & Geiss 1998). This D/H corresponds to $`\mathrm{\Omega }_bh^2=0.028`$ and $`Y_p=0.2500.252`$. The helium abundance $`Y_p`$ directly depends on the theoretical calculation of the helium recombination line, and the discrepancy between the estimate ($`Y_p=0.244\pm 0.002`$, Izotov & Thuan 1998) and the traditional estimates ($`Y_p=0.234\pm 0.003`$) largely arises from the two different calculations, Smits (1996) and Blocklehurst (1972). The helium abundances derived from three recombination lines He I $`\lambda `$4471, $`\lambda `$5876 and $`\lambda `$6678 for a given HII region differ fractionally by a few percent. Also, the effect of underlying stellar absorption by hot stars is unclear: Izotov & Thuan (1998) use the departure of the He I $`\lambda 6678/\lambda 5876`$ strengths from the Smits calculation as an estimator, but a calculation is not actually available for the He line absorption effect. While these variations are usually included as random errors in the nucleosynthesis literature, we suspect that the error in the helium abundance is dominated by systematics and a further change by a few percent in excess of the quoted range is not excluded. The interpretation of the Li abundance rests on a simplistic model of stars. It seems that our understanding of the <sup>7</sup>Li abundance evolution is still far from complete: for instance we do not understand the temperature gradient of the Li/H ratio in halo dwarfs, which shows a trend opposite to what is expected with <sup>7</sup>Li destruction due to diffusion. Hence we do not view the primordial <sup>7</sup>Li abundance determinations as rock-solid. We therefore consider two cases $`\mathrm{\Omega }_bh^2<0.023`$ as a widely accepted upper limit and $`\mathrm{\Omega }_bh^2<0.028`$ as a very conservative upper limit based on interstellar deuterium. When combined with the limit of eq. (10), the latter constraint becomes $`n<1.16`$ for $`\mathrm{\Omega }_mh^2<0.2`$ (as appropriate for setting a lower bound on $`\mathrm{\Omega }_m`$ in the next section, see also Fig. 9); if we instead take $`\mathrm{\Omega }_bh^2<0.023`$ (Olive et al. 1999), the limit becomes $`n<0.98`$. In conjunction with the constraint from $`H_1`$, the allowed range for the tilt becomes $$0.85<n<1.16.$$ (11) ### 3.2. Cluster Abundance The next external constraint we consider is the abundance of clusters of galaxies, which constrains the matter power spectrum at intermediate scales. We adopt the empirical fit of Eke et al. (1996) for a flat universe $`\sigma _8=(0.52\pm 0.08)\mathrm{\Omega }_m^{0.52+0.13\mathrm{\Omega }_m}`$. The value of $`\sigma _8`$ is well converged within 1 $`\sigma `$ among different authors (Viana & Liddle 1999; Pen 1998). This is because the cluster abundance depends strongly on $`\sigma _8`$ due to its appearance in the exponential of a Gaussian in the Press-Shechter formalism. We take the amplitude at COBE scales with a 14 % normalization uncertainty (95 % confidence) together with the 95 % confidence range coming from the cluster abundance to obtain an allowed region that is a function of $`\mathrm{\Omega }_m`$, $`h`$ and $`n`$ and can be roughly described by $$0.27<\mathrm{\Omega }_m^{0.76}hn<0.35,$$ (12) assuming no tensor contribution to COBE and $`\mathrm{\Omega }_bh^2=0.028`$. These assumptions lead to the most conservative constraints on the ($`\mathrm{\Omega }_m`$,$`h`$) plane. The lower limit comes from undershooting $`\sigma _8`$ which is only exacerbated with the inclusion of tensors. It also depends on the upper limit on $`n`$, which is maximized at the highest acceptable baryon density $`\mathrm{\Omega }_bh^2=0.028`$. The upper limit comes from overshooting $`\sigma _8`$ and depends on the lower limit on $`n`$, which only tightens with the inclusion of tensors and lowering of the baryon density. ### 3.3. Baryon Fraction The third external constraint we consider is the baryon fraction in rich clusters, derived from X-ray observations. The observed baryon fraction shows a slight increase outwards, and the true baryon fraction inferred for the entire cluster depends on the extrapolation. The estimates range from $`(0.052\pm 0.0025)h^{3/2}`$ (White & Fabian 1995; lowest estimate) to $`(0.076\pm 0.008)h^{3/2}`$ (Arnaud & Evrard 1999; highest estimate) for rich clusters. We take the 2$`\sigma `$ limits to correspond to these two extreme values. We remark that very similar constraints are derived from the Sunyaev Zeldovich effect for clusters as long as $`h=0.51.0`$: Myers et al. (1997) derive $`(0.061\pm 0.011)h^1`$, and Grego et al. (2000) give $`(0.074\pm 0.009)h^1`$. Adding baryons locked into stars to that in gas inferred by $`X`$-ray observations, and assuming the cluster baryon fraction represents the global value (White et al. 1993), we have $`f_b\mathrm{\Omega }_b/\mathrm{\Omega }_m`$ constrained as $$0.052h^{3/2}+0.006h^1<f_b<0.076h^{3/2}+0.015h^1$$ (13) This relation is used to convert the constraints on $`\mathrm{\Omega }_bh^2`$ into the $`\mathrm{\Omega }_m`$ vs $`h`$ plane. ### 3.4. Age We take the lower limit on the age of the universe to be $`t_0>11`$Gyr based on stellar evolution. While this is not based on statistical analysis, no authors have ever claimed a cosmic age less than this value (Gratton et al. 1997; Reid 1997; Chaboyer et al. 1998). ### 3.5. Allowed Region We display all our constraints in the ($`\mathrm{\Omega }_m`$,$`h`$) plane in Fig. 10. Combining the range $`0.019\mathrm{\Omega }_bh^20.028`$ from BBN and the CMB, together with the constraint on $`\mathrm{\Omega }_b/\mathrm{\Omega }_m`$ from the baryon fraction (13) leads to the range $$\frac{0.019}{0.076h^{1/2}+0.015h}<\mathrm{\Omega }_m<\frac{0.028}{0.052h^{1/2}+0.006h},$$ (14) which is plotted as the solid contours labeled $`f_b`$ in Fig. 10. We also plot the more conservative limits derived from taking the $`2\sigma `$ extremes of the extreme baryond fraction measures ($`0.0760.092`$ and $`0.0520.047`$) as dashed lines. We convert the cluster abundance constraint using the range in tilts acceptable from the CMB constraints and the limit $`\mathrm{\Omega }_bh^2<0.028`$ from BBN ($`0.85<n<1.16`$) and find $$0.15<\mathrm{\Omega }_mh^{1.3}<0.32.$$ (15) This range is displayed by the contours labelled “$`\sigma _8`$”. Finally the constraints $`t_0>11`$Gyr and $`\mathrm{}_1<218`$ are labelled as “$`t_0`$” and “$`\mathrm{}_1`$” respectively. The shading indicates the parameter space within which a model consistent with the CMB and external constraints may be constructed. Dark shading indicates the region that is also consistent with the stronger nucleosynthesis bound of $`\mathrm{\Omega }_bh^2<0.023`$. This does not mean that all models in this region are consistent with the CMB data. To construct a viable model for a given ($`\mathrm{\Omega }_m`$,$`h`$) in this region, one picks a tilt $`n`$ between $`0.851.16`$ consistent with the cluster abundance constraint (12) and then a baryon density consistent with $`H_2`$ (10) and $`\mathrm{\Omega }_bh^2<0.028`$ (or $`0.023`$). In Fig. 1 (solid lines), we verify that the power spectrum prediction of a model so constructed is a good fit to the data. Here $`\chi ^2=28.5`$ for the 30 data points to be compared with $`\chi ^2=27.2`$ for the model optimized for the CMB alone. ### 3.6. Consistency Checks There are a variety of other cosmological measurements that provide alternate paths to constraints in the ($`\mathrm{\Omega }_m`$,$`h`$) plane. We do not use these measurements as constraints since a proper error analysis requires a detailed consideration of systematic errors that is beyond the scope of this paper. We instead use them as consistency checks on the adiabatic CDM framework. Hubble constant. A combined analysis of secondary distance indicators gives $`h=0.71\pm 0.04`$ for an assumed LMC distance of 50 kpc (Mould et al. 2000). Allowing for a generous uncertainty in the distance to the LMC (see Fukugita 2000 for a review) these values may be multiplied by $`0.951.15`$ and this should be compared with our constraint of $`0.6<h<0.9`$. Cosmic acceleration. The luminosity distance to distant supernovae requires $`\mathrm{\Omega }_m<0.48`$ for flat $`\mathrm{\Lambda }`$ models if the systematic errors are no worse than they are claimed (Riess et al 1998; Perlmutter et al. 1999). This limit should be compared with our constraint of $`\mathrm{\Omega }_m<0.6`$. Mass-to-Light Ratio. The $`\mathrm{\Omega }_m`$ constraint we derived using the range $`0.019<\mathrm{\Omega }_bh^2<0.023`$ (see Fig. 10, dark shaded region) and the cluster baryon fraction corresponds to $`M/L_B=(350600)h`$, which is roughly consistent with $`M/L_B`$ for rich clusters (e.g. Carlberg et al. 1997). A yet larger $`\mathrm{\Omega }_m`$ ($`\mathrm{\Omega }_m>0.45`$) would imply the presence of a substantial amount of matter outside clusters and galaxies, whereas we have some evidence indicating the contrary (Kaiser et al. 1998). Power Spectrum. The shape parameter of the transfer function is $`\mathrm{\Gamma }\mathrm{\Omega }_mh\mathrm{exp}[\mathrm{\Omega }_b(11/\mathrm{\Omega }_m)]0.220.33`$ for our allowed region (Sugiyama 1995). This is close to the value that fits the galaxy power spectrum; $`\mathrm{\Gamma }=0.20.25`$ (Efstathiou et al. 1990; Peacock & Dodds 1994). On smaller scales, the Ly$`\alpha `$ forest places constraints on the amplitude and slope of the power spectrum near $`k1`$ $`h`$ Mpc<sup>-1</sup> at $`z3`$ (Croft et al. 1999; McDonald et al 1999). McDonald et al (1999) map these constraints onto cosmological parameters within $`\mathrm{\Lambda }`$CDM as $`n=0.93\pm 0.10`$ and $`\sigma _8=0.68+1.16(0.95n)\pm 0.04`$. Cluster Abundance Evolution. The matter density $`\mathrm{\Omega }_m`$ can be inferred from evolution of the rich cluster abundance (Oukbir & Blanchard 1992), but the result depends sensitively on the estimates of the cluster masses at high redshift. Bahcall & Fan (1998) argue for a low density universe $`\mathrm{\Omega }_m=0.2\genfrac{}{}{0pt}{}{+0.3}{0.1}`$; Blanchard & Bartlett (1998) and Reichart et al. (1999) favor a high density $`\mathrm{\Omega }1`$, while Eke et al. (1998) obtain a modestly low density universe $`\mathrm{\Omega }_m=0.36\pm 0.25`$. Peculiar Velocities. The results from peculiar velocity flow studies are controversial: they vary from $`\mathrm{\Omega }_m=0.15`$ to 1 depending on scale, method of analysis and the biasing factor (see e.g. Dekel 1999 for a recent review). Local Baryons. The CMB experiments require a high baryon abundance. The lower limit (together with a modest red-tilt of the spectrum) is just barely consistent with the high baryon abundance option from nucleosynthesis. The required baryon abundance is still below the maximum estimate of the baryon budget in the local universe $`0.029h^1`$ (Fukugita et al. 1998), but this requires 3/4 of the baryons to reside near groups of galaxies as warm and cool gas. ## 4. Future Directions and Implicit Assumptions A useful aspect of our approach is that one can ask how the allowed parameter space might evolve as the data evolves. More specifically, what aspect of the data can make the allowed region qualitatively change or vanish altogether? If the data are taken at face value, what theoretical assumptions might be modified should that come to pass? An increase in the precision with which the acoustic scale is measured may lead to a new age crisis. It is noteworthy that the secondary peaks will eventually provide a substantially more precise determination of the scale due to sample variance limitations per patch of sky, the multiplicity of peaks, and the effects of driving forces and tilt on the first peak \[see Appendix eq. (A7)\]. Indeed, consistency between the determinations of this scale from the various peaks will provide a strong consistency check on the underlying framework. If the measurements were to determine an equivalent $`\mathrm{}_1200`$, then $`t_0<1011`$Gyrs in a flat $`\mathrm{\Lambda }`$ cosmology with $`\mathrm{\Omega }_bh^2=0.019`$; taking $`\mathrm{\Omega }_bh^2=0.03`$ decreases the age by 1 Gyr and exacerbates the problem. Such a crisis, should it occur, can only be mildly ameliorated by replacing the cosmological constant with a dynamical “quintessence” field. Because increasing the equation of state $`w`$ from $`1`$ reduces both $`\mathrm{}_1`$ and the age, only a relatively extreme choice of $`w1/3`$ can help substantially \[see eqn. (9)\]. This option would also imply that the universe is not accelerating and is in conflict with evidence from distant supernovae. However, other solutions may be even more unpalatable: a small positive curvature and a cosmological constant or a delay in recombination. As constraints on the tilt improve by extending the dynamic range of the CMB observations and those on $`H_2`$ by resolving the second peak, one might be faced with a baryon crisis. Already $`\mathrm{\Omega }_bh^2=0.019`$ is only barely allowed at the 95% CL. Modifications of big-bang nucleosynthesis that allow a higher baryon density for the same deuterium abundance are difficult to arrange: current directions of study include inhomogeneous nucleosynthesis (e.g. Kainulainen et al. 1999) and lepton asymmetry (e.g. Lesgourgues & Peloso 2000; Esposito et al. 2000). On the CMB side, there are two general alternatives. The first possibility is that there is a smooth component that boosts the relative height of the first peak (Bouchet et al. 2000). That possibility can be constrained in the same way as tilt: by extending the dynamic range, one can distinguish between smooth and modulated effects. The direct observable in the modulation is the ratio of energy densities in non-relativistic matter that is coupled to the CMB versus the CMB itself \[$`R_{}`$, see eqn. (A4)\], times the gravitational potential, all evaluated at last scattering. The second possibility is that one of the links in the chain of reasoning from the observables to the baryon and matter densities today is broken in some way. It is noteworthy that there is no aspect of the CMB data today that strongly indicates missing energy in the form of a cosmological constant or quintessence. An Einstein-de Sitter universe with a high baryon density is still viable unless external constraints are introduced. Under the assumption of a flat $`\mathrm{\Lambda }`$ cosmology, tight constraints on $`\mathrm{\Omega }_mh^{3.8}`$ from the peak locations and $`\mathrm{\Omega }_mh^2`$ from the third and higher peak heights should allow $`\mathrm{\Omega }_m`$ and $`h`$ to be separately measured. It will be important to check whether the CMB implications for $`\mathrm{\Omega }_m`$ are consistent with external constraints. Aside from acceleration measurements from distant supernovae, the missing energy conclusion finds its strongest support from the cluster abundance today through $`\sigma _8`$ and the cluster baryon fraction. Changes in the interpretation of these measurements would affect the viability of the Einstein-de Sitter option. The interpretation of the cluster abundance is based on the assumption of Gaussian initial conditions and the ability to link the power spectrum today to that of the CMB through the usual transfer functions and growth rates. One possibility is that the primordial power spectrum has strong deviations from power-law behavior (e.g. Adams et al. 1997). Just like tilt, this possibility can be constrained through the higher peaks. A more subtle modification would arise if the neutrinos had a mass in the eV range. Massive neutrinos have little effect on the CMB itself (Dodelson et al. 1996; Ma & Bertschinger 1995) but strongly suppress large scale structure through growth rates (Jing et al. 1993; Klypin et al. 1993). A total mass (summed over neutrino species) of $`m_{\nu _i}=1`$eV would be sufficient to allow an Einstein-de Sitter universe in the cluster abundance. One still violates the cluster baryon fraction constraint. In fact, even for lower $`\mathrm{\Omega }_m`$ one can only find models consistent with both the cluster abundance and baryon fraction if $`m_{\nu _i}<4`$eV. These constraints could be weakened if some unknown form of support causes an underestimate of the dark mass in clusters through the assumption of hydrostatic equilibrium. They could also be evaded if modifications in nucleosynthesis weaken the upper limit on the baryons. ## 5. Conclusions We find that the current status of CMB power spectrum measurements and their implications for cosmological parameters can be adequately summarized with four numbers: the location of the first peak $`\mathrm{}_1=206\pm 6`$ and the relative heights of the first three peaks $`H_1=7.6\pm 1.4`$, $`H_2=0.38\pm 0.04`$ and $`H_3=0.43\pm 0.07`$. When translated into cosmological parameters, they imply $`\mathrm{\Omega }_mh^{3.8}>0.079`$ (or $`t_0<1314`$ Gyr), $`n>0.85`$, $`\mathrm{\Omega }_bh^2>0.019`$, $`\mathrm{\Omega }_mh^2<0.42`$ for flat $`\mathrm{\Lambda }`$CDM models. Other constraints mainly reflect the implicit (with priors) or explicit use of information from other aspects of cosmology. For example, our consideration of nucleosynthesis, the cluster abundance, the cluster baryon fraction, and the age of the universe leads to an allowed region where $`0.6<h<0.9`$, $`0.25<\mathrm{\Omega }_m<0.45`$, $`0.85<n<0.98`$, $`0.019<\mathrm{\Omega }_bh^2<0.023`$. The region is narrow, but there clearly are adiabatic CDM models viable at the $`95\%`$ CL as exemplified in Fig. 1. The region widens and the quality of the fit improves if one allows somewhat higher baryons $`\mathrm{\Omega }_bh^2<0.028`$ as discussed in this paper. With this extension the tilt can be larger than unity $`n<1.16`$ and $`\mathrm{\Omega }_m`$ as high as 0.6. We note that in both cases our limits reflect conservative assumptions about tensors and reionization, specifically that they are negligible effects in the CMB.<sup>1</sup><sup>1</sup>1This assumption is not conservative when considering likelihood constraints from the CMB alone. The presence of tensors substantially weaken the upper limit on $`n`$. The constraints on these and other CMB observables are expected to rapidly improve as new data are taken and analyzed. We have identified sets of observables that should provide sharp consistency tests for the assumptions that underly their translation into cosmological parameters in the adiabatic CDM framework. With the arrival of precision data sets, the enterprise of measuring cosmological parameters from the CMB has entered a new era. Whether the tension between the observations that is confining the standard parameters to an ever tightening region is indicating convergence to a final solution or hinting at discord that will challenge our underlying assumptions remains to be seen. We would like to thank M. Hudson, S. Landau, G. Steigman, M. Turner, and S. Weinberg for useful discussions. WH is supported by the Keck Foundation; MF by the Raymond and Beverly Sackler Fellowship in Princeton; MZ by the Hubble Fellowship HF-01116-01-98A from STScI, operated by AURA, Inc. under NASA contract NAS5-26555; MT by NASA grant NAG5-9194 and NSF grant AST00-71213. ## Appendix A Scaling Relations The phenomenology of the peaks can be understood through three fundamental scales which vary with cosmological parameters: the acoustic scale $`\mathrm{}_A`$, the equality scale $`\mathrm{}_{\mathrm{eq}}`$ and the damping scale $`\mathrm{}_D`$. We begin by employing an idealized picture of the photon-baryon fluid before recombination that neglects dissipation and time variation of both the sound speed $`c_s`$ and the gravitational driving forces. Simple acoustic physics then tells us that the effective temperature perturbation in the wavemode $`k`$ oscillates as (Hu & Sugiyama 1995) $$\mathrm{\Delta }T(\eta _{},k)=[\mathrm{\Delta }T(0,k)+R_{}\mathrm{\Psi }]\mathrm{cos}(ks_{})R_{}\mathrm{\Psi }.$$ (A1) where the sound horizon at the last scattering surface $`sc_s𝑑\eta =c_s𝑑t/a`$ with $`c_s^2=1/3(1+R)`$ and $`R=3\rho _b/4\rho _\gamma `$. $`\mathrm{\Psi }`$ is the gravitational potential. The asterisk denotes evaluation at last scattering. Baryons modulate the amplitude of the oscillation by shifting the zero point by $`R_{}\mathrm{\Psi }`$. The result is that the modes that reach maximal compression inside potential wells at last scattering are enhanced over those that reach maximal rarefaction. Note that this amplitude modulation is not equivalent to saying that the hot spots are enhanced over cold spots as the same reasoning applies to potential “hills”. The oscillator equation (A1) predicts peaks in the angular power spectrum at $`\mathrm{}_m=m\mathrm{}_A`$ where $`\mathrm{}_A`$ is related to $`s_{}`$ through its projection on the sky today via the comoving angular diameter distance (Hu & Sugiyama 1995) $`D`$ $``$ $`2{\displaystyle \frac{[1+\mathrm{ln}(1\mathrm{\Omega }_\mathrm{\Lambda })^{0.085}]^{1+1.14(1+w)}}{\sqrt{\mathrm{\Omega }_mH_0^2\mathrm{\Omega }_\mathrm{t}^{(1\mathrm{\Omega }_\mathrm{\Lambda })^{0.76}}}}}`$ (A2) $``$ $`{\displaystyle \frac{2d}{\sqrt{\mathrm{\Omega }_mH_0^2}}},`$ where $`\mathrm{\Omega }_\mathrm{\Lambda }`$ refers to the density in dark energy with a fixed equation of state $`w=p_\mathrm{\Lambda }/\rho _\mathrm{\Lambda }`$ ($`w=1`$ for a true cosmological constant) and the total density is $`\mathrm{\Omega }_\mathrm{t}=\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }`$. For convenience, we have defined the dimensionless angular diameter distance $`d`$ which scales out the effect of the expansion rate during matter domination; hence it is equal to unity for an Einstein-de Sitter cosmology. More specifically, $`\mathrm{}_A\pi D/s_{}`$ or $`\mathrm{}_A`$ $``$ $`172d\left({\displaystyle \frac{z_{}}{10^3}}\right)^{1/2}`$ $`\times \left({\displaystyle \frac{1}{\sqrt{R_{}}}}\mathrm{ln}{\displaystyle \frac{\sqrt{1+R_{}}+\sqrt{R_{}+r_{}R_{}}}{1+\sqrt{r_{}R_{}}}}\right)^1,`$ where the radiation-to-matter and baryon-to-photon ratios at last scattering are $`r_{}`$ $``$ $`\rho _r(z_{})/\rho _m(z_{})=0.042\omega _m^1(z_{}/10^3),`$ $`R_{}`$ $``$ $`3\rho _b(z_{})/4\rho _\gamma (z_{})=30\omega _b(z_{}/10^3)^1`$ (A4) with a redshift of last scattering given by $`z_{}`$ $``$ $`1008(1+0.00124\omega _b^{0.74})(1+c_1\omega _m^{c_2}),`$ $`c_1`$ $`=`$ $`0.0783\omega _b^{0.24}(1+39.5\omega _b^{0.76})^1,`$ $`c_2`$ $`=`$ $`0.56(1+21.1\omega _b^{1.8})^1.`$ (A5) Here we use the shorthand convention $`\omega _b=\mathrm{\Omega }_bh^2`$ and $`\omega _m=\mathrm{\Omega }_mh^2`$. Baryon drag works to enhance $`m=`$odd over $`m=`$even peaks in the power. These simple relations are modified by driving and dissipative effects. The driving effect comes from the decay of the gravitational potential in the radiation dominated epoch which enhances the oscillations and leads to an increase in power of approximately a factor of 20 for $`\mathrm{}>\mathrm{}_{\mathrm{eq}}`$ (Hu & Sugiyama 1995) where $$\mathrm{}_{\mathrm{eq}}(2\mathrm{\Omega }_mH_0^2z_{\mathrm{eq}})^{1/2}D438d\omega _m^{1/2}.$$ (A6) It also introduces a phase shift to the oscillations such that the $`m`$th peak of a scale invariant ($`n=1`$ model) is at<sup>2</sup><sup>2</sup>2The coefficients are from fits to the first peak at $`\mathrm{\Omega }_bh^2=0.02`$. For better accuracy, replace the coefficient $`0.267`$ with $`0.24`$ for $`\mathrm{}_2`$ or $`0.35`$ for $`\mathrm{}_3`$. Note that the fractional change made by the phase shift decreases with $`m`$. $`\mathrm{}_m`$ $`=`$ $`\mathrm{}_A(m\varphi )`$ $`\varphi `$ $``$ $`0.267\left({\displaystyle \frac{r_{}}{0.3}}\right)^{0.1}.`$ (A7) Tilt also mildly affects the location of the peaks especially the first which is broadened by radiation effects; around $`n=1`$ (and $`\mathrm{\Omega }_mh^2=0.15`$), the change is approximately $`{\displaystyle \frac{\mathrm{\Delta }\mathrm{}_1}{\mathrm{}_1}}`$ $``$ $`0.17(n1),`$ $`{\displaystyle \frac{\mathrm{\Delta }\mathrm{}_2}{\mathrm{}_2}}`$ $``$ $`0.033(n1),`$ $`{\displaystyle \frac{\mathrm{\Delta }\mathrm{}_3}{\mathrm{}_3}}`$ $``$ $`0.012(n1).`$ (A8) The matter dependence is weak: for $`\mathrm{\Omega }_mh^2=0.25`$ the coefficient $`0.17`$ is reduced to $`0.15`$ for $`\mathrm{}_1`$. The other effect of radiation driving is to reduce the baryon drag effect by reducing the depth of the potential wells at $`z_{}`$. The baryon drag effect is fractionally of order $`R_{}\mathrm{\Psi }(\eta _{})/\mathrm{\Psi }(\eta _{\mathrm{initial}})R_{}T(k)`$ where $`T(k)`$ is the matter transfer function and $`k`$ is the comoving wavenumber in Mpc<sup>-1</sup>. The transfer function quantifies the decay of the potential in the radiation dominated epoch (see e.g. Eisenstein & Hu 1999 for a fit). The break in the transfer function is also given by the horizon scale at matter-radiation equality so that it appears on the sky at $`\mathrm{}_{\mathrm{eq}}`$. Shifting the equality scale to raise $`\mathrm{}_{\mathrm{eq}}`$ by raising the matter content decreases the overall amplitude of the oscillations but increases the odd-even modulation leading to somewhat counterbalancing effects on the peak heights. Finally, the acoustic oscillations are dissipated on small scales. The quantitative understanding of the effect requires numerical calculation but its main features can be understood through qualitative arguments. Since the oscillations dissipate by the random walk of the photons in the baryons, the characteristic scale for the exponential damping of the amplitude is the geometric mean between the mean free path $`\lambda _C=(x_en_e\sigma _Ta)^1`$ and the horizon scale $$\eta _{}=2(\mathrm{\Omega }_mH_0^2)^{1/2}z_{}^{1/2}[\sqrt{1+r_{}}\sqrt{r_{}}].$$ (A9) under the Saha approximation $`x_e\omega _b^{1/2}`$ so that $`k_D(\eta _{}\lambda _C)^{1/2}z_{}^{5/4}\omega _b^{1/4}\omega _m^{1/4}`$. Numerically, the scaling is slightly modified to (refitting values from Hu & White 1997) $$\mathrm{}_Dk_DD\frac{2240d}{[(1+r_{})^{1/2}r_{}^{1/2}]^{1/2}}\left(\frac{z_{}}{1000}\right)^{5/4}\omega _b^{0.24}\omega _m^{0.11}.$$ (A10) Compared with the acoustic scale $`\mathrm{}_A`$, it has a much stronger dependence on $`\omega _b`$ and the redshift of recombination $`z_{}`$. We show a model spectrum<sup>3</sup><sup>3</sup>3The model spectrum is obtained by following a construction based on Hu & White (1997): the damping envelope is $$𝒟_{\mathrm{}}=\mathrm{exp}[(\mathrm{}/\mathrm{}_D)^{1.2}],$$ (A11) yielding acoustic oscillations of the form $$𝒟_{\mathrm{}}𝒜_{\mathrm{}}=[1+R_{}T(\mathrm{}/D)]𝒟_{\mathrm{}}\mathrm{cos}[\pi (\mathrm{}/\mathrm{}_A+\varphi )]R_{}T(\mathrm{}/D);$$ (A12) the potential driving envelope is $$𝒫_{\mathrm{}}1+19\mathrm{exp}(1.4\mathrm{}_{\mathrm{eq}}/\mathrm{}).$$ (A13) The spectrum is then constructed as $`(\mathrm{\Delta }T_{\mathrm{}})^2`$ $``$ $`\left({\displaystyle \frac{\mathrm{}}{10}}\right)^{n1}𝒫_{\mathrm{}}𝒟_{\mathrm{}}^2{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{𝒜_{\mathrm{}}^21}{1+(\mathrm{}_A/2\mathrm{})^6}}+2\right),`$ (A14) where we have added an offset to the oscillations to roughly account for projection smoothing and the Doppler effect and forced the form to return to $`𝒫_{\mathrm{}}`$ above the first peak to account for the early ISW effect from the radiation (Hu & Sugiyama 1995). This mock spectrum should only be used to understand the qualitative behavior of the spectrum. obtained this way for a $`\mathrm{\Lambda }`$CDM cosmology with $`\mathrm{\Omega }_m=0.35`$, $`\mathrm{\Omega }_\mathrm{t}=1`$, $`h=0.65`$, $`\omega _b=0.02`$, and $`w=1`$ in Fig. 11. For this cosmology, the three fundamental scales are $`\mathrm{}_{\mathrm{eq}}=149`$ ($`\mathrm{}_1=221`$), $`\mathrm{}_A=301`$, and $`\mathrm{}_D=1332`$. The dependence of the morphology of the acoustic peaks on cosmological parameters is controlled by these three scales. Around the fiducial $`\mathrm{\Lambda }`$CDM model with the parameters given above $`{\displaystyle \frac{\mathrm{\Delta }\mathrm{}_A}{\mathrm{}_A}}`$ $``$ $`0.11\mathrm{\Delta }w0.24{\displaystyle \frac{\mathrm{\Delta }\omega _m}{\omega _m}}+0.07{\displaystyle \frac{\mathrm{\Delta }\omega _b}{\omega _b}}`$ (A15) $`0.17{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{\Lambda }}{\mathrm{\Omega }_\mathrm{\Lambda }}}1.1{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{t}}{\mathrm{\Omega }_\mathrm{t}}}`$ $``$ $`0.11\mathrm{\Delta }w0.48{\displaystyle \frac{\mathrm{\Delta }h}{h}}+0.07{\displaystyle \frac{\mathrm{\Delta }\omega _b}{\omega _b}}0.15{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_m}{\mathrm{\Omega }_m}}`$ $`1.4{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{t}}{\mathrm{\Omega }_\mathrm{t}}},`$ where the leading order dependence is on $`\mathrm{\Omega }_\mathrm{t}`$ and $`h`$ $`{\displaystyle \frac{\mathrm{\Delta }\mathrm{}_{\mathrm{eq}}}{\mathrm{}_{\mathrm{eq}}}}`$ $``$ $`0.11\mathrm{\Delta }w+0.5{\displaystyle \frac{\mathrm{\Delta }\omega _m}{\omega _m}}0.17{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{\Lambda }}{\mathrm{\Omega }_\mathrm{\Lambda }}}1.1{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{t}}{\mathrm{\Omega }_\mathrm{t}}}`$ (A16) $``$ $`0.11\mathrm{\Delta }w+{\displaystyle \frac{\mathrm{\Delta }h}{h}}+0.59{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_m}{\mathrm{\Omega }_m}}1.4{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{t}}{\mathrm{\Omega }_\mathrm{t}}},`$ which depends more strongly on $`\mathrm{\Omega }_m`$, and $`{\displaystyle \frac{\mathrm{\Delta }\mathrm{}_D}{\mathrm{}_D}}`$ $``$ $`0.11\mathrm{\Delta }w0.21{\displaystyle \frac{\mathrm{\Delta }\omega _m}{\omega _m}}+0.20{\displaystyle \frac{\mathrm{\Delta }\omega _b}{\omega _b}}`$ (A17) $`0.17{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{\Lambda }}{\mathrm{\Omega }_\mathrm{\Lambda }}}1.1{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{t}}{\mathrm{\Omega }_\mathrm{t}}}`$ $``$ $`0.11\mathrm{\Delta }w0.42{\displaystyle \frac{\mathrm{\Delta }h}{h}}+0.20{\displaystyle \frac{\mathrm{\Delta }\omega _b}{\omega _b}}`$ $`0.12{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_m}{\mathrm{\Omega }_m}}1.4{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Omega }_\mathrm{t}}{\mathrm{\Omega }_\mathrm{t}}}.`$ which depends more strongly on the baryon abundance $`\omega _b`$. Note that the sensitivity to $`\mathrm{\Omega }_\mathrm{t}`$ increases from the often quoted $`0.5\mathrm{\Delta }\mathrm{\Omega }_\mathrm{t}/\mathrm{\Omega }_\mathrm{t}`$ as $`\mathrm{\Omega }_\mathrm{\Lambda }`$ increases (Weinberg 2000; M. Turner, private communication). Ideally one would like to extract these three numbers and the baryon-photon ratio $`R_{}`$ directly from the data. The acoustic scale is readily extracted via the position of the first and/or other higher peaks. The other quantities however are less directly related to the observables. We instead choose to translate the parameter dependence into the space of the observations: in particular the height of the first three peaks. The height of the first peak $$H_1\left(\frac{\mathrm{\Delta }T_\mathrm{}_1}{\mathrm{\Delta }T_{10}}\right)^2.$$ (A18) may be raised by increasing the radiation driving force (lowering $`\mathrm{}_{\mathrm{eq}}`$ or $`\omega _m`$) or the baryon drag (raising $`\omega _b`$). However it may also be lowered by filling in the anisotropies at $`\mathrm{}10`$ through the ISW effect (raising $`\mathrm{\Omega }_\mathrm{\Lambda }`$ or $`w`$, or lowering $`\mathrm{\Omega }_\mathrm{t}`$), reionization (raising the optical depth $`\tau `$), or inclusion of tensors. Each of the latter effects leaves the morphology of the peaks essentially unchanged. Because $`H_1`$ depends on many effects, there is no simple fitting formula that describes it. Around the $`\mathrm{\Lambda }`$CDM model with $`H_1=7.4`$, it is crudely $`{\displaystyle \frac{\mathrm{\Delta }H_1}{H_1}}`$ $``$ $`0.5{\displaystyle \frac{\mathrm{\Delta }\omega _m}{\omega _m}}+0.4{\displaystyle \frac{\mathrm{\Delta }\omega _b}{0.02}}0.5\mathrm{\Delta }\mathrm{\Omega }_\mathrm{\Lambda }+0.7\mathrm{\Delta }\mathrm{\Omega }_\mathrm{t}`$ $`2.5\mathrm{\Delta }n1\mathrm{\Delta }\tau 0.3\mathrm{\Delta }w0.76{\displaystyle \frac{\mathrm{\Delta }r}{1+0.76r}}.`$ where the tensor contribution $`r1.4(\mathrm{\Delta }T_{10}^{(\mathrm{T})}/\mathrm{\Delta }T_{10}^{(\mathrm{S})})^2`$. This scaling should only be used for qualitative purposes. The height of the second peak relative to the first is written $$H_2\left(\frac{\mathrm{\Delta }T_\mathrm{}_2}{\mathrm{\Delta }T_\mathrm{}_1}\right)^2\frac{0.925\omega _m^{0.18}(2.4)^{n1}}{[1+(\omega _b/0.0164)^{12\omega _m^{0.52}}]^{1/5}},$$ (A19) where $`n`$ is the scalar tilt and $`\mathrm{}_2/\mathrm{}_12.4`$. This approximation breaks down at high $`\omega _b`$ and $`\omega _m`$ as the second peak disappears altogether. In the $`\mathrm{\Lambda }`$CDM model with $`n=1`$, $`H_2=0.51`$ and parameter variations yield $$\frac{\mathrm{\Delta }H_2}{H_2}0.88\mathrm{\Delta }n0.64\frac{\mathrm{\Delta }\omega _b}{\omega _b}+0.14\frac{\mathrm{\Delta }\omega _m}{\omega _m}.$$ (A20) The effect of tilt is obvious. Baryons lower $`H_2`$ by increasing the modulation that raises all odd peaks. The dependence on the matter comes from two competing effects which nearly cancel around the $`\mathrm{\Lambda }`$CDM: increasing $`\omega _m`$ (lowering $`\mathrm{}_{\mathrm{eq}}`$) decreases the radiation driving and increases $`H_2`$ but also increases the depth of potential wells and hence the modulation that lowers $`H_2`$. For the third peak, these effects add rather than cancel. When scaled to the height of the first peak, which is also increased by raising the baryon density, the $`\mathrm{\Omega }_bh^2`$ dependence weakens leaving a strong dependence on the matter density $`H_3`$ $``$ $`\left({\displaystyle \frac{\mathrm{\Delta }T_\mathrm{}_3}{\mathrm{\Delta }T_\mathrm{}_1}}\right)^2`$ (A21) $``$ $`2.17\left[1+\left({\displaystyle \frac{\omega _b}{0.044}}\right)^2\right]^1\omega _m^{0.59}(3.6)^{n1}`$ $`\times \left[1+1.63\left(1{\displaystyle \frac{\omega _b}{0.071}}\right)\omega _m\right]^1`$ where $`\mathrm{}_3/\mathrm{}_13.6`$. Around the fiducial $`\mathrm{\Lambda }`$CDM model where $`H_3=0.50`$ $$\frac{\mathrm{\Delta }H_3}{H_3}0.41\mathrm{\Delta }n0.31\frac{\mathrm{\Delta }\omega _b}{\omega _b}+0.53\frac{\mathrm{\Delta }\omega _m}{\omega _m}.$$ (A22) We emphasize that the phenomenology in terms of $`\mathrm{}_A`$, $`\mathrm{}_{\mathrm{eq}}`$ and $`\mathrm{}_D`$ is relatively robust, predictive of morphology beyond the first three peaks and readily generalizable to models outside the adiabatic cold dark matter paradigm. The specific scalings of $`H_1`$, $`H_2`$ and $`H_3`$ with cosmological parameters are only valid within the family of adiabatic CDM models. Furthermore, as the data continue to improve, the fits must also be improved from their current few percent level accuracy. The number of phenomenological parameters must also increase to include at least both the heights of the peaks and the depths of the troughs for all observed peaks.
warning/0006/hep-th0006201.html	ar5iv	text	# Sigma-model symmetry in orientifold models ## 1 Introduction Recently, renewed interest has been devoted to orientifold vacua of Type IIB string theory, constructed by projecting out a standard toroidal compactification by the combined action of a discrete spacetime orbifold symmetry $`G`$ and the world-sheet parity $`\mathrm{\Omega }`$ . These unoriented string theories contain both open and closed strings, and constitute the perhaps most important and concrete example of models in which gauge interactions are localized on D-branes . They are therefore the natural arena for the realization of the “brane-world” scenario. Furthermore, this kind of models have proven to offer surprisingly attractive possibilities from a phenomenological point of view (see for instance ). In the following, we will be concerned with compact $`D=4`$ $`N=1`$ Type IIB orientifold models . These vacua represent a simple and tractable prototype of more general and possibly non-supersymmetric orientifold models. Some of them are also phenomenologically appealing and constitute a viable alternative to their more traditional heterotic analogues. In fact, a weak-weak Type IIB - heterotic duality has been conjectured for several pairs of vacua<sup>1</sup><sup>1</sup>1Notice that, although Type IIB $`/\mathrm{\Omega }`$ = Type I , this duality is not a trivial consequence of Type I - heterotic duality in $`D=10`$ , because in general $`\mathrm{\Omega }`$ and $`G`$ do not commute (see for a discussion in the $`D=6`$ models of ), and therefore Type IIB $`/\{\mathrm{\Omega },G\}`$ Type I $`/G`$.. In particular, $`𝐙_N`$ orientifolds with $`N`$ odd do contain D9-branes but no D5-branes, and could be dual to the corresponding perturbative $`𝐙_N`$ heterotic orbifold. Models with $`N`$ even do instead contain both D9-branes and D5-branes, and could be dual to heterotic orbifolds with a perturbative sector corresponding to D9-branes and a non-perturbative instantonic sector corresponding to D5-branes . At the classical level, evidence for the duality is suggested by the almost perfect matching of the low-energy spectra and the fact that the orientifold models seem to possess the same classical symmetries as their heterotic companions . In particular, they both possess a so-called “sigma-model” symmetry<sup>2</sup><sup>2</sup>2In the following, we shall often use the abbreviation “sigma” for “sigma-model”., naturally emerging from $`N=1`$ supergravity. More precisely, this symmetry consists of $`SL(2,R)`$ transformations for the untwisted $`T^i`$ moduli and the other chiral superfields, implemented as the combination of a Kähler transformation and a reparametrization of the scalar Kähler manifold. On the heterotic side, a discrete $`SL(2,Z)`$ subgroup of these transformations is known to correspond to the well-known T-duality symmetry, valid to all orders of string perturbation theory, and is therefore expected to be exact. On the orientifold side, instead, sigma-model transformations do not seem to correspond to any known underlying string symmetry, and it is not clear whether the symmetry is exact. At the quantum level, the comparison becomes much more involved and several subtleties arise. In particular, it has been argued in that the one-loop corrected gauge couplings in orientifold models seem to be incompatible with any duality map (see also for further discussion). There is however an important issue which can be addressed even without knowing the detailed duality map: whether or not the classical sigma-model symmetry is anomalous at the quantum level. The latter continuous symmetry is indeed associated to a composite connection in the low-energy effective supergravity theory, and acts as chiral rotations on all the fermions. There are therefore anomalous triangular diagrams involving gluons, gravitons and composite connections, leading in general to a non-vanishing one-loop anomaly. On the heterotic side, this one-loop anomaly is canceled by a universal tree-level Green-Schwarz (GS) mechanism mediated by the dilaton , and the appearance of the appropriate GS term has been explicitly checked through a string theory computation . On the orientifold side, it was proposed in that a similar GS mechanism involving both the dilaton and twisted RR axions could cancel the anomalies. This observation was motivated by the factorizability of mixed sigma-gauge anomalies computed from the low-energy spectra. However, it was subsequently argued in that requiring a similar mechanism also for mixed sigma-gravitational anomalies would lead to an apparent contradiction with the known results for gauge-gravitational anomaly cancellation . The question of whether the sigma-model symmetry is anomalous or not in orientifold models is therefore still unclear and of extreme relevance for their duality to heterotic theories. Note however that even if the presence of anomalies would pose serious problems to the duality, it would not be fatal for the consistency of the orientifold models in themselves<sup>3</sup><sup>3</sup>3Even in the worse case in which all types of mixed anomalies arise, it is always possible to redefine the conserved currents and energy-momentum tensor in such a way to eliminate mixed gauge or gravitational anomalies, and push all the anomaly in the sigma-model symmetry only. The latter is not fatal, since the emerging longitudinal states are composite and not elementary, and so cannot violate unitarity in higher-loop diagrams as would do longitudinal gluons or gravitons resulting from gauge or gravitational anomalies.. The aim of this paper is to study the cancellation of all possible (pure or mixed) sigma-gauge-gravitational anomalies in orientifold models through a string theory computation. Such an analysis is interesting by itself even beyond the context of Type IIB - heterotic duality, since it can provide useful informations about the low-energy effective action. For instance, the GS couplings that will be derived are related by supersymmetry to other couplings in the Lagrangian and determine under suitable assumptions the gauge kinetic functions and the Fayet-Iliopoulos terms. For simplicity, the analysis will be restricted to the models with $`N`$ odd. These are indeed simpler than models with $`N`$ even for a variety of reasons; in particular, they do not present threshold corrections . The only consistent models with $`N`$ odd are the $`𝐙_3`$ and the $`𝐙_7`$ models. We follow the strategy developed in for standard gauge-gravitational anomalies, and compute both the quantum anomaly and the classical inflow in all possible channels. By factorization, it is then possible to extract all the anomalous couplings for D-branes and fixed-points present in each model, and the GS term given by their sum. A major ingredient of our computation is an effective vertex operator for the composite sigma-model connection, which results from a pair of untwisted Kähler moduli. We provide arguments that such a vertex is in fact the same as that of an “internal graviton” associated to a deformation of the Kähler structure of the orbifold respecting its rigid complex-structure. This suggests that there is a close relation between sigma-model symmetry and invariance under reparametrizations of the internal part of the spacetime manifold. In particular, potential anomalies in these symmetries seem to coincide. Assuming the relation above to be valid and using this common vertex, we are able to compute the complete anomaly polynomial as a function of the gauge, gravitational and sigma-model curvatures. We find that all the anomalies are canceled through a GS mechanism mediated by twisted RR axions only, extending the results of for gauge-gravitational anomalies. The dilaton does not play any role in the anomaly cancellation mechanism, contrarily to what proposed in and in agreement with . The results of our string computation disagree with the field theory analysis of on a crucial sign in the contribution of the twisted modulini to the one-loop sigma-gravitational anomaly. Contrary to what assumed in , it seems that these twisted closed string states must have a non-vanishing “effective” modular weight, that is responsible for the full cancellation of all the anomalies. Although we do not have yet a complete understanding of the field theory interpretation of our results and their implications on the low-energy effective action, we believe that they rise some questions about the actual form of the Kähler potential for twisted fields. As far as we know, this potential has not yet been unambiguously determined. The only available proposal about its form is that of , and it was indeed assumed in . However, this potential implies vanishing modular weight for twisted fields, in apparent contradiction with our results, at least if one does not include possible tree-level corrections to it induced by the GS mechanism. Whether our string results might be explained by taking into account the GS terms in the potential proposed in , or they imply a different form for the Kähler potential of twisted fields, has still to be understood . Independently of their actual field theory explanation, we think that our results provide strong evidence for the occurrence of this cancellation mechanism, generalizing it moreover to all the other types of anomalies, like in particular pure sigma-model anomalies. Unfortunately, although we provide several convincing arguments on the correctness of the effective vertex operator for the composite sigma-model connection, a rigorous proof is missing. Therefore, the only safe statement that we are in position to make is that the associated symmetry is preserved at the one-loop level thanks to a generalized GS mechanism. Whether or not this is really the sigma-model symmetry remains strictly speaking to be proven, although we believe that it is quite unlikely that this is not the case since all the anomalies we compute have precisely the structure expected for sigma-gauge-gravitational anomalies. Notice also that thanks to the alternative interpretation of this vertex as an internal graviton, these canceled anomalies can be inequivocably interpreted as relative to internal reparametrizations. As such, they admit a topological interpretation in terms of equivariant indices of the spin and signature complexes, and it is possible to verify the results obtained through the direct string computation by applying suitable index theorems, as we will see. The structure of the paper is the following. In Section 2, we briefly review the notion of sigma-model symmetry. In Section 3, we set up the general strategy of the string computation and propose a possible path-integral derivation of the effective vertex for the composite connection. In Section 4 we perform the string computation on the four surfaces appearing at the one-loop order. In Section 5, we reproduce the same results from a mathematical point of view as topological indices. In Section 6, we discuss in more detail the obtained quantum anomalies and perform the factorization of the classical inflow to get all the RR anomalous couplings and the total GS term. In Section 7, we discuss possible field-theory interpretations of our results and their implications. Finally, we give conclusions in Section 8. In Appendix A, we report useful conventions about $`\vartheta `$-functions. In Appendix B, we discuss the cancellation of anomalies in Type IIB string theory (this completes the analysis in ). Finally, Appendix C contains some useful details about the string computation. ## 2 Sigma-model symmetry In this section, we briefly review some well-known facts about the sigma-model symmetry, and discuss its potential anomalies in $`D=4`$, $`N=1`$ supergravity models. These general concepts are useful for the considerations that will follow, in particular in Section 7. The scalar manifold $``$ of any generic $`D=4`$ $`N=1`$ supergravity model is known to be a Kähler manifold, described by a Kähler potential $`K`$. At the classical level the Lagrangian presents two distinct symmetries (beside possible local gauge symmetries): * Kähler symmetry, under which the Kähler potential transforms as<sup>4</sup><sup>4</sup>4The Lagrangian is invariant under (1) if also the superpotential $`W`$ transforms as $`We^FW`$. Since $`W`$ is irrelevant in the considerations that will follow, we will neglect it. $$\kappa ^2K(\mathrm{\Phi }^M,\overline{\mathrm{\Phi }}^M)\kappa ^2K(\mathrm{\Phi }^M,\overline{\mathrm{\Phi }}^M)+F(\mathrm{\Phi }^M)+\overline{F}(\overline{\mathrm{\Phi }}^M).$$ (1) * Global isometries of $``$, under which $$\varphi ^M\varphi ^M(\varphi ^N).$$ (2) Here $`\mathrm{\Phi }^M`$ and $`\varphi ^M`$ denote all the chiral multiplets in the model and their lowest components, $`F(\mathrm{\Phi })`$ is a generic chiral superfield, and $`\kappa ^2`$ is Newton’s gravitational constant. The fermions $`\psi ^M`$ in the chiral multiplet $`\mathrm{\Phi }^M`$ transform also under (1) and (2). Correspondingly, the fermionic kinetic terms contain a covariant derivative involving the following “Kähler” and “isometry” connections : $`A_\mu ^{(K)}={\displaystyle \frac{i}{2}}\kappa ^2(K_M_\mu \varphi ^MK_{\overline{M}}_\mu \varphi ^{\overline{M}}),`$ (3) $`A_{\mu N}^{(I)M}=i(\mathrm{\Gamma }_{KN}^M_\mu \varphi ^K\mathrm{\Gamma }_{\overline{K}\overline{N}}^{\overline{M}}_\mu \varphi ^{\overline{K}}).`$ (4) Here $`\varphi ^{\overline{M}}\overline{\varphi }^M`$, $`K_M`$ and $`K_{\overline{M}}`$ denote the derivative of $`K`$ with respect to the corresponding fields and $`\mathrm{\Gamma }`$ is the usual Christoffel connection on the Kähler manifold $``$. Notice that the above connections are not new fundamental states, but composites of the scalar fields. At the quantum level, the symmetries associated to (1) and (2) might be spoiled by triangular one-loop graphs involving as external states the connections (3) and (4), as well as gluons or gravitons. A direct evaluation of these mixed anomalies is not an easy issue, because of the compositeness of the connections. One can however use indirect arguments that rely on the similarity of the structure of the associated anomalous one-loop amplitudes with that of standard $`U(1)`$-gauge and $`U(1)`$-gravitational anomalies . We shall briefly review this analogy in the following, focusing on the case in which a single composite connection enters as external state in the anomalous diagram. The considerations made so far are quite general and apply to any $`D=4`$ $`N=1`$ model. We specialize now to the low-energy Lagrangians arising from the Type IIB orientifolds we want to analyze, i.e. the $`𝐙_3`$ and the $`𝐙_7`$ model (see for more details on these string vacua). The massless closed string spectrum of these models contain the gravitational multiplet, a universal chiral multiplet $`S`$, three chiral multiplets $`T^i`$ corresponding to the (complexified) Kähler deformations of the three internal two-tori<sup>5</sup><sup>5</sup>5Actually, additional “off-diagonal” untwisted moduli survive the orientifold projection in the special $`𝐙_3`$ model. We do not consider them here for simplicity, and all the considerations that follow are independent of the presence of these fields., and a given number of chiral multiplets $`M^\alpha `$ arising from the twisted sectors of the orbifold. The open string spectrum (from D9 branes only in these models) contains vector multiplets and three groups of charged chiral multiplets $`C^a`$. In order to distinguish the different coordinates of $``$, we use the index $`M=(i,a,\alpha )`$ for $`T^i,C^a`$ and $`M^\alpha `$ respectively. As we will see in next sections, the dilaton field $`S`$ does not participate at all to the GS mechanism canceling the anomalies, and is inert under any gauge, diffeomorphism or sigma-model transformations. Up to quadratic order in the charged fields, the total Kähler potential of these orientifolds is believed to be <sup>6</sup><sup>6</sup>6See also for further considerations on the Kähler potential of $`D=4`$ orientifold models.: $`\kappa ^2K_{tot}(\mathrm{\Phi }^M,\overline{\mathrm{\Phi }}^M)`$ $`=`$ $`\mathrm{ln}(S+\overline{S}){\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{ln}(T^i+\overline{T}^i)+{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _i^a{\displaystyle \frac{\overline{C}^aC^a}{T^i+\overline{T}^i}}`$ (5) $`+\kappa ^2K^{(M)}(M^\alpha ,\overline{M}^\alpha ,T^i,\overline{T}^i),`$ where $`K^{(M)}`$ is an unknown potential for the twisted fields $`M^\alpha `$. As mentioned in the introduction, the sigma-model symmetry we want to study in these orientifold models is the dual of heterotic T-duality. It acts on the fields $`T^i`$, $`C^a`$ and $`M^\alpha `$ through the following $`SL(2,R)_i`$ transformations (no sum over $`i`$, $`adbc=1`$): $`T^i{\displaystyle \frac{a_iT^iib_i}{ic_iT^i+d_i}},`$ (6) $`C^a{\displaystyle \frac{\delta _i^a}{(ic_iT^i+d_i)}}C^a,`$ (7) $`M^\alpha M^\alpha (M^\beta ,T^i),`$ (8) and similarly for the complex conjugate fields. The transformation (7) leaves the corresponding (third) term of the Kähler potential (5) invariant and (8) is chosen in such way to preserve the last contribution $`K^{(M)}`$. On the other hand, (6) produces a non-trivial transformation of the second term. In total, the complete Kähler potential (5) undergoes the following Kähler transformation under (6), (7) and (8): $$\kappa ^2K_{tot}(\mathrm{\Phi }^M,\overline{\mathrm{\Phi }}^M)\kappa ^2K_{tot}(\mathrm{\Phi }^M,\overline{\mathrm{\Phi }}^M)+\lambda ^i(T^i)+\overline{\lambda }^i(\overline{T}^i),$$ (9) with $$\lambda ^i(T^i)=\mathrm{ln}(ic_iT^i+d_i).$$ (10) The sigma-model symmetry in question is therefore the combination of an isometry and a Kähler transformation, and potential anomalies will therefore involve both connections (3) and (4). In order to be able to derive an explicit formula at least for mixed sigma-gauge/gravitational anomalies, we need to make some extra assumptions on the potential $`K^{(M)}`$ and the transformations (8). We take here the one usually considered in the literature, that indeed holds generically for heterotic models : $`\kappa ^2K^{(M)}(M^\alpha ,\overline{M}^\alpha ,T^i,\overline{T}^i)={\displaystyle \underset{\alpha }{}}{\displaystyle \underset{i=1}{\overset{3}{}}}(T^i+\overline{T}^i)^{n_i^\alpha }\overline{M}^\alpha M^\alpha +\mathrm{},`$ $`M^\alpha (ic_iT^i+d_i)^{n_i^\alpha }M^\alpha ,`$ (11) where the dots stand for possible higher order terms in $`M^\alpha ,\overline{M}^\alpha `$. The numbers $`n_i^\alpha `$ are the so-called “modular weights” of the fields $`M^\alpha `$. It is straightforward to see that for the reparametrizations (6), (7) and (8), and the Kähler transformation (9) and (10) ($`F=\lambda ^i`$), the total connection $`Z_\mu ^MA_\mu ^{(K)}+A_{\mu M}^{(I)M}`$ transforms as a $`U(1)`$ connection<sup>7</sup><sup>7</sup>7In deriving (12) we assumed that the orbifold limit corresponds to $`C^a=M^\alpha =0`$. The orbifold limit, however, is generically assumed to be given by $`m^\alpha =0`$, where the scalars $`m^\alpha `$ belong to the linear multiplets $`L^\alpha `$, dual of the chiral multiplets $`M^\alpha `$ . So we are assuming that at leading order $`m^\alpha =0`$ corresponds to $`M^\alpha =0`$.: $$Z_\mu ^MZ_\mu ^M+(1+2n_i^M)_\mu \mathrm{Im}\lambda ^i,$$ (12) where $`n_i^\alpha `$ are the coefficients defined in (11), $`n_i^a=\delta _i^a`$, and $`n_i^j=2\delta _i^j`$. The sigma-model symmetry can therefore be viewed as a $`U(1)_i`$ symmetry with “modular charge” $`Q_i^M=(1+2n_i^M)`$. The explicit form of $`Z^M`$ and its field-strength $`G^M=dZ^M`$ can be easily evaluated. It is actually convenient to disentangle the modular charges $`Q_i^M`$ from the connection and define the three connections $`Z_{\mu ,i}`$ and their field-strength $`G_{\mu ,i}`$ so that $`Z_\mu ^M=_iQ_i^MZ_{i,\mu }`$ and $`G_{\mu \nu }^M=_iQ_i^MG_{i,\mu \nu }`$. One finds: $`Z_{i,\mu }={\displaystyle \frac{i}{2}}{\displaystyle \frac{_\mu (t^i\overline{t}^i)}{t^i+\overline{t}^i}},`$ (13) $`G_{i,\mu \nu }=2i{\displaystyle \frac{_{[\mu }t^i_{\nu ]}\overline{t}^i}{(t^i+\overline{t}^i)^2}}.`$ (14) Sigma-gauge/gravitational anomalies can then be computed by treating them as $`U(1)_i`$-gauge/gravitational anomalies (in the following denoted briefly by $`FFG_i`$ and $`RRG_i`$ anomalies respectively). Explicit formulae for the anomaly coefficients can be found for example in eqs.(2.8) and (2.12) of . ## 3 Anomalies in orientifold models In this section, we will set up the general strategy for studying all types of anomalies in chiral orientifold models, and investigate their cancellation. We will begin by reviewing the main aspects of the approach developed in for standard anomalies (see also for a similar analysis in non-geometric models), and generalize it to sigma-gauge-gravitational anomalies. To begin, we shall briefly recall some basic but important facts about anomalies for the convenience of the reader. Anomalies in a quantum field theory effective action have to satisfy the Wess-Zumino (WZ) consistency condition. These in turn imply that any anomaly in $`D`$ dimensions is uniquely characterized by a gauge-invariant and closed $`(D+2)`$-form $`I`$. Using the standard WZ-descent notation<sup>8</sup><sup>8</sup>8The invariant closed $`(D+2)`$-form $`I`$ defines locally a non-invariant Chern-Simons $`(D+1)`$-form $`I^{(0)}`$ such that $`I=dI^{(0)}`$, whose gauge variation then defines a $`(D)`$-form $`I^{(1)}`$ through $`\delta I^{(0)}=dI^{(1)}`$.: $`𝒜=2\pi iI^{(1)}`$. The anomaly polynomial $`I`$ is a characteristic class of the gauge and tangent bundles, of degree $`(D+2)/2`$ in the curvature two-forms. ### 3.1 The strategy The cancellation of anomalies in string theory is achieved in a very natural and elegant way, and is intimately related to more general consistency requirements, like modular invariance and tadpole cancellation. Possible anomalies arise exclusively from boundaries of the moduli space of one-loop string world-sheets. Moreover, direct computations have shown that the whole tower of massive string states contribute in general to anomalies in such a way that these vanish for consistent models, even if the massless spectrum is generically anomalous on its own. From a low-energy effective field theory point of view, where massive states are integrated out and only the resulting effective dynamics of the light modes is considered, the total one-loop anomaly is canceled by an exactly opposite anomaly arising in tree-level processes involving the magnetic exchange of tensor fields . This is the celebrated Green-Schwarz (GS) mechanism , and is an absolutely crucial ingredient for the existence of consistent supersymmetric chiral gauge theories in higher dimensions. In the following, we will focus on the $`CP`$-odd part of the one-loop effective action, where anomalies arise. For consistent models, the exact string theory computation is expected to yield a vanishing anomaly. However, as discussed above, this is interpreted as a non-trivial GS mechanism of anomaly cancellation in a low-energy effective theory valid at energies $`E1/\sqrt{\alpha ^{}}`$. In order to get directly this low-energy approximation, one can take the limit $`\alpha ^{}0`$ from the beginning, before integrating over the world-sheet moduli. The motivation to pursue this strategy, instead of the more direct full string theory computation, is threefold. First, the required computations simplify dramatically. Furthermore, one gets an improved understanding of the low-energy mechanism of anomaly cancellation. Finally, one can extract important WZ couplings appearing in the effective action by factorization . Consider now orientifold models. The relevant anomalous string diagrams are the annulus ($`A`$), the Möbius strip ($`M`$) and the Klein bottle ($`K`$). These world-sheet surfaces lead to potential divergences due to possible tadpoles for massless particles propagating in the transverse channel. Consequently, they also lead to potential anomalies. In addition, also the torus ($`T`$) surface can be anomalous, in the limit under consideration. We will see that there are contributions to the anomaly from this diagram, but they turn out to always cancel among themselves. The most general situation which is allowed by the property that anomalous amplitudes are boundary terms in moduli space is the following. The $`A`$, $`M`$ and $`K`$ surfaces are parametrized by a real modulus $`t[0,\mathrm{}]`$. The contribution from the boundary at $`t\mathrm{}`$ is interpreted as the standard quantum anomaly, whereas the contributions from the other boundary at $`t0`$ is interpreted as classical inflow of anomaly. The $`T`$ amplitude is instead parametrized by a complex modulus $`\tau `$, where $``$ is the fundamental domain. Again, the contribution from the component $`_{\mathrm{}}=[1/2+i\mathrm{},1/2+i\mathrm{}]`$ of the boundary $``$ at infinity is interpreted as the standard quantum anomaly, whereas the contribution from the remaining component $`_0`$ should be associated to the classical inflow of anomaly. Summing up, one would therefore get a quantum anomaly $`𝒜=(A+M+K+T)\|_{\mathrm{}}`$ and a classical inflow $`=(A+M+K+T)\|_0`$. It should be however mentioned that the above interpretation for the $`T`$ surface involves some conceptual subtleties related to modular invariance, that might mix different contributions. Luckily, we will see that the $`T`$ amplitude gives a vanishing contribution anyhow: the pieces in the $`_0`$ component cancel pairwise thanks to modular invariance , that still holds in the $`\alpha ^{}0`$ limit, whereas the $`_{\mathrm{}}`$ component vanishes by itself. Moreover, the $`A`$, $`M`$ and $`K`$ contributions are topological and independent of the modulus. Correspondingly, $`𝒜`$ and $``$ are identical to each other and cancel. As last important remark, notice that in four dimensions even in non-planar diagrams the closed string state exchanged in the transverse channel is always on-shell, due to the conservation of momentum. Strictly speacking, this means that the usual argument for the cancellation of anomalies at the string level does not apply in this case, giving further motivation for a detailed analysis. ### 3.2 Set-up of the computation The computation of the $`A`$, $`M`$, $`K`$ and $`T`$ amplitudes proceeds along the lines of , that we shall briefly review and extend. For the time being, we shall assume that the composite connections (13) are described by suitable effective vertex operators, postponing a detailed discussion of this issue to next subsection. An anomaly of the type discussed above, in the $`CP`$-odd part of the effective action, is encoded in a one-loop correlation function in the odd spin-structure on the $`A`$, $`M`$ and $`K`$ surfaces, and in the odd-even and even-odd spin-structures on the $`T`$ surface, involving gluons, gravitons and composite connections. Denoting by $`\rho `$ the modulus of the surface and by $``$ its integration domain, one has on a given surface and spin-structure $$𝒜_{1\mathrm{}n}=_{}𝑑\rho V_{1}^{}{}_{}{}^{}V_2\mathrm{}V_nJ.$$ (15) The insertion of the supercurrent $`J`$ is due to the existence of a world-sheet gravitino zero-mode; more precisely, $`J=T_F+\stackrel{~}{T}_F`$ in the odd spin-structure on the $`A`$, $`M`$ and $`K`$ surfaces, and $`J=T_F,\stackrel{~}{T}_F`$ in the odd-even and even-odd spin-structures respectively on $`T`$. The vertex $`V^{}`$ is taken in the $`1`$-picture in the odd sector and represents an unphysical particle. Taking the latter to be a longitudinally polarized gluon, graviton or composite connection, one computes the variation of the one-loop effective action under gauge, diffeomorphisms or sigma-model transformations. The remaining vertices $`V`$ are taken in the $`0`$-picture and represent physical background gluons, gravitons or composite connections. Thanks to world-sheet supersymmetry and the limit $`\alpha ^{}0`$, one can use effective vertex operators which are simpler to handle. After some formal manipulations, the correlation function above can be rewritten as boundary terms in moduli space $$𝒜_{1\mathrm{}n}=_{}𝑑\rho W_1V_2\mathrm{}V_n,$$ (16) where $`W`$ is an auxiliary vertex defined out of $`V^{}`$ for the unphysical particle. Importantly, the vertices $`V`$’s contain two tangent fermionic zero-modes, whereas $`W`$ does not contain any of them. The insertion of $`W`$, rather than $`V`$, for the unphysical particle representing the gauge variation of the one-loop effective action corresponds to the fact that the anomaly $`𝒜`$ is given by the WZ descent of the anomaly polynomial $`I`$: $`A=2\pi iI^{(1)}`$. More precisely, one can show that the latter is obtained simply by substituting back $`V`$ instead of $`W`$, that is $$I_{1\mathrm{}n}=_{}𝑑\rho V_1V_2\mathrm{}V_n,$$ (17) with the convention of working in two more dimensions and omitting the integration over bosonic zero-modes. Finally, it is possible to define the generating functional of all the possible anomalies by exponentiating one representative vertex for each type of particle and compute the resulting deformed partition function $`Z^{}`$. Finally, the total anomaly polynomial is given just by $$I=_{}𝑑\rho Z^{}.$$ (18) ### 3.3 Effective vertices The fact that one can use effective vertices in the computation of the partition function yielding the anomaly polynomial is due to the $`\alpha ^{}0`$ limit and to certain special properties of correlation functions in supersymmetric spin-structures like those of relevance here. One way to understand this is to notice that the partition functions to be computed are related to topological indices which are almost insensitive to any continuous parameter deformation. From a more technical point of view, there is always a fermionic zero-mode for each spacetime direction. The corresponding Berezin integral in the partition function yields a vanishing result unless the interaction vertices provide one of each fermionic zero-mode. Infact, products of these fermionic zero modes provide a basis of forms of all degrees in the target spacetime, the Berezin integral selecting the appropriate total degree. On general grounds, it is expected that the effective vertices depend only on the corresponding curvature. Since these behave as two-forms, they must be contracted with two tangent fermionic zero-modes. Moreover, the vertices must be world-sheet supersymmetric. Finally, thanks to the $`\alpha ^{}0`$ limit, they cannot contain additional momenta, beside from those defining the curvature. These three basic requirements, together with the index structure of the curvatures and conformal invariance, turn out to severely constrain the effective vertices in each case. For gluons and gravitons, they can be derived in a straightforward way as in , but for the composite connections (13), the analysis is much more involved since the latter are not fundamental fields but composite of the scalar fields of the theory, and there are therefore no vertex operators directly associated to them. Our main observation is that the field-strengths (14) have a quadratic dependence on the untwisted $`t^i`$ and $`\overline{t}^i`$ moduli fluctuations. Correspondingly, suitable amplitudes with the insertion of the vertex operators associated to these scalars should reproduce the insertion of the composite connections (13). The untwisted $`t^i`$ moduli are defined as $$t^i=e^{\varphi _{10}}g_{i\overline{i}}+i\theta _i,$$ (19) where $`\varphi _{10}`$ is the ten-dimensional dilaton, $`g_{i\overline{i}}`$ is the metric component along the $`T_i^2`$ torus and $`\theta _i`$ is a RR axion. The real part of these moduli is therefore represented by a NSNS vertex operator, whereas the imaginary part is described by a RR vertex, involving spin-fields and particularly unpleasant to deal with. Notice for the moment that these vertex operators can provide at most one spacetime fermionic zero-mode. Since physical gluons and gravitons bring each two fermionic zero-modes, correlations with an odd number of moduli vanish, as expected from the fact these should come in pairs reconstructing composite connections. Moreover, in the limit of interest, the correlation functions under analysis factorize into an internal correlation among moduli fields and a spacetime correlation among gluons and gravitons. We now propose an approach to the derivation of the effective vertex for the composite connection, which is not exhaustive but will allow us to emphasize a few important points. Focus for simplicity on a single internal torus only, for which the composite curvature (14) becomes (no sum over the indices) $`G_{i,\mu \nu }=2iK_{i\overline{i}}_{[\mu }t^i_{\nu ]}\overline{t}^i`$, with $`K_{i\overline{i}}=(t^i+\overline{t}^i)^2`$. On general grounds, one expects the moduli to pair and reconstruct only composite curvatures of this form. At leading order in the momenta, the structure of the internal correlation between two moduli must therefore be as follows: $`V_{t^i}(p_1)V_{\overline{t}^i}(p_2)=\alpha _iK_{i\overline{i}}p_{1\mu }t^ip_{2\nu }\overline{t}^i\psi _0^\mu \psi _0^\nu ,`$ (20) $`V_{t^i}V_{t^i}=V_{\overline{t}^i}V_{\overline{t}^i}=0,`$ (21) where $`\alpha _i`$ are some coefficients and $`V_{t^i}`$ and $`V_{\overline{t}^i}`$ are the vertex operators for the scalars $`t^i`$ and $`\overline{t}^i`$. As already mentioned, correlations such as (20) are potentially difficult to compute in orientifold models, because the moduli vertices have a simple NSNS real part, but a complicated RR imaginary part. More precisely, the sigma-model curvature can be rewritten as $`G_{i\mu \nu }=iK_{i\overline{i}}_{[\mu }(t^i\overline{t}^i)_{\nu ]}(t^i+\overline{t}^i)`$, and one has in principle to use one RR vertex $`V_{t^i}V_{\overline{t}^i}`$ and one NSNS vertex $`V_{t^i}+V_{\overline{t}^i}`$. One could then proceed by contracting the NSNS and RR vertex, take the $`\alpha ^{}0`$ limit and try to figure out which is the effective vertex that, inserted in the correlation function, gives the same result. This procedure is however complicated, so we prefer to use a trick that will allow us to deduce the effective vertex in a quicker (although not rigorous) way. The point is that correlations involving only pairs of $`V_{t^i}+V_{\overline{t}^i}`$ vertices are formally proportional to the corresponding correlations involving pairs of $`V_{t^i}+V_{\overline{t}^i}`$ and $`V_{t^i}V_{\overline{t}^i}`$ vertices. Indeed, using (20) and (21), one gets: $$(V_{t^i}\pm V_{\overline{t}^i})(p_1)(V_{t^i}+V_{\overline{t}^i})(p_2)=\alpha _iK_{i\overline{i}}(p_{1\mu }t^ip_{2\nu }\overline{t}^i\pm (12))\psi _0^\mu \psi _0^\nu .$$ Due to the symmetrization in $`12`$, one gets a vanishing result for two NSNS vertices (upper sign), but a non vanishing one for one RR and one NSNS vertices (lower sign). Nevertheless, both of them encode the same non-vanishing coupling $`\alpha _i`$, and by careful inspection it is possible to extract the latter also from the vanishing correlation involving only NSNS vertices, after having recognized the zero corresponding to the unavoidable symmetrization. A convenient way to properly remove the zero is to flip the crucial sign by hand in the final result, reconstructing the sigma-model curvature. A similar analysis goes through for correlations involving more than two moduli. Indeed, as will now become clear, the moduli vertices do indeed always contract in pairs associated to composite curvatures, and all of them can be represented by the NSNS real part, keeping track of the zeroes arising by symmetrization. We are now in position to attempt a derivation of the effective vertex operator for the composite connections (13), by considering a correlation involving an even number of moduli real parts and using the trick discussed above. The corresponding NSNS vertex operator can be easily deduced from (19), and is given by $$V_{t^i}+V_{\overline{t}^i}=(t^i+\overline{t}^i)d^2z(X^i+ip\psi \psi ^i)(\overline{}\overline{X}^i+ip\stackrel{~}{\psi }\overline{\stackrel{~}{\psi }^i})e^{ipX}+\mathrm{c}.\mathrm{c}.,$$ (22) where $`\mathrm{c}.\mathrm{c}.`$ stands for complex conjugate<sup>9</sup><sup>9</sup>9Notice that the vertex (22) is actually the right one for $`g_{i\overline{i}}`$, that differs from $`t^i+\overline{t}^i`$ for a factor $`g_S=e^{\varphi _{10}}`$. This difference, possibly important for a careful understanding and comparison of string and field theory results (see e.g. footnote 12), is however irrelevant for most of the considerations that will follow. Correspondingly, we effectively identify $`g_{i\overline{i}}`$ with $`t^i+\overline{t}^i`$.. This vertex can be further simplified case by case thanks to the limit $`\alpha ^{}0`$, and to the presence of fermionic zero-modes. But contrarily to the simpler case of gluons and gravitons, it might happen that pieces of the vertex which are apparently subleading for small momenta, give nevertheless a leading contribution when contracted. We proceed separately for the $`A`$, $`M`$, $`K`$ and the $`T`$ surfaces. $`A`$, $`M`$ and $`K`$ surfaces In this case, one can start with the following effective vertex: $$V_{\overline{t}^i}+V_{t^i}=ip\psi _0(t^i+\overline{t}^i)d^2z\left[\psi ^i\overline{}\overline{X}^i+\stackrel{~}{\psi }^i\overline{X}^i+\overline{\psi }^i\overline{}X^i+\overline{\stackrel{~}{\psi }^i}X^i+\mathrm{}\right],$$ (23) where the dots represent possibly important fermionic terms, that are difficult to fix unambiguously in the present approach. By exponentiating two of these vertices with momentum $`p_{1,2}^i`$, and performing a shift on the internal fermions, one gets an effective interaction for the internal bosons. Rescaling then $`(X^i,\overline{X}^i)g_{i\overline{i}}^{1/2}(X^i,\overline{X}^i)`$ so that the bosonic kinetic terms are normalized, one finds $$S_{int}=K_{i\overline{i}}(p_{1\mu }t^ip_{2\nu }\overline{t}^i+(12))\psi _0^\mu \psi _0^\nu d^2z[\overline{X}^i(+\overline{})X^i+\mathrm{}].$$ (24) As expected, the factor in front of the effective vertex (24) has precisely the same form as (3.3) with the $`+`$ sign, and this interaction term vanishes due to the $`12`$ symmetrization. According to the previous discussion, by flipping the sign of the second term in the brackets, one generates a non-vanishing interaction which can be interpreted as an effective vertex operator for the composite connection. Notice that one would expect such an effective vertex to be world-sheet supersymmetric, whereas the expression obtained above is not. We conclude from this that the expression (24) is incomplete, and that additional purely fermionic terms must indeed be present in (23) and (24). By requiring a world-sheet supersymmetric vertex, it is then easy to deduce the right form for these fermionic terms, and one finds finally $$V_G^{eff.}=\frac{1}{2}G_{i,\mu \nu }\psi _0^\mu \psi _0^\nu d^2z\left[\overline{X}^i(+\overline{})X^i+(\overline{\psi }\overline{\stackrel{~}{\psi }})^i(\psi \stackrel{~}{\psi })^i\right].$$ (25) $`T`$ surface In this case, on can effectively take: $$V_{\overline{t}^i}+V_{t^i}=ip\psi _0(t^i+\overline{t}^i)d^2z\left[\psi ^i\overline{X}^i+\overline{\psi }^iX^i+\mathrm{}\right].$$ (26) The dots represent again possible fermionic terms. By exponentiating and performing a shift on the left-moving internal fermions, one gets an effective interaction for the bosons given by $$S_{int}=K_{i\overline{i}}(p_{1\mu }t^ip_{2\nu }\overline{t}^i+(12))\psi _0^\mu \psi _0^\nu d^2z[\overline{X}^i\overline{}X^i+\mathrm{}].$$ (27) As before, this interaction term vanishes and one has to perform the discussed sign flip to obtain a non-vanishing interaction to be interpreted as an effective vertex operator for the composite connection. Again, since such effective vertex should be world-sheet supersymmetric, we conclude that (27) is indeed incomplete, and fix again the missing fermionic terms thanks to world-sheet supersymmetry. Finally, one gets $$V_G^{eff.}=\frac{1}{2}G_{i,\mu \nu }\psi _0^\mu \psi _0^\nu d^2z\left[\overline{X}^i\overline{}X^i+\overline{\stackrel{~}{\psi }^i}\stackrel{~}{\psi }^i\right].$$ (28) There is an alternative way to deduce the form of the effective vertices above. Since the NSNS $`\mathrm{Re}t^i`$ scalar is related to the metric of the corresponding internal two-torus, the exponentiation of its vertex induces a geometric deformation of the orbifold along the $`i`$-th internal torus. This can be analyzed directly from a $`\sigma `$-model point of view. By doing that, with standard techniques, it is easy to see that the metric deformation associated to the internal $`T_i^2`$ torus is represented by (25) and (28) on the corresponding surfaces, where $`G_i`$ is now replaced by the geometric curvature of $`T_i^2`$. By exploiting the tensorial structure of this curvature, one easily realizes that the components whose derivatives are all along the spacetime directions, like in (24), vanish due to a symmetrization, exactly like before. As expected, one is therefore led to use the same trick as above to get a non-vanishing composite field-strength. However, in this way one gets automatically the fermionic terms in (25) and (28) and also a first clue of the close relation between the field-strength $`G`$ and the curvature of the internal space. We postpone to Section 5 a more precise analysis of this relationship. Notice also that in heterotic models, where the untwisted moduli consist of NSNS fields only, the correspondence between Kähler deformations of the orbifold and sigma-model symmetry can be unambiguously established. The net result is again that the effective vertex for the composite connection has the same form as that of an internal graviton, like in (28). We think that this gives some extra evidence for the relation between sigma-model symmetry and orbifold Kähler deformations also in Type IIB orientifolds. Indeed, although in the latter case the pseudo-scalars $`\mathrm{Im}t^i`$ are RR fields, from a purely geometrical point of view there is no difference with respect to heterotic models, since in both theories $`\mathrm{Im}t^i`$ simply complexifies the geometric Kähler structure of the orbifold/orientifold. ## 4 String computation The computation of the partition functions entering the anomaly polynomial closely follow . We proceed separately for the various surfaces. The $`A`$, $`M`$ and $`K`$ amplitudes are generalizations of the results of to a non-trivial “composite” background. The $`T`$ amplitude was instead irrelevant in , as shown in Appendix B for the six-dimensional case, and has therefore to be computed in detail. As already said, we restrict to the simplest $`𝐙_3`$ and $`𝐙_7`$ models, which do not contain D5-branes neither $`N=2`$ sub-sectors. In these models, the $`k`$-th element of $`𝐙_N`$ is $`g^k=(\theta ^k,\gamma _k)`$, where $`\theta ^k`$ is a rotation of angles $`2\pi kv_i`$ in the internal two-tori $`i=1,2,3`$, and $`\gamma _k`$ is a non-trivial twist matrix, acting on the Chan-Paton bundle. The Chan-Paton representation of the twist is fixed by the tadpole cancellation condition. For future convenience, and in order to get contact with the notation used in the literature, we define $$C_k=\underset{i=1}{\overset{3}{}}(2\mathrm{sin}\pi kv_i).$$ (29) and its sign $`ϵ_k=\mathrm{sign}C_k`$. For $`N`$ odd, the tadpole cancellation condition can then be written as $$\frac{1}{4}\mathrm{tr}(\gamma _{2k})=\frac{C_{2k}}{C_k}=\frac{C_k}{C_{2k}}$$ (30) and holds because actually all the quantities in the equality are equal to a sign, namely $`ϵ_{2k}/ϵ_k`$ which is equal to $`1`$ for $`𝐙_3`$ and $`+1`$ for $`𝐙_7`$. Let us define the characteristic classes which will appear in the polynomial associated to generic sigma-gauge-gravitational anomalies. For the gauge bundle, one has the natural $`𝐙_N`$ Chern character, function of the gauge curvature $`F`$, defined as a trace over the Chan-Paton representation: $$\mathrm{ch}_k(F)=\mathrm{tr}[\gamma _ke^{iF/2\pi }].$$ (31) This factor appears in the anomaly from charged chiral spinors. For the tangent bundle, the relevant characteristic classes are the Roof-genus, G-polynomial and Hirzebruch polynomial, functions of the gravitational curvature $`R`$ and defined in terms of the skew eigenvalues $`\lambda _a`$ of $`R`$ as: $`\widehat{A}(R)={\displaystyle \underset{a=1}{\overset{D/2}{}}}{\displaystyle \frac{\lambda _a/4\pi }{\mathrm{sinh}\lambda _a/4\pi }},`$ (32) $`\widehat{G}(R)={\displaystyle \underset{a=1}{\overset{D/2}{}}}{\displaystyle \frac{\lambda _a/4\pi }{\mathrm{sinh}\lambda _a/4\pi }}\left(2{\displaystyle \underset{b=1}{\overset{D/2}{}}}\mathrm{cosh}\lambda _b/2\pi \mathrm{\hspace{0.17em}1}\right),`$ (33) $`\widehat{L}(R)={\displaystyle \underset{a=1}{\overset{D/2}{}}}{\displaystyle \frac{\lambda _a/2\pi }{\mathrm{tanh}\lambda _a/2\pi }}.`$ (34) These factors appear respectively in the anomaly from chiral spinors, chiral Rarita-Schwinger fields, and self-dual tensor fields. We also introduce three new characteristic classes depending on the composite curvature $`G=dZ`$, defined in terms of the curvatures $`G_i`$ in the three internal tori as $`\widehat{A}_k(G)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\mathrm{sin}(\pi kv_i)}{\mathrm{sin}(\pi kv_i+G_i/2\pi )}},`$ (35) $`\widehat{G}_k(G)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\mathrm{sin}(\pi kv_i)}{\mathrm{sin}(\pi kv_i+G_i/2\pi )}}\left(2{\displaystyle \underset{j=1}{\overset{3}{}}}\mathrm{cos}(2\pi kv_j+G_j/\pi )\mathrm{\hspace{0.17em}1}\right),`$ (36) $`\widehat{L}_k(G)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\mathrm{tan}(\pi kv_i)}{\mathrm{tan}(\pi kv_i+G_i/\pi )}}.`$ (37) These characteristic classes will appear in the anomaly from states transforming as chiral spinors, Rarita-Schwinger fields and self-dual tensors with respect to sigma-model transformations. A last preliminary comment relevant to all the surfaces is the following. Due to the universal six bosonic zero-modes in the four non-compact spacetime directions and the two extra auxiliary dimensions introduced to deal with the WZ descent, the partition functions will always contain a free-particle contribution proportional to $`\rho ^3`$. Moreover, the curvatures will always appear multiplied by $`\rho `$ as twists in the partition function. An important simplification occurs using the fact that only the 6-form component of the partition function is relevant for our purposes: one can scale out the above explicit dependences on the modulus $`\rho `$. This will be important also for the modular invariance of the string amplitudes yielding the anomaly in the torus surface, as we shall see. ### 4.1 $`A`$, $`M`$ and $`K`$ surfaces On the $`A`$, $`M`$ and $`K`$ surfaces, the boundary of moduli space is given by the component $`t\mathrm{}`$ encoding the quantum anomaly, minus the component $`t0`$ encoding the classical GS inflow. The contribution of each surface to the total anomaly polynomial is given by $$I_\mathrm{\Sigma }=\left(\underset{t\mathrm{}}{lim}\underset{t0}{lim}\right)Z_\mathrm{\Sigma }(t).$$ (38) The partition functions $`Z_\mathrm{\Sigma }(t)`$ are in the RR odd spin-structure, and their operatorial representation is $`Z_A(t)={\displaystyle \frac{1}{4N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}\mathrm{Tr}_R[g^k(1)^Fe^{tH}],`$ $`Z_M(t)={\displaystyle \frac{1}{4N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}\mathrm{Tr}_R[\mathrm{\Omega }g^k(1)^Fe^{tH}],`$ $`Z_K(t)={\displaystyle \frac{1}{8N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}\mathrm{Tr}_{RR}[\mathrm{\Omega }g^k(1)^{F+\stackrel{~}{F}}e^{tH}].`$ (39) Here $`H=H(R,F,G)`$ is the Hamiltonian associated to the two-dimensional supersymmetric non-linear $`\sigma `$-model in a gauge, gravitational and composite background defined by the effective vertex operators below, with Neumann boundary conditions. Due to supersymmetry, (39) are generalized Witten indices in which only massless modes can contribute . Indeed, it can be verified explicitly that massive world-sheet fermionic and bosonic modes exactly cancel. As a consequence, the partition functions (39) are independent of $`t`$, and (38) vanishes, reflecting anomaly cancellation through the GS mechanism. The background dependence of the action is encoded in the effective vertices for external particles. In the odd spin-structure on the $`A`$, $`M`$ and $`K`$ surfaces, the sum $`Q+\stackrel{~}{Q}`$ of the left and right world-sheet supersymmetries is preserved, and there are space-time fermionic zero-modes $`\psi _0^\mu =\stackrel{~}{\psi }_0^\mu `$. In the limit $`\alpha ^{}0`$, we use the following effective vertex operators for gluons, gravitons and composite sigma-model connections: $`V_F^{eff.}=F^a{\displaystyle 𝑑\tau \lambda ^a},`$ (40) $`V_R^{eff.}=R_{\mu \nu }{\displaystyle d^2z\left[X^\mu (+\overline{})X^\nu +(\psi \stackrel{~}{\psi })^\mu (\psi \stackrel{~}{\psi })^\nu \right]},`$ (41) $`V_G^{eff.}=G_i{\displaystyle d^2z\left[\overline{X}^i(+\overline{})X^i+(\overline{\psi }\overline{\stackrel{~}{\psi }})^i(\psi \stackrel{~}{\psi })^i\right]},`$ (42) in terms of the curvature two-forms $$F^a=\frac{1}{2}F_{\mu \nu }^a\psi _0^\mu \psi _0^\nu ,R_{\mu \nu }=\frac{1}{2}R_{\mu \nu \rho \sigma }\psi _0^\rho \psi _0^\sigma ,G_i=\frac{1}{2}G_{i,\mu \nu }\psi _0^\mu \psi _0^\nu .$$ (43) It is now straightforward to compute the partition functions (39) on the $`A`$, $`M`$ and $`K`$ surfaces. The composite background modifies only the internal partition functions, whereas the spacetime contribution has only the standard dependence on the gauge and gravitational backgrounds. The spacetime part can be computed exactly as in , and one finds the same results as in . The computation of the internal part is also similar to that in , the curvature $`G`$ entering as a twist. Using $`\zeta `$-function regularization, one finds $`Z_A={\displaystyle \frac{i}{4N}}{\displaystyle \underset{k=1}{\overset{N1}{}}}C_k\widehat{A}_k(G)\mathrm{ch}_k^2(F)\widehat{A}(R),`$ $`Z_M={\displaystyle \frac{i}{4N}}{\displaystyle \underset{k=1}{\overset{N1}{}}}C_k\widehat{A}_k(G)\mathrm{ch}_{2k}(2F)\widehat{A}(R),`$ $`Z_K={\displaystyle \frac{i}{16N}}{\displaystyle \underset{k=1}{\overset{N1}{}}}C_{2k}\widehat{L}_k(G)\widehat{L}(R),`$ (44) in terms of the characteristic classes defined before. As anticipated, the partition functions (44) are independent of the modulus $`t`$. Consequently, the quantum anomaly encoded in the $`t\mathrm{}`$ boundary, and the classical inflow associated to $`t0`$ boundary, are precisely opposite to each other and cancel on each of the $`A`$, $`M`$ and $`T`$ surfaces. ### 4.2 $`T`$ surface On the $`T`$ surface, the boundary $``$ of moduli space splits into the component at infinity, $`_{\mathrm{}}=[1/2+i\mathrm{},1/2+i\mathrm{}]`$, minus the remaining component, $`_0`$, and the contribution to the total anomaly polynomial is given by $$I_T=\frac{1}{2}\left[\left(_{_{\mathrm{}}}__0\right)d\tau Z_T(\tau )+\left(_{_{\mathrm{}}}__0\right)d\overline{\tau }Z_T(\overline{\tau })\right].$$ (45) The quantities $$Z_T(\tau )=\underset{\alpha }{}(1)^\alpha Z_T^{S_\alpha }(\tau ),Z_T(\overline{\tau })=\underset{\stackrel{~}{\alpha }}{}(1)^{\stackrel{~}{\alpha }}Z_T^{S_{\stackrel{~}{\alpha }}}(\overline{\tau }).$$ (46) are the total partition functions in the odd-even and even-odd sector respectively. More precisely, $`\alpha =2,3,4`$ represent the RR, $`\mathrm{RNS}_+`$ and $`\mathrm{RNS}_{}`$ odd-even spin-structures, and similarly $`\stackrel{~}{\alpha }=2,3,4`$ represent the RR, $`\mathrm{RNS}_+`$ and $`\mathrm{RNS}_{}`$ even-odd spin-structures. Their operatorial representation is $`Z_T^{RR}(\tau )={\displaystyle \frac{1}{8N}}{\displaystyle \underset{k,l=0}{\overset{N1}{}}}\mathrm{Tr}_{RR}^{(l)}[g^k(1)^F\stackrel{~}{g}^ke^{\tau H}e^{\overline{\tau }\stackrel{~}{H}}],`$ $`Z_T^{RNS_+}(\tau )={\displaystyle \frac{1}{8N}}{\displaystyle \underset{k,l=0}{\overset{N1}{}}}\mathrm{Tr}_{RNS}^{(l)}[g^k(1)^F\stackrel{~}{g}^ke^{\tau H}e^{\overline{\tau }\stackrel{~}{H}}],`$ $`Z_T^{RNS_{}}(\tau )={\displaystyle \frac{1}{8N}}{\displaystyle \underset{k,l=0}{\overset{N1}{}}}\mathrm{Tr}_{RNS}^{(l)}[g^k(1)^F\stackrel{~}{g}^k(1)^{\stackrel{~}{F}}e^{\tau H}e^{\overline{\tau }\stackrel{~}{H}}].`$ (47) The expression for the even-odd spin-structures is perfectly similar, with left and right movers exchanged. In this case, $`H`$ and $`\stackrel{~}{H}=\stackrel{~}{H}(R,F,G)`$ are the left and right-moving Hamiltonians associated to the two-dimensional supersymmetric non-linear $`\sigma `$-model in a gauge, gravitational and composite background defined by the effective vertex operators below. Notice that whereas the even part of the partition functions is influenced by the backgrounds, the odd part remains trivial. This will lead to holomorphic and anti-holomorphic results in the odd-even and even-odd spin structures. Furthermore, only the odd parts of (47) are supersymmetric indices, whereas the even parts receive contributions from all the tower of string states and will therefore depend on $`\tau `$. Again, the background dependence of the action is encoded in the effective vertices for external particles. In the odd-even spin-structure on the $`T`$ surface, the left-moving world-sheet supersymmetry $`Q`$ is preserved, and there are space-time fermionic zero-modes $`\psi _0^\mu `$. In the limit $`\alpha ^{}0`$, we use the following effective vertex operators for gravitons and composite connections: $`V_R^{eff.}=R_{\mu \nu }{\displaystyle d^2z\left[X^\mu \overline{}X^\nu +\stackrel{~}{\psi }^\mu \stackrel{~}{\psi }^\nu \right]},`$ (48) $`V_G^{eff.}=G_i{\displaystyle d^2z\left[\overline{X}^i\overline{}X^i+\overline{\stackrel{~}{\psi }^i}\stackrel{~}{\psi }^i\right]},`$ (49) in terms of the curvature two-forms defined in (43). It is then easy to evaluate the partition function on the $`T`$ surface. The gravitational background influences bosons and left-moving fermions, in a similar way to the cases discussed in Appendix B. The composite background influences instead only the internal bosons and left-moving fermions. The evaluation of the internal partition functions is very similar to that reported in Appendix B for the six-dimensional case of Type IIB on $`T^4/𝐙_N`$, the curvature $`G`$ being responsible for a twist. In total, one gets: $`Z_T(R,G,\tau )`$ $`=`$ $`{\displaystyle \frac{i}{8N}}{\displaystyle \underset{\alpha =2}{\overset{4}{}}}(1)^\alpha {\displaystyle \underset{k,l=0}{\overset{N1}{}}}N_{k,l}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\theta _\alpha \left[\genfrac{}{}{0pt}{}{lv_i}{kv_i}\right](G_i/\pi ^2\|\tau )}{\theta _1\left[\genfrac{}{}{0pt}{}{lv_i}{kv_i}\right](G_i/\pi ^2\|\tau )}}`$ (50) $`\times {\displaystyle \underset{a=1}{\overset{2}{}}}\left[{\displaystyle \frac{(ix_a)}{\theta _1(ix_a/\pi \|\tau )}}\theta _\alpha (ix_a/\pi \|\tau )\right]{\displaystyle \frac{\eta ^3(\tau )}{\theta _\alpha (0\|\tau )}},`$ where $`x_a=\lambda _a/2\pi `$ and $`N_{k,l}`$ is the number of fixed-points that are at the same time $`k`$ and $`l`$-fixed ($`N_{0,0}=0`$). The result for the odd-even spin-structures is the complex conjugate of (50). It is a lengthy but straightforward exercise to show that the partition function (50) is modular invariant. Indeed, one gets $`Z_T(R,G,\tau +1)`$ $`=`$ $`Z_T(R,G,\tau )`$ $`Z_T(R,G,1/\tau )`$ $`=`$ $`{\displaystyle \frac{1}{\tau }}Z_T(R\tau ,G\tau ,\tau )=\tau ^2Z_T(R,G,\tau ),`$ (51) where the last step in the second equation is valid for the relevant 6-form component of $`Z_T`$. Thanks to the modular invariance of $`Z_T(\tau )`$ and $`Z_T(\overline{\tau })`$, their integral on various components of $``$ are related to each other. In fact, only the component $`_{\mathrm{}}`$ at infinity can give a non-vanishing contributions, the remaining four pieces of the remaining component $`_0`$ canceling pairwise, as in . The potential contribution from $`_{\mathrm{}}`$ is interpreted as a quantum sigma-gravitational anomaly. In order to evaluate the contribution from $`_{\mathrm{}}`$, one has to take the limit $`\tau _2\mathrm{}`$ of the partition function. This is easy to take in untwisted sectors, but in twisted sectors one has to pay attention to the range of the twists. For $`l0`$, one gets for instance: $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\theta _2\left[\genfrac{}{}{0pt}{}{lv_i}{kv_i}\right](G_i/\pi ^2\|\tau )}{\theta _1\left[\genfrac{}{}{0pt}{}{lv_i}{kv_i}\right](G_i/\pi ^2\|\tau )}}iϵ_l,`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\theta _{3,4}\left[\genfrac{}{}{0pt}{}{lv_i}{kv_i}\right](G_i/\pi ^2\|\tau )}{\theta _1\left[\genfrac{}{}{0pt}{}{lv_i}{kv_i}\right](G_i/\pi ^2\|\tau )}}iϵ_lq^{1/8}{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{exp}(iϵ_lG_i/\pi ).`$ where the quantity $`ϵ_k`$ was defined as the sign of $`C_k`$ in (29)<sup>10</sup><sup>10</sup>10It arises here as $`ϵ_l=(1)^{_i\theta _i(lv_i)}`$ and $`ϵ_l=2_i\theta _i(lv_i)3`$ in terms of the representative $`\theta _i(lv_i)=lv_i\mathrm{int}(lv_i)`$ of the twist $`lv_i`$ in the interval $`[0,1]`$.. The above expressions already show that the anomaly from RR twisted states does not depend on the composite curvature, and therefore trivially vanishes in $`D=4`$. On the contrary, the anomaly from RNS twisted states does depend on the curvature $`G`$, and is non-vanishing. The corresponding apparently complex internal contribution to the partition function (50) turns out to be actually real, and using the fact that only odd powers of $`G_i`$ are relevant in $`D=4`$, one gets: $$iϵ_lq^{1/8}\underset{i=1}{\overset{3}{}}\mathrm{exp}(iϵ_lG_i/\pi )=q^{1/8}(i)\mathrm{ch}(2G).$$ Finally, the total odd-even and even-odd spin structure partition functions (46) are found to behave both in the same way in the limit $`\tau _2\mathrm{}`$, giving: $`Z_T`$ $``$ $`{\displaystyle \frac{i}{16N}}{\displaystyle \underset{k=1}{\overset{N1}{}}}C_{2k}\widehat{L}_k(G)\widehat{L}(R)`$ (52) $`+{\displaystyle \frac{i}{2N}}{\displaystyle \underset{k=1}{\overset{N1}{}}}C_k\left[\widehat{A}_k(G)\widehat{G}(R)+\widehat{G}_k(G)\widehat{A}(R)\right]`$ $`+{\displaystyle \frac{i}{2N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \underset{l=1}{\overset{N1}{}}}N_{k,l}(i)\mathrm{ch}(G)\widehat{A}(R).`$ Since this expression is independent of $`\tau `$, the remaining integral over $`\tau _1`$ in $`_{\mathrm{}}`$ is trivial, and according to (45), this is also the final result for the $`T`$ contribution to the anomaly. The first line of (52) corresponds to the RR untwisted sector, whose contribution precisely cancels that of the Klein bottle in (44). The second line corresponds to the RNS/NSR untwisted sectors and encodes the contributions of the gravitino, dilatino and untwisted modulini. Finally, the last line encodes those of twisted RNS/NSR moduli; notice that all the twisted sectors $`l=1,\mathrm{},N1`$ give the same contribution, since $`N_{k,l}`$ takes the same value for all $`\{k,l\}\{0,0\}`$ for $`N`$ odd. Actually, one can check that the relevant 6-form component of the result (52) vanishes identically. Some useful details in this respect are reported in Appendix C. In conclusion, the total anomaly from the $`T`$ surface exactly cancels: $$Z_T0.$$ (53) Note that whereas the vanishing of the $`T`$ amplitude is expected from modular invariance in a full string context, it has to be explicitly checked in the particular $`\alpha ^{}0`$ limit we consider. Because of the importance of this result and since we are not aware of any similar computation in the literature, we report in Appendix B a similar computation of gravitational anomalies on the $`T`$ surface for Type IIB string theory in $`D=10`$ and $`D=6`$ on an orbifold. ## 5 Topological interpretation Probably it is interesting to point out that all the anomalies considered so far, eqs.(44), have a nice topological interpretation in terms of the $`G`$-index of the Dirac operator ($`A`$ and $`M`$) and the $`G`$-index of the signature complex ($`K`$), with $`G=𝐙_N`$ for a $`𝐙_N`$ orbifold (see for a nice introduction and more details on the $`G`$-index)<sup>11</sup><sup>11</sup>11A relation between anomalous couplings and the $`𝐙_2`$ signature complex was already exploited in in the case of smooth manifolds.. The $`𝐙_N`$ group can be thought to act on the whole ten dimensional spacetime $`X=R^{1,3}\times T^6`$, as well as on the gauge bundle. As before, we denote by $`g_k=(\theta ^k,\gamma _k)`$ the $`k`$-th element of the complete $`𝐙_N`$ group. Among other things, this will twist the Chern classes appearing in index theorems. The subspace $`X_k`$ left invariant by the geometric $`\theta ^k`$ is $`X_k=_{i=1}^{N_k}R^{1,3}`$, that is $`N_k`$ copies of spacetime. When restricted to $`X_k`$, the tangent bundle of $`X`$ decomposes into the tangent and normal bundles $`𝒯_k`$ and $`𝒩_k`$ of $`X_k`$ in $`X`$. Moreover, the normal bundle $`𝒩_k`$ further decomposes naturally into three components $`𝒩_k^i`$, in which $`\theta `$ acts as $`2\pi v_i`$ rotations. The cotangent and spin bundles, which will be relevant for spinor and self-dual tensor fields, have a similar decomposition. The Dirac-$`G`$ index theorem is then given by (see e.g. ) $$\text{index}(D_{g^k})=_{X^G}\frac{\mathrm{ch}(S_{𝒯_k}^+S_{𝒯_k}^{})\mathrm{ch}_k(S_{𝒩_k}^+S_{𝒩_k}^{})\mathrm{ch}_k(F)}{\mathrm{ch}_k(\stackrel{~}{𝒩}_k)e(𝒯_k)}\mathrm{Td}(𝒯_k^C)$$ (54) where $`S_{𝒯_k}^\pm `$ and $`S_{𝒩_k}^\pm `$ are the positive and negative chirality spin bundles lifted from the tangent and normal bundles, and $`\stackrel{~}{𝒩}_k=_i()^i^i𝒩_k^{}`$ in terms of the conormal bundle $`𝒩_k^{}`$. $`e(𝒯_k)`$ and $`\mathrm{Td}(𝒯_k^C)`$ are the usual Euler and (complexified) Todd classes: $$\mathrm{Td}(𝒯_k^C)=\underset{a=1}{\overset{2}{}}\frac{x_a}{1e^{x_a}}\frac{(x_a)}{1e^{x_a}},e(𝒯_k)=\underset{a=1}{\overset{2}{}}x_a.$$ By expliciting the other terms appearing in (54), one gets $`\mathrm{ch}(S_{𝒯_k}^+S_{𝒯_k}^{})={\displaystyle \underset{a=1}{\overset{2}{}}}(e^{x_a/2}e^{x_a/2}),`$ $`\mathrm{ch}_k(S_{𝒩_k}^+S_{𝒩_k}^{})={\displaystyle \underset{i=1}{\overset{3}{}}}(e^{x_i/2}e^{i\pi kv_i}e^{x_i/2}e^{i\pi kv_i}),`$ $`\mathrm{ch}_k(\stackrel{~}{𝒩}_k)={\displaystyle \underset{i=1}{\overset{3}{}}}(1e^{x_i}e^{2i\pi kv_i})(1e^{x_i}e^{2i\pi kv_i}),`$ (55) where $`x_a`$ and $`x_i`$ are the eigenvalues of the curvature two-form on $`𝒯_k`$ and $`𝒩_k`$. $`\mathrm{ch}_k(F)`$ is precisely the twisted Chern character defined in (31), in terms of the twist matrix $`\gamma _k`$. The trace is in the bifundamental or fundamental representation of the gauge group, for the $`A`$ and $`M`$ surfaces respectively. As previously discussed, the composite field-strength $`G`$ is closely associated to the curvature two-form of the normal bundle $`𝒩_k`$. More precisely, $`x_i=iG_i/\pi `$, and by plugging in the relations (55) above, one gets after some simple algebra $$\text{index}(D_{g^k})=i_{R^{3,1}}C_k\widehat{A}_k(G)\mathrm{ch}_k(F)\widehat{A}(R),$$ (56) which corresponds to the k-th term in the partition functions (44) on $`A`$ and $`M`$. The case of the $`G`$-index of the signature complex can be treated similarly. The $`G`$-signature index theorem is $$\text{index}(𝒟_{g^k}^+)=_{X^G}\frac{\mathrm{ch}(𝒯_k^+𝒯_k^{})\mathrm{ch}_k(𝒩_k^+𝒩_k^{})}{\mathrm{ch}_k(\stackrel{~}{𝒩}_k)e(𝒯_k)}\mathrm{Td}(𝒯_k^C)$$ (57) where $`𝒯_k^\pm ={}_{}{}^{\pm }𝒯_k^{}`$, $`𝒩_k^\pm ={}_{}{}^{\pm }𝒩_k^{}`$, in terms of the cotangent and conormal bundles $`𝒯_k^{}`$ and $`𝒩_k^{}`$. More explicitly, we have $`\mathrm{ch}(𝒯_k^+𝒯_k^{})={\displaystyle \underset{a=1}{\overset{2}{}}}(e^{x_a}e^{x_a}),`$ $`\mathrm{ch}_k(𝒩_k^+𝒩_k^{})={\displaystyle \underset{i=1}{\overset{3}{}}}(e^{x_i}e^{2i\pi kv_i}e^{x_i}e^{2i\pi kv_i}).`$ (58) Similarly to the previous case, the index can then be written as $$\text{index}(𝒟_{g^k}^+)=i_{R^{3,1}}C_{2k}\widehat{L}_k(G)\widehat{L}(R),$$ (59) which corresponds to the k-th term in the partition functions (44) on $`K`$. ## 6 Factorization Having computed all the four amplitudes contributing to the anomaly, we are now in the position of facing the interpretation in terms of quantum anomalies and classical inflows, and understand the mechanism allowing their cancellation. We will also extract all the anomalous couplings to twisted RR fields by factorization. ### 6.1 Quantum anomalies The anomaly arising from open string states is given by the $`A`$ and $`M`$ partition functions: $`𝒜_{open}=𝒜_A+𝒜_M`$. In total, one has: $$𝒜_{open}=\frac{i}{4N}\underset{k=1}{\overset{N1}{}}C_kA_k(G)\left[\mathrm{ch}_k^2(F)\mathrm{ch}_{2k}(2F)\right]\widehat{A}(R).$$ (60) The anomaly from closed string states comes instead from the $`K`$ and $`T`$ partition functions: $`𝒜_{closed}=𝒜_K+𝒜_T`$, where $`𝒜_T`$ denotes all the contributions in (52). It turns out that the $`𝒜_K`$ precisely cancels against the untwisted RR part of $`𝒜_T`$. This reflects the fact that all the descendants of the anti-self-dual 4-form of the original Type IIB theory are projected out by the $`\mathrm{\Omega }`$-projection. One is then left with: $`𝒜_{closed}`$ $`=`$ $`{\displaystyle \frac{i}{2N}}{\displaystyle \underset{k=1}{\overset{N1}{}}}C_k\left[\widehat{A}_k(G)\widehat{G}(R)+\widehat{G}_k(G)\widehat{A}(R)\right]`$ (61) $`+{\displaystyle \frac{1}{2N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \underset{l=1}{\overset{N1}{}}}N_{k,l}\mathrm{ch}(G)\widehat{A}(R).`$ The quantum anomalies (60) and (61) can be qualitatively understood in their alternative interpretation as anomalies involving internal reparametrizations. Indeed, in that context it is easy to discuss the representation of each state under all the symmetries. In particular, all the open string states and the untwisted closed string states transforms under tangent and internal reparametrizations in a way which is dictated essentially by dimensional reduction. This is easily made precise after recalling that the characteristic classes (32), (33) and (34) signal spinor, gravitino and self-dual representations under tangent reparametrizations, and similarly (35), (36) and (37) correspond to spinor, gravitino and self-dual representations under internal reparametrizations. The open string contribution (60) comes clearly from a chiral spinor in $`D=10`$, which once dimensionally reduced to $`D=4`$ gives rise to a multiplet of chiral spinors transforming as an internal spinor. Similarly, the untwisted part (first two terms) of the closed string contribution (61) come from a chiral gravitino in $`D=10`$, which when dimensionally reduced to $`D=4`$ gives rise to a multiplet of chiral gravitinos transforming as an internal spinor (first term), plus a multiplet of chiral spinors transforming as an internal gravitino (second term). Even the canceled contribution of the states projected out by the orientifold projection in summing the $`K`$ and $`T`$ surfaces can be understood. They come, as anticipated, from a self-dual form in $`D=10`$, which is eventually projected out, but would give rise in $`D=4`$ to a multiplet of self-dual forms transforming as an internal self-dual from. The only contribution which cannot be understood in this way is the twisted part (third term) of (61). It is clear that the correponding states must be chiral spinors, and one can argue intuitively that they should transform in a simpler way than untwisted fields under internal reparametrizations (not as tensors), since they arise at given fixed-points in the internal space. Indeed, it is clear from the Chern character in their contribution that they transform with a common $`U(1)`$ charge. The interpretation and analysis of (60) and (61) as sigma-model anomalies is postponed to Section 7. ### 6.2 Classical inflows The GS inflow, which cancels the anomalies computed in previous section, is given by the $`t0`$ limit of the $`A`$, $`M`$ and $`K`$ partition functions (44). By factorization, it is then possible to obtain the anomalous couplings responsible for the inflows. As in the case without composite background , the $`A`$, $`M`$ and $`K`$ partition functions have to factorize exactly. This is made possible by the following non-trivial identities among the characteristic classes defined in Section 4: $`\sqrt{\widehat{A}(R)}\sqrt{\widehat{L}(R/4)}=\widehat{A}(R/2),`$ (62) $`\sqrt{\widehat{A}_{2k}(G)}\sqrt{\widehat{L}_k(G/4)}=\widehat{A}_k(G/2).`$ (63) Indeed, by performing suitable rescalings (allowed by the fact that only the 6-form component of all the polynomials is relevant) and summing the $`k`$-th and the $`Nk`$-th terms in the sums since they correspond to the same closed string twisted sector, the partition functions (44) can be rewritten in the factorized form $`Z_A={\displaystyle \frac{i}{2}}{\displaystyle \underset{k=1}{\overset{(N1)/2}{}}}N_kY_{(k)}Y_{(k)},`$ $`Z_M=i{\displaystyle \underset{k=1}{\overset{(N1)/2}{}}}N_kY_{(2k)}Z_{(2k)},`$ $`Z_K={\displaystyle \frac{i}{2}}{\displaystyle \underset{k=1}{\overset{(N1)/2}{}}}N_kZ_{(2k)}Z_{(2k)},`$ (64) where $`N_k=C_k^2`$ is the number of fixed-points and $`Y_{(k)}={\displaystyle \frac{ϵ_k}{\sqrt{N}}}\sqrt{\left\|{\displaystyle \frac{1}{C_k}}\right\|}\mathrm{ch}_k(ϵ_kF)\sqrt{\widehat{A}_k(ϵ_kG)}\sqrt{\widehat{A}(R)},`$ $`Z_{(2k)}={\displaystyle \frac{4ϵ_k}{\sqrt{N}}}\sqrt{\left\|{\displaystyle \frac{C_{2k}}{C_k^2}}\right\|}\sqrt{\widehat{L}_k(ϵ_{2k}G/4)}\sqrt{\widehat{L}(R/4)}.`$ (65) This implies the following anomalous couplings : $`S_D=\sqrt{2\pi }{\displaystyle \underset{k=1}{\overset{(N1)/2}{}}}{\displaystyle \underset{i_k=1}{\overset{N_k}{}}}{\displaystyle C_{(k)}^{i_k}}Y_{(k)},`$ (66) $`S_F=\sqrt{2\pi }{\displaystyle \underset{k=1}{\overset{(N1)/2}{}}}{\displaystyle \underset{i_k=1}{\overset{N_k}{}}}{\displaystyle C_{(2k)}^{i_k}}Z_{(2k)}.`$ (67) In these couplings, $`C_{(k)}^{i_k}`$ denotes the sum of all the RR forms in the $`k`$-twisted sector and at the fixed-point $`i_k`$; it contains a 4-form plus a 2-form $`\stackrel{~}{\chi }_{(k)}^{i_k}`$ and its dual 0-form $`\chi _{(k)}^{i_k}`$. The relevant components of the charges (65) are therefore the 0, 2 and 4-forms. Thanks to the tadpole condition (30), all irreducible terms in the inflow (64) vanish, and no unphysical negative RR forms propagate in the transverse channel. The total GS couplings can be obtained by summing the D-brane and fixed point contributions (66) and (67), after sending $`k`$ into $`2k`$ in (66). This is allowed for $`N`$ odd, and also in agreement with the fact that in the transverse channel one finds $`k`$-twisted states on $`A`$, and $`2k`$ twisted states on $`M`$ and $`K`$ for the $`k`$-th term in the partition function; in order to add the two consistently one is therefore led to the above substitution. Defining the quantities $`X_{(2k)}=Y_{(2k)}+Z_{(2k)}`$, one has $$S_{GS}=\sqrt{2\pi }\underset{k=1}{\overset{(N1)/2}{}}\underset{i_k=1}{\overset{N_k}{}}C_{(2k)}^{i_k}X_{(2k)}.$$ (68) Using the explicit form (65) of the charges and the tadpole cancellation condition (30), one can check that the total charges $`X_{(2k)}^{(0)}`$ with respect to the RR 4-forms are zero, and the following results for the total charges $`X_{(2k)}^{(2)}`$ and $`X_{(2k)}^{(4)}`$ with respect to the RR 2-forms $`\stackrel{~}{\chi }_{(2k)}^{i_k}`$ and the RR 0-forms $`\chi _{(2k)}^{i_k}`$ are found: $`X_{(2k)}^{(2)}`$ $`=`$ $`{\displaystyle \frac{N_k^{1/4}}{\sqrt{N}(2\pi )}}\left\{i\mathrm{tr}(\gamma _{2k}F)+{\displaystyle \frac{1}{2}}\mathrm{tr}(\gamma _{2k}){\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{tan}(\pi kv_i)G_i\right\},`$ (69) $`X_{(2k)}^{(4)}`$ $`=`$ $`{\displaystyle \frac{ϵ_{2k}N_k^{1/4}}{2\sqrt{N}(2\pi )^2}}\{\mathrm{tr}(\gamma _{2k}F^2){\displaystyle \frac{1}{32}}\mathrm{tr}(\gamma _{2k})\mathrm{tr}R^2+i\mathrm{tr}(\gamma _{2k}F){\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{cot}(2\pi kv_i)G_i`$ $`{\displaystyle \frac{1}{4}}\mathrm{tr}(\gamma _{2k})[{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{tan}^2(\pi kv_i)(G_i)^2`$ $`+\mathrm{\hspace{0.17em}2}{\displaystyle \underset{ij=1}{\overset{3}{}}}{\displaystyle \frac{\mathrm{cos}(2\pi kv_i)\mathrm{cos}(2\pi kv_j)1}{\mathrm{sin}(2\pi kv_i)\mathrm{sin}(2\pi kv_j)}}G_iG_j]\}.`$ Finally, one arrives at a very simple factorized expression for the 6-form encoding the complete sigma-gauge-gravitational anomaly and its opposite inflow: $$𝒜^{(6)}=^{(6)}=i\underset{k=1}{\overset{(N1)/2}{}}N_kX_{(2k)}^{(2)}X_{(2k)}^{(4)}.$$ (71) ## 7 Field theory outlook In this section, we shall address the interpretation of the results found through the string computation within the low-energy supergravity. The 2-form couplings (69) will be responsible for a modification of the kinetic terms of the twisted RR axions, and will force the latter to transform non-homogeneously under gauge and modular transformations. The 4-form couplings will then become anomalous and generate the GS inflow required to cancel all the anomalies. In the following, we focus on $`FFG_i`$ and $`RRG_i`$ anomalies, since these can be compared to field theory expectations. $`FFG_i`$ anomalies This kind of anomalies arise only from open string states. To get an explicit expression from the expansion of (60), it is convenient to transform $`k`$ into $`2k`$ in the annulus contribution. Using (30), one finds for the non-Abelian part: $$𝒜^{FFG_i}=\frac{i}{2N(2\pi )^3}\underset{k=1}{\overset{N1}{}}C_k\mathrm{tan}(\pi kv_i)\mathrm{tr}(\gamma _{2k}F^2)G_i.$$ (72) As we have seen in Section 2, $`FFG_i`$ anomalies are encoded in some coefficients $`b_a^i`$ defined through $$𝒜^{FFG_i}=\frac{i}{2(2\pi )^3}b_a^i\mathrm{tr}(F_a^2)G_i,$$ (73) where the index $`a`$ label the various factors of the gauge group. The coefficients $`b_a^i`$ are found to be in agreement with those computed in for any $`a,i`$ and for both the $`𝐙_3`$ and $`𝐙_7`$ models. This confirms the conjectured anomaly cancellation of mixed $`FFG_i`$ anomalies through a GS mechanism involving RR axions, as proposed in . Indeed, the same anomaly polynomial is reproduced and by factorization the expected couplings are obtained, i.e. the second term in (69). The one-forms $`(X_{(2k)}^{(2)})^{(0)}`$ modify the kinetic terms for the axions $`\chi _{(2k)}^{i_k}`$. Strictly speaking it is only the combination $$\chi _{2k}=\frac{1}{\sqrt{N_k}}\underset{i_k=1}{\overset{N_k}{}}\chi _{(2k)}^{i_k}$$ (74) that gets modified, since all the axions enter in a completely symmetric way in the GS mechanism . Whereas the first term in (69) induces a non-homogeneous $`U(1)`$ transformation for $`\chi _{2k}`$ that eventually leads to a Higgs mechanism through which $`\chi _{2k}`$ itself is eaten by the $`U(1)`$ field, the second term in (69) leads to a non-homogeneous modular transformation for $`\chi _{2k}`$. Note that the WZ descent for $`G_i`$ is $`G_i^{(1)}=i/2[\lambda ^i(t^i)\overline{\lambda }^i(\overline{t}^i)]`$, with $`\lambda ^i(t^i)`$ the lowest component of (10). Correspondingly, the (normalized) kinetic term for $`\chi _{2k}`$ will be invariant under sigma-model transformations if the associated superfields $`M_{2k}`$ transforms, under $`SL(2,R)^i`$, in the following non-homogeneous way: $$M_{2k}M_{2k}\frac{1}{8\pi ^2}\alpha _{2k}^i\lambda ^i(T^i),$$ (75) with $$\alpha _{2k}^i=\frac{(2\pi )^{3/2}}{\sqrt{N}}N_k^{1/4}\mathrm{tr}(\gamma _{2k})\mathrm{tan}(\pi kv_i).$$ (76) $`RRG_i`$ anomalies This kind of anomalies gets contribution both from open and closed string states. Performing the same manipulation as before in the sum over $`k`$ for the annulus contribution, and summing the contributions (60) and (61) from open and closed strings, one finds $`𝒜^{RRG_i}`$ $`=`$ $`{\displaystyle \frac{i}{96N(2\pi )^3}}\{{\displaystyle \underset{k=1}{\overset{N1}{}}}C_k[\mathrm{tan}(\pi kv_i){\displaystyle \frac{1}{2}}\mathrm{cot}(\pi kv_i)]\mathrm{tr}(\gamma _{2k})`$ (77) $`+{\displaystyle \underset{k=1}{\overset{N1}{}}}C_k\mathrm{cot}(\pi kv_i)\left[21+12\left(4\mathrm{sin}^2(\pi kv_i)+{\displaystyle \underset{j=1}{\overset{3}{}}}\mathrm{cos}(2\pi kv_j)\right)\right]`$ $`+{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \underset{l=1}{\overset{N1}{}}}N_{k,l}\}\mathrm{tr}R^2G_i.`$ The first line comes from the open strings, and the second and third line from untwisted and twisted closed strings. As expected the untwisted RNS contribution in the first line encodes the anomaly of the gravitino ($`21`$), the dilatino ($`1`$), and the fermionic partners of the three untwisted moduli ($`2(4\mathrm{sin}^2(\pi kv_i)+_j\mathrm{cos}(2\pi kv_j))`$). The RNS twisted sector contribution in the second line corresponds instead to the anomaly of the neutralini. By explicit evaluation one finds finally: $$𝒜^{RRG_i}=\frac{i}{48(2\pi )^3}\left[\genfrac{}{}{0pt}{}{\mathrm{\hspace{0.17em}10}+\mathrm{\hspace{0.17em}21}+\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}3}\mathrm{\hspace{0.17em}27}}{\mathrm{\hspace{0.17em}6}+\mathrm{\hspace{0.17em}21}+\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}21}}\right]\mathrm{tr}R^2G_i.$$ (78) The coefficient in the square brackets has to be compared with $`b_{grav.}^i=b_{open}^i+b_{closed}^i`$ of , the upper and lower raws corresponding to the $`𝐙_3`$ and $`𝐙_7`$ models respectively. In the notation of , the explicit form of these coefficients is: $`b_{closed}^i=21+\mathrm{\hspace{0.17em}1}+\delta _T^i+{\displaystyle \underset{\alpha }{}}(1+2n_i^\alpha ),`$ $`b_{open}^i=\mathrm{dim}G+{\displaystyle \underset{a=1}{\overset{3}{}}}(1+2n_i^a)\eta _a.`$ (79) In $`b_{closed}^i`$, $`\delta _T^i`$ is the total contribution of the untwisted moduli (nine for the $`𝐙_3`$ model and three for the $`𝐙_7`$ model) and $`\alpha `$ runs over all the twisted massless states. These are assumed to have modular weight $`n_i^\alpha `$ as defined in Section 2, and in it was assumed that $`n_i^\alpha =0`$. In $`b_{open}^i`$, the first term is the contribution of the gaugini, where $`G`$ is the total gauge group of the model, $`n_a^i=\delta _a^i`$ are the modular weights of the charged fields $`C^a`$ and $`\eta _a`$ simply counts the number of charged states belonging to the group $`a`$. Comparing the string result (78) with the field theory expectations given by (79), one finds agreement for $`b_{open}^i`$ (first number in (78)) and for the untwisted contribution in $`b_{closed}^i`$ (next three numbers), but opposite signs for the twisted contribution (last number). Assuming the validity of (79), agreement with the string results would predict twisted modular weights $`n_i^\alpha =1`$, $`i,\alpha `$. This is in apparent contradiction with the non-homogeneous transformation (75) required for the cancellation of sigma-gauge anomalies. The sign that we find for the contribution of twisted modulini is crucial for the realization of the GS anomaly cancellation mechanism, since it directly influences the factorizability of the quantum anomaly. We do not have a full understanding of this discrepancy; rather, we would like to revisit the assumptions at the origin of the above field theory analysis and point out a few delicate points. A first point to observe, in comparing string results with field theory expectations in $`D=4`$ $`N=1`$ models, is that the first are believed to be expressed in terms of linear multiplets, whereas the latter are often given in terms of the usual chiral multiplets, as is the case for (79). The two multiplets are related by the so called linear multiplet - chiral multiplet duality, that is basically the extension to superfields of the duality between a two-form and a scalar in four dimensions. It is also known that the GS terms modify the above duality . Correspondingly, particular attention has to be paid in comparing the results (78) with field theory formulae obtained using the chiral multiplet basis as (79) (see for instance footnote 7). A second very important point is that the expression (79) for the anomaly coefficients are valid only under the assumption that the Kähler potential $`K^{(M)}`$ for twisted fields and their modular transformations have the form (11). Unfortunately, the potential $`K^{(M)}`$ has not been computed yet in Type IIB orientifold models, and therefore it is not possible to verify directly these assumptions. We propose that the Kähler potential for twisted fields does in fact not satisfy the assumptions at the origin of (79), so that the whole field-theory derivation of sigma-model anomalies, as reviewed in Section 2 and expressed in (79), has to be revisited . A first possibility is that $`K^{(M)}(M+\overline{M})^2`$, as proposed in . This potential satisfies the assumptions behind (79) (and leads to $`n_i^\alpha =0`$ as assumed in ), but only if one neglects the correction induced by the GS couplings (69) and (6.2). These are indeed present, as described in , and it might be that they must be considered on equal footing with the rest of the potential<sup>12</sup><sup>12</sup>12This seems quite strange from a string theory point of view, but we believe it might be reasonable in light of the string coupling dependence of the definition (19) for the $`T^i`$ moduli.. A second possibility is that $`K^{(M)}`$ is a different function of the twisted moduli, invariant under sigma-model transformations and the shift (75), whose form does not satisfy the assumptions leading to (79). An explicit string computation of $`K^{(M)}`$ would therefore be extremely interesting and could give a definite answer to the problems raised above. Unfortunately, such a computation appears to be quite complicated. ## 8 Conclusions In this paper, we have studied along the lines of the pattern of sigma-gauge-gravitational anomaly cancellation in compact Type IIB $`D=4`$ $`N=1`$ $`𝐙_N`$ orientifolds with $`N`$ odd. Our main result is that all the anomalies are cancelled through a generalized GS mechanism. The starting point of our analysis is the definition of the effective vertex operator corresponding to the sigma-model connection. We provided several general arguments for identifying it with the vertex encoding Kähler deformations of the orbifold, but we were able to give only a not completely rigorous derivation which cannot be taken as a proof. A posteriori, this identification is strongly supported also by the results obtained for the anomalies using this vertex. Under the assumption that the effective vertex is indeed correct, we generalize the known results for gauge-gravitational anomalies and show that all possible sigma-gauge-gravitational anomalies are cancelled through a GS mechanism. This is essentially what was proposed in for sigma-gauge anomalies, and seems to evade the arguments of against a field theory mechanism for the cancellation of sigma-gravitational anomalies. We interpret this discrepancy as evidence that the comparison of the string results with the field theory expectations is probably more subtle than expected. In particular, we propose that the actual Kähler potential for twisted fields does not satisfy the usual assumptions made in the literature, so that the interpretation of our string results remains actually open. We would like to stress that the present results imply a full cancellation of anomalies in all possible channels. The torus contribution presents a surprising cancellation and yields vanishing anomalies and inflows. This implies in particular that the dilaton field does not play any role in the GS mechanism. The annulus, Möbius strip and Klein bottle contributions are instead topological, guaranteeing an exact cancellation between quantum anomalies and classical inflows mediated by twisted RR axions. ###### Acknowledgments. It is a pleasure to thank L.E. Ibáñez and M. Klein for continuous support and collaboration. We would like also to thank C.P. Bachas, G.L. Cardoso, J.-P. Derendinger, J. Fröhlich and H. Nilles for interesting discussions and useful comments. We acknowledge the Universidad Autonoma de Madrid for hospitality. C.A.S also thanks the Edwin Schrödinger Institute of Vienna and the Institute for Theoretical Physics of the ETH in Zürich. This work has been supported by the EC under TMR contract ERBFMRX-CT96-0045 and by the Fundamenteel Onderzoek der Materie (FOM). ## Appendix A $`\vartheta `$-functions For convenience, we introduce here a convenient notation for the twisted $`\theta `$-functions appearing in orbifold and orientifold partition functions. In particular, in order to keep manifest the origin of each of these, we shall define $`\theta _1\left[{\displaystyle \genfrac{}{}{0pt}{}{\alpha }{\beta }}\right](z\|\tau )=\theta \left[{\displaystyle \genfrac{}{}{0pt}{}{\frac{1}{2}+\alpha }{\frac{1}{2}+\beta }}\right](z\|\tau ),`$ (80) $`\theta _2\left[{\displaystyle \genfrac{}{}{0pt}{}{\alpha }{\beta }}\right](z\|\tau )=\theta \left[{\displaystyle \genfrac{}{}{0pt}{}{\frac{1}{2}+\alpha }{0+\beta }}\right](z\|\tau ),`$ (81) $`\theta _3\left[{\displaystyle \genfrac{}{}{0pt}{}{\alpha }{\beta }}\right](z\|\tau )=\theta \left[{\displaystyle \genfrac{}{}{0pt}{}{0+\alpha }{0+\beta }}\right](z\|\tau ),`$ (82) $`\theta _4\left[{\displaystyle \genfrac{}{}{0pt}{}{\alpha }{\beta }}\right](z\|\tau )=\theta \left[{\displaystyle \genfrac{}{}{0pt}{}{0+\alpha }{\frac{1}{2}+\beta }}\right](z\|\tau ),`$ (83) in terms of the usual twisted $`\theta `$-functions $`\theta \left[{\displaystyle \genfrac{}{}{0pt}{}{\alpha }{\beta }}\right](z\|\tau )={\displaystyle \underset{n}{}}q^{\frac{1}{2}(n\alpha )^2}e^{2\pi i(z\beta )(n\alpha )}.`$ (84) All the properties and identities relevant to the usual $`\theta `$-functions (84) easily translate into analogous properties of (80)-(83). ## Appendix B Anomalies in Type IIB string theory The cancellation of gravitational anomalies in Type IIB supergravity theories requires non-trivial identities involving the anomalies of dilatinos, gravitinos and self-dual forms. These are given by $$I_{1/2}=\widehat{A}(R),I_{3/2}=\widehat{G}(R),I_A=\frac{1}{8}\widehat{L}(R),$$ (85) in terms of the characteristic classes (32)-(34). From a Type IIB string theory point of view, anomaly freedom is more manifest since the corresponding torus amplitude is perfectly finite. However, it is clear that in the low-energy field theory limit one has to reproduce in string theory the same non-trivial identity. This can be regarded as a technique to compute anomalous Feynman diagrams using a string regularization. Due to the relevance of the torus amplitude in the mixed sigma-gauge-gravitational anomalies considered in this paper, we find useful to report here some details on how to reproduce in Type IIB string theory the aforementioned identity. As explained in Section 3 and 4, the only potentially anomalous contributions on the torus come from the three odd-even and the three even-odd spin-structures, and the total anomaly is given by the expression (45), in terms of the partition functions (46) defined through the deformation vertices (48) and (49). It turns out that the two partition functions (46) will always be modular invariant, so that the $`_0`$ component of the boundary gives a vanishing contribution. Moreover, on the other component $`_{\mathrm{}}`$ of the boundary, the odd-even and even-odd partition functions become equal and sum. In the following, we will therefore restrict to the odd-even spin-structures. $`D=10`$ In the ten dimensional case, the partition functions (46) are particularly easy to compute. One gets $$Z_T^{S_\alpha }=\frac{1}{4}\underset{a=1}{\overset{5}{}}\left[\frac{ix_a}{\theta _1(ix_a/\pi \|\tau )}\theta _\alpha (ix_a/\pi \|\tau )\right]\frac{\eta ^3(\tau )}{\theta _\alpha (0\|\tau )}.$$ (86) Here $`\alpha =2,3,4`$ represent respectively the RR, $`\mathrm{RNS}_+`$ and $`\mathrm{RNS}_{}`$ spin-structures, the factor of $`1/4`$ is due to the left and right GSO projections and $`x_a=\lambda _a/2\pi `$, in terms of the skew eigenvalues $`\lambda _a`$ of the gravitational curvature $`R`$. The first fraction is the contribution of the bosonic and fermionic fields, whereas the last fraction is due to ghosts and superghosts. Taking the limit $`\tau _2\mathrm{}`$, one obtains in the RR spin-structure $$Z_T^{RR}\frac{1}{8}\underset{a=1}{\overset{5}{}}\frac{x_a}{\mathrm{tanh}x_a}.$$ (87) In the $`\mathrm{RNS}_\pm `$ spin-structures, similarly $$Z_T^{RNS}=Z_T^{RNS_+}Z_T^{RNS_{}}\underset{a=1}{\overset{5}{}}\frac{x_a/2}{\mathrm{sinh}x_a/2}\left(2\underset{b=1}{\overset{5}{}}\mathrm{cosh}x_b2\right),$$ (88) where we rescaled by a factor of 2 the $`x_a`$’s, exploiting the fact that only the 12-form of (88) is relevant. Notice that the leading “tachyonic” terms in $`Z_T^{RNS_\pm }`$ cancel in the combination $`Z_T^{RNS_+}Z_T^{RNS_{}}`$. By summing the three contributions one finds as expected the anomaly of an anti-chiral gravitino and of a chiral dilatino from the RNS/NSR sector and that of an (anti)self-dual tensor from the RR sector. In total, one gets $$I_T=I_{3/2}+I_{1/2}I_A=0,$$ (89) ensuring the absence of pure gravitational anomalies in $`D=10`$ Type IIB supergravity and superstring theory . $`D=6`$ on $`T^4/𝐙_N`$ As usual in orbifold theories, the partition functions (46) contain a sum over orbifold twisted sectors $`l`$, as well as a projection on $`𝐙_N`$-invariant states; see (47). In the following, we will further distinguish between the contributions coming from untwisted and twisted sectors. The twist vector is $`v_i=(1/N,1/N)`$, $`C_k=_i(2\mathrm{sin}(\pi kv_i))`$, and $`N_{k,l}`$ are the number of points that are at the same time $`k`$ and $`l`$-fixed. The total partition function is $$Z_T=\underset{\alpha }{}()^\alpha \underset{l=0}{\overset{N1}{}}Z_T^{S_\alpha (l)},$$ (90) where $$Z_T^{S_\alpha (l)}=\frac{1}{4N}\underset{k=0}{\overset{N1}{}}N_{k,l}\underset{i=1}{\overset{2}{}}\frac{\theta _\alpha \left[\genfrac{}{}{0pt}{}{lv_i}{kv_i}\right](0\|\tau )}{\theta _1\left[\genfrac{}{}{0pt}{}{lv_i}{kv_i}\right](0\|\tau )}\underset{a=1}{\overset{3}{}}\left[\frac{ix_a}{\theta _1(ix_a/\pi \|\tau )}\theta _\alpha (ix_a/\pi \|\tau )\right]\frac{\eta ^3(\tau )}{\theta _\alpha (0\|\tau )}.$$ (91) In the $`\tau _2\mathrm{}`$ limit, one finds in the RR spin-structures: $`Z_T^{RR(0)}{\displaystyle \frac{1}{8N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}C_{2k}{\displaystyle \underset{a=1}{\overset{3}{}}}{\displaystyle \frac{x_a}{\mathrm{tanh}x_a}},`$ $`Z_T^{RR(l0)}{\displaystyle \frac{1}{8N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}N_{k,l}{\displaystyle \underset{a=1}{\overset{3}{}}}{\displaystyle \frac{x_a}{\mathrm{tanh}x_a}}.`$ (92) One can easily check that for any $`N=2,3,4,6`$, the total is given by $$Z_T^{RR}2\underset{a=1}{\overset{3}{}}\frac{x_a}{\mathrm{tanh}x_a}.$$ (93) In the $`\mathrm{RNS}_\pm `$ spin-structures one has to pay particular attention in taking the limit, because when $`l=N/2`$, the fields in the internal directions have zero modes. One finds the following results: $`Z_T^{RNS_\pm (0)}\pm {\displaystyle \frac{1}{2N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}C_k{\displaystyle \underset{a=1}{\overset{3}{}}}{\displaystyle \frac{x_a/2}{\mathrm{sinh}x_a/2}}\left(2{\displaystyle \underset{b=1}{\overset{3}{}}}\mathrm{cosh}x_b2+{\displaystyle \underset{i=1}{\overset{2}{}}}(2\mathrm{cos}2\pi kv_i)\right),`$ $`Z_T^{RNS_\pm (l0,N/2)}\pm {\displaystyle \frac{1}{N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}N_{k,l}{\displaystyle \underset{a=1}{\overset{3}{}}}{\displaystyle \frac{x_a/2}{\mathrm{sinh}x_a/2}},`$ $`Z_T^{RNS_+(N/2)}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}N_{k,N/2}{\displaystyle \underset{a=1}{\overset{3}{}}}{\displaystyle \frac{x_a/2}{\mathrm{sinh}x_a/2}}{\displaystyle \underset{i=1,2}{}}(2\mathrm{cos}\pi kv_i)^2`$ $`Z_T^{RNS_{}(N/2)}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}N_{k,N/2}{\displaystyle \underset{a=1}{\overset{3}{}}}{\displaystyle \frac{x_a/2}{\mathrm{sinh}x_a/2}}{\displaystyle \underset{i=1,2}{}}(2\mathrm{sin}\pi kv_i)^2.`$ (94) We omitted the leading “tachyonic” term that, as in the previous case, will cancel in taking the sum $`Z_T^{RNS_+}Z_T^{RNS_{}}`$. One can easily verify that the total result in the RNS sectors, obtained by summing over the two $`\mathrm{RNS}_\pm `$ contributions and over all twisted and untwisted sectors, is the same for any $`N=2,3,4,6`$ and given by $$Z_T^{RNS}=\mathrm{\hspace{0.17em}2}\underset{a=1}{\overset{3}{}}\frac{x_a/2}{\mathrm{sinh}x_a/2}\left(2\underset{b=1}{\overset{3}{}}\mathrm{cosh}x_b22\right)$$ (95) Putting all together, one gets finally $$I_T=2(I_{3/2}21I_{1/2}8I_A)=0$$ (96) ensuring the absence of purely gravitational anomalies in Type IIB theory on $`T^4/𝐙_N`$. ## Appendix C Vanishing of the torus amplitude We show here that the whole 6-form component of the torus amplitude, including $`G^3`$ anomalies, vanishes. For the $`RRG_i`$ terms, one gets $`Z_T^{R^2G}`$ $`=`$ $`{\displaystyle \frac{i}{96N(2\pi )^3}}{\displaystyle \underset{i=1}{\overset{3}{}}}\{\mathrm{\hspace{0.17em}4}{\displaystyle \underset{k=1}{\overset{N1}{}}}C_{2k}\mathrm{sin}^1(2\pi kv_i)`$ (97) $`+{\displaystyle \underset{k=1}{\overset{N1}{}}}C_k\mathrm{cot}(\pi kv_i)\left[21+12\left(4\mathrm{sin}^2(\pi kv_i)+{\displaystyle \underset{j=1}{\overset{3}{}}}\mathrm{cos}(2\pi kv_j)\right)\right]`$ $`+{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \underset{l=1}{\overset{N1}{}}}N_{k,l}\}\mathrm{tr}R^2G_i`$ $`=`$ $`{\displaystyle \frac{i}{48(2\pi )^3}}\left[{\displaystyle \genfrac{}{}{0pt}{}{8+\mathrm{\hspace{0.17em}21}+\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}3}\mathrm{\hspace{0.17em}27}}{0+\mathrm{\hspace{0.17em}21}+\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}21}}}\right]\mathrm{tr}R^2\left({\displaystyle \underset{i=1}{\overset{3}{}}}G_i\right)`$ $`=`$ $`0,`$ where we reported in square bracket the explicit values for both the $`𝐙_3`$ (up) and $`𝐙_7`$ (down) orientifolds. Consider next the $`G_iG_jG_p`$ terms. The RR twisted contributions vanish as before, whereas the RNS twisted ones are present and can be easily read from the last line of (52). On the contrary, the untwisted RR and RNS contributions requires more work. However, one can now put to zero the gravitational curvature. By doing so, the contribution of the superghosts cancels that of one of the two complex spacetime fermions in (50) and one can therefore use the Riemann identity to simplify the result. In the $`\tau _2\mathrm{}`$ limit, one has then: $`{\displaystyle \frac{8_i\mathrm{cos}(\pi kv_i+G_i/2\pi )2_i\mathrm{cos}2(\pi kv_i+G_i/2\pi )2}{8_i\mathrm{sin}(\pi kv_i+G_i/2\pi )}}`$ $`=2\mathrm{sin}\left({\displaystyle \underset{p=1}{\overset{3}{}}}G_p/4\pi \right){\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\mathrm{sin}\left[(\pi kv_i+G_i/2\pi )(_pG_p/4\pi )\right]}{\mathrm{sin}(\pi kv_i+G_i/2\pi )}}.`$ The relevant cubic term of the partition function are now easily computed, and one finds: $`Z_T^{G^3}`$ $`=`$ $`{\displaystyle \frac{i}{24N(2\pi )^3}}\{{\displaystyle \underset{k=1}{\overset{N1}{}}}N_k[(3{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \underset{ji=1}{\overset{3}{}}}\mathrm{cot}(\pi kv_i)5)\left({\displaystyle \underset{p=1}{\overset{3}{}}}G_p\right)^3`$ (98) $`+\mathrm{\hspace{0.17em}6}{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{sin}^2(\pi kv_i)G_i\left({\displaystyle \underset{p=1}{\overset{3}{}}}G_p\right)^2]`$ $`\mathrm{\hspace{0.17em}2}{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \underset{l=1}{\overset{N1}{}}}N_{k,l}\left({\displaystyle \underset{p=1}{\overset{3}{}}}G_p\right)^3\}`$ $`=`$ $`{\displaystyle \frac{i}{4N(2\pi )^3}}\left[{\displaystyle \genfrac{}{}{0pt}{}{9\mathrm{\hspace{0.17em}15}+\mathrm{\hspace{0.17em}24}\mathrm{\hspace{0.17em}18}}{3\mathrm{\hspace{0.17em}5}+\mathrm{\hspace{0.17em}16}\mathrm{\hspace{0.17em}14}}}\right]\left({\displaystyle \underset{p=1}{\overset{3}{}}}G_p\right)^3`$ $`=`$ $`0.`$
warning/0006/cond-mat0006223.html	ar5iv	text	# The nondeterministic Nagel-Schreckenberg traffic model with open boundary conditions ## I Introduction Asymmetric exclusion processes (ASEP) play an important role in non-equilibrium statistical mechanics. The one-dimensional ASEP is a lattice model which describes particles hopping in one direction with stochastic dynamics and hard core exclusion. It was first introduced in 1968 to provide a qualitative understanding of the kinetics of the protein synthesis on RNA templates . It turned out, however, that \- despite its simplicity - there are numerous further applications of the ASEP on the field of interface growth, polymer dynamics, and traffic flow -. Unfortunately, as far as traffic is concerned, the ASEP yields rather unrealistic results, because essential phenomena like acceleration or slowing down cannot be reproduced in this model. As a consequence, Nagel and Schreckenberg developed an extension of the ASEP resulting in a one-dimensional probabilistic cellular automaton model . According to the Nagel-Schreckenberg model the road consists of a single lane which is divided into L cells of equal size numbered by i = 1, 2, $`\mathrm{}`$, L and the time is also discrete. Each site can be either empty or occupied by a car with integer velocity v = 0, 1, $`\mathrm{}`$, $`\text{v}_{max}`$. All sites are simultaneously updated according to four successive steps: 1. Acceleration: increase v by 1 if v $`<`$ $`\text{v}_{max}`$. (1) 2. Slowing down: decrease v to v = d if necessary (d: number of empty cells in front of the car). (2) 3. Randomization: decrease v by 1 with randomization probability p if p $`>`$ 0. (3) 4. Movement: move car v sites forward. (4) Either ring (periodic boundary conditions) or open (open boundary conditions) geometry is considered. In the case of ring geometry cars move on a ring and the car density in the system keeps constant. Open systems, on the other hand, are characterized by the injection (extinction) rate $`\alpha `$ ($`\beta `$), that means, by the probability $`\alpha `$ ($`\beta `$) that a car moves into (out of) the system. For the maximum velocity $`\text{v}_{max}`$ = 1 the model is identical with the ASEP with parallel update - which has been solved exactly with periodic boundary conditions and recently also with open boundary conditions , . In this special case three regimes (free flow, jamming, and maximum current) can be distinguished from each other. The transition from the free flow to the jamming phase at the $`\alpha `$ = $`\beta `$ \- line for $`\alpha `$, $`\beta `$ $`<`$ 1-$`\sqrt{\text{p}}`$ is of first order. The transition from the free flow (jamming) to the maximum current phase is continuous and takes place at the injection (extinction) rate $`\alpha _c`$ ($`\beta _c`$) with $`\alpha _c`$(p) = $`\beta _c`$(p) = 1-$`\sqrt{\text{p}}`$. As it is common for traffic simulations cars are updated in parallel in the Nagel-Schreckenberg model, too, because this update scheme is the only one which models the formation of spontaneous jams occurring in real traffic . Systems with parallel update are furthermore characterized by strong short-range correlations, and therefore, short-range correlation functions play an important role here , . Most of the work dealing with the Nagel-Schreckenberg model for $`\text{v}_{max}`$ $`>`$ 1 impose periodic boundary conditions -. Much attention has been paid to the question of the transition from the free flow to the jamming regime. According to the state-of-the-art this is a crossover rather than a sharp transition \- . Systems with periodic boundary conditions are furthermore characterized by a trivial density profile $`\rho `$(i) = $`\rho `$ with 1 $``$ i $``$ L due to translational invariance. In this context it should be mentioned that short-range correlation functions are well-suited for the description of the free flow - jamming transition : The free flow regime is characterized by anticorrelations around a propagating peak, that is, in free flow cars are surrounded by empty space. At the critical density $`\rho _c`$ the anticorrelations are maximally developed, and for higher densities they vanish. Simultaneously, a jamming peak develops according to the fact that the back car is strongly slowed down in a jam. In the following, systems for maximum velocities $`\text{v}_{max}`$ $``$ 10 are investigated for the more realistic case of open boundaries. Boundary conditions are defined as in : At site i = 0, that means out of the system a vehicle with the probability $`\alpha `$ and with the velocity v = $`\text{v}_{max}`$ is created. This car immediately moves according to the Nagel-Schreckenberg rules. If the velocity of the injected car on i = 0 is v = 0 (because site i = 1 is occupied by another car or because the front car is on site i = 2 and the injected car is slowed down by 1 due to randomization) then the injected car is deleted. At i = L+1 a ”block” occurs with probability 1 - $`\beta `$ and causes a slowing down of the cars at the end of the system. Otherwise, with probability $`\beta `$, the cars simply move out of the system. In systems with open boundaries have been already analyzed for the randomization probability p = 0, i.e., by ignoring the randomization step. The most interesting feature of the deterministic Nagel-Schreckenberg model for maximum velocities $`\text{v}_{max}`$ $``$ 3 and open boundaries is the existence of so-called buffers: As a consequence of the parallel updating and the hindrance an injected car feels from the front car at the beginning of the system spaces larger than $`\text{v}_{max}`$ develop between two neighbouring cars for high injection rates. This can be easily demonstrated by considering i = 0 and the first sites of the system i = 1, $`\mathrm{}`$, 5 for $`\alpha `$ = $`\beta `$ = 1 and $`\text{v}_{max}3`$. At an arbitrary time t = $`\text{t}_0`$ a car with velocity $`\text{v}_{max}`$ is injected on i = 0, i.e. $`injection:\text{v}_{max}\mathrm{.\hspace{0.17em}\; .\hspace{0.17em}\; 2}..`$ (5) After application of the NaSch rules on the system we have injected on i = 0, i.e. $`movement:.\mathrm{.\hspace{0.17em}\; 2}...`$ (6) Correspondingly, we get for t = $`\text{t}_0+1`$ $`injection:\text{v}_{max}\mathrm{.\hspace{0.17em}\; 2}...`$ (7) $`movement:.\mathrm{1\hspace{0.17em}\; .\hspace{0.17em}\; .\hspace{0.17em}\; .\hspace{0.17em}\; 3}`$ (8) and for t = $`\text{t}_0+2`$ $`injection:\text{v}_{max}\mathrm{1\hspace{0.17em}\; .\hspace{0.17em}\; .\hspace{0.17em}\; .\hspace{0.17em}\; 3}`$ (9) $`movement:0\mathrm{.\hspace{0.17em}\; .\hspace{0.17em}\; 2}..`$ (10) As the car on site i = 0 cannot move it is deleted and the situation starts over in the next time step: $`injection:\text{v}_{max}\mathrm{.\hspace{0.17em}\; .\hspace{0.17em}\; 2}..`$ (11) $`movement:.\mathrm{.\hspace{0.17em}\; 2}...`$ (12) and so on. Obviously, one car is lost out of three injection possibilities. In the system, far from the boundaries, the distance between two neighbouring cars turns out to be alternately $`\text{d}_1=\text{v}_{max}`$ and $`\text{d}_2=2(\text{v}_{max}1)`$. That means, in addition to the expected $`\text{v}_{max}`$ empty sites larger gaps occur in the $`\alpha `$ $``$ 1, $`\beta `$ $``$ 1 limit. We call these additional sites ”buffers” because they have a buffer effect at the end of the system: Due to these sites the development of jamming waves is suppressed even for $`\beta `$ $`<`$ 1. The transition from the free flow to the jamming phase is of first order and accompanied by the collapse of the buffers. The effect resulting from the buffers do not depend on the maximum velocity if $`\text{v}_{max}`$ $``$ 5 (for $`\text{v}_{max}`$ = 3,4 the buffer effect is not so strong as the buffers are not completely developed for that case). However, randomization is indispensable for the analysis of real traffic as it takes human behaviour into account: The behaviour of a car driver is not like that of a machine but rather contains unpredictable elements. In traffic, over-reactions when slowing down can be found as well as delays when accelerating, furthermore fluctuations when following a car (follow-the-leader situation) and so on. Besides this motivation it is of interest to compare the general p $`>`$ 0, $`\text{v}_{max}`$ $``$ 1 case with the previously investigated models (p $`>`$ 0, $`\text{v}_{max}`$ = 1 , and p = 0, $`\text{v}_{max}`$ $``$ 1 ). The presented results were obtained by simulating a L = 1024 sample with at least 1000 runs with $`10^4`$ time steps each. In order to investigate the influence of randomization on the system we proceed similarly to : Section II considers the behaviour of current and average occupation number in the middle of the system. Section III deals with density profiles, Section IV with short-range correlation functions. Finally, the results are summarized in Section V. ## II Current and Average Occupation number at the middle of the system The phase diagram for systems with probability p = 0.5 and maximum velocity $`\text{v}_{max}`$ = 2 and $`\text{v}_{max}`$ $``$ 5 is shown in Figs 1a,b (the case $`\text{v}_{max}`$ = 3,4 is similar to the latter case). For $`\text{v}_{max}`$ = 2 the phase diagram (Fig 1a) is qualitatively the same as for the case $`\text{v}_{max}`$ = 1 , : The free flow and the jamming regime are divided by a straight line and for $`\alpha `$ $`>`$ 0.35 and $`\beta `$ $`>`$ 0.8 the system is in the maximum current phase. Obviously, the maximum current regime is smaller than for $`\text{v}_{max}`$ = 1 and there is no symmetry along the $`\alpha `$=$`\beta `$-line. For $`\text{v}_{max}`$ = 3 the maximum current regime is even smaller ( for $`\alpha `$ $`>`$ 0.35 and $`\beta `$ $`>`$ 0.85 ) and the free flow/jamming border shows a slight bending. These tendencies are even stronger developed for higher maximum velocities $`\text{v}_{max}`$ with a maximum current regime for for $`\alpha `$ $`>`$ 0.35, $`\beta `$ $`>`$ 0.87, and $`\text{v}_{max}`$ = 4 (for $`\alpha `$ $`>`$ 0.35 and $`\beta `$ $`>`$ 0.89, and $`\text{v}_{max}`$ $``$ 5, see Fig 1b). Another interesting feature of the nondeterministic case is that the course of the free flow/jamming border is totally different from that for p = 0. In Fig 1c it can be clearly seen that this difference is a consequence of the vanishing buffer effect due to randomization. For randomization probabilities p $`>`$ $`\text{p}_c`$ ($`\text{p}_c`$ = 0.1172 $`\pm `$ 0.008 for $`\text{v}_{max}`$ = 5) there is no sign of the buffer effect any more, and a (rectangular) maximum current regime develops for $`\alpha `$ $`>`$ $`\alpha _{FreeFlow}`$ and $`\beta `$ $`>`$ $`\beta _{Jamming}`$. In order to understand the nature of the transition between the phases we consider the average occupation number on the site i = $`\frac{\text{L}}{\text{2}}`$, $`\rho `$ (i = $`\frac{\text{L}}{\text{2}}`$), as proposed in . Fig 2 shows $`\rho `$ (i = $`\frac{\text{L}}{\text{2}}`$) for $`\text{v}_{max}`$ = 5 as a function of the injection and the extinction rates (the average occupation number on i = $`\frac{\text{L}}{\text{2}}`$ for any $`\text{v}_{max}`$ $`>`$ 1 behaves similarly). It turns out that phase transitions in systems with maximum velocity $`\text{v}_{max}`$ $`>`$ 1 show the following features: The transition from free flow to jamming is of first order indicated by a jump in $`\rho `$ (i = $`\frac{\text{L}}{\text{2}}`$). At the maximum current/free flow and the maximum current/jamming transition there is a jump in the derivative of $`\rho `$ (i = $`\frac{\text{L}}{\text{2}}`$) what is a hint at a continuous phase transition. These conclusions will be confirmed in the next section which deals with the investigation of the corresponding density profiles. In the following we analyze the current due to the influence of the \- boundaries (13) \- maximum velocity $`\text{v}_{max}`$ (14) \- randomization probability p (15) The best way to investigate the influence of the left boundary is to consider the case $`\beta `$ = 1 for p = 0.5 where cars simply move out of the system. In Fig 3a free flow and maximum current phase can be clearly distinguished from each other. For maximum velocities $`\text{v}_{max}`$ $``$ 5 the curves are nearly the same with a maximum at $`\alpha `$ $``$ 0.35 becoming stronger with increasing $`\text{v}_{max}`$ (the occurrence of the maximum will be explained below). The dependence of the current on $`\alpha `$ and $`\beta `$ for $`\text{v}_{max}`$ = 2, 3, 4 is qualitatively the same as for $`\text{v}_{max}`$ = 1. For the investigation of the influence of the right boundary we consider the case $`\alpha `$ = 1 for p = 0.5 (Fig 3b). Here, the current for $`\text{v}_{max}`$ $``$ 5 does not depend on the maximum velocity. Furthermore, it seems to increase monotonously with increasing $`\beta `$ also in the maximum current phase ($`\beta `$ $`>`$ 0.89). Investigations for system sizes L $``$ 4096, however, show that the latter observation is just a finite size effect and that the current for $`\text{v}_{max}`$ $``$ 5 and $`\beta `$ $`>`$ 0.89 is constant. Apart from this, the curves for the current do not change in an essential way with increasing system size L, and therefore it is sufficient to investigate systems with L = 1024 in the following. From the observations so far we can conclude that - as for the deterministic case - the behaviour of the system only negligibly changes when maximum velocities $`\text{v}_{max}`$ $``$ 5 are considered. Therefore we restrict ourselves to the case $`\text{v}_{max}`$ = 5 in the following observations. We will now investigate the influence of the randomization probability on the behaviour of the system. It can be seen from Figs 4a,b that the buffer effect observed for the deterministic case vanishes with increasing p: For $`\beta `$ = 1 in Fig 4a the maximum at $`\alpha `$ $``$ 0.81 resulting from the existence of the buffers moves to the left and becomes weaker and weaker (In Fig 4a we make an exception and consider the current at $`\text{v}_{max}`$ = 10 instead of $`\text{v}_{max}`$ = 5, because for the latter case this effect is nearly invisible). For the injection rate $`\alpha `$ = 1, on the other hand, the buffer effect vanishes as soon as randomization probabilities p $`>`$ $`\text{p}_c`$ are considered. To sum it up it can be said that as a consequence of the buffer effect, the course of the current in the maximum current phase deviates from the expected (constant) behaviour for $`\beta `$ = 1, p = 0.5, and maximum velocities $`\text{v}_{max}`$ $``$ 5 showing a slight maximum at $`\alpha `$ $``$ 0.35. Besides, there are strong indications that a continuous transition from the free flow (jamming) to a maximum current phase develops with increasing randomization probability on the $`\beta `$ = 1 - line ($`\alpha `$ = 1 - line). More convincing arguments for the existence of a maximum current phase, however, will be given in the following sections. ## III Density Profiles Our observations so far consider the behaviour of the whole system and of the site i = $`\frac{\text{L}}{\text{2}}`$. For the analysis of what happens on the other sites it is useful to investigate the density profiles. Of special interest in this context is the question in how far the density profiles reflect the transition between the phases. For that purpose we consider the density profiles for $`\beta `$ = 1 ($``$ transition from free flow to maximum current), for $`\beta `$ = 0.7 ($``$ transition from free flow to jamming), and for $`\alpha `$ = 1 ($``$ transition from maximum current to jamming). It turns out that – as in the case of $`\text{v}_{max}`$ = 1 – the free flow (jamming) phase can be divided into the regime AI and AII (BI and BII). The following investigations are confined to the randomization probability p = 0.5, L = 1024 and $`\text{v}_{max}`$ = 5 (for $`\text{v}_{max}`$ $``$ 5 and L $`>`$ 1024 the density profiles are qualitatively the same). In Fig 5a the transition from free flow to maximum current for $`\beta `$ = 1 can be clearly seen. The free flow regime is characterized by oscillations at the beginning of the system dying out for i $``$ 100 (if p = 0.5 and $`\text{v}_{max}`$ = 5) due to randomization and the density profile becomes constant. It can therefore be said that randomization blurs the influence of the left boundary. In the maximum current regime we do not have any oscillations at all. Instead, an analytic decrease of the density is observed becoming stronger with increasing $`\alpha `$. This phenomenon can be easily understood as cars hinder each other at the beginning of the system for high injection rates: The higher the injection rate the stronger the hindrance. As a consequence the density profiles in the maximum current phase do not depend at all on $`\alpha `$ in the middle and at the end of the system. At the beginning of the system, however, the density profiles decay as $`\text{i}^\gamma `$ with $`\gamma `$ $``$ 0.66 what is valid for all $`\text{v}_{max}`$ $`>`$ 1. We conjecture that $`\gamma `$ = $`\frac{2}{3}`$ because the exponent converges to this value with increasing system sizes. That means that the cases $`\text{v}_{max}`$ = 1 and $`\text{v}_{max}`$ $`>`$ 1 belong to different universality classes since the corresponding density profiles for the ASEP with parallel update decay as $`\text{i}^\gamma `$ with $`\gamma `$ = $`\frac{1}{2}`$ at the beginning of the system , . The free flow/maximum current transition is nicely reflected by the density at the end of the system: $`\rho `$(i=L) is proportional to the injection rate when the cars move freely up to $`\alpha _c`$ = 0.35 and becomes constant in the maximum current regime. The situation is different when density profiles at the free flow - jamming border for $`\beta `$ = 0.7 are considered (Fig 5b). It can be easily seen that the transition is of first order as the density profile at $`\alpha _c`$ = 0.278 is linear which is a typical feature of a first-order phase transition (see and references therein). The density profiles for the injection rate $`\alpha `$ = 1 are shown in Fig 5c. The course of the curves for $`\beta `$ $`<`$ 0.89 ($`\beta `$ $`>`$ 0.89) is typical for the maximum current (BII jamming) phase with an algebraic (exponential) decay at the beginning of the system due to the hindrance already described for $`\beta `$ = 1. The decay becomes weaker with decreasing $`\beta `$ and finally vanishes as the repercussion resulting from the blockage at i = L+1 increasingly superimposes the hindrance effect at the beginning of the system. For extinction rates 0.6 $``$ $`\beta `$ $``$ 0.9 there is a slight increase at the end of the system which indicates a hindrance due to the blockage. It just remains a border effect, however, and is of no relevance for the considerations in this article. ## IV Correlation Functions In this chapter we consider short-range correlation functions $`\text{C(i,t) = }<\eta \text{(i’,t’)}\eta \text{(i+i’,t+t’)}>_{\text{i’,t’}}\text{ - }<\eta \text{(i’,t’)}>_{\text{i’,t’}}^2`$ (16) where $`\eta \text{(i’,t’) = 1}\text{ if site i’ is occupied at time t ’}`$ (17) $`\eta \text{(i’,t’) = 0}\text{ else}`$ (18) $`<\mathrm{}>_{\text{i’,t’}}`$ describes the spatial and temporal average over all L sites i’ and over times t’ taken from our simulation of the steady state. The correlation functions are measured in the middle of the system where the influence of the boundaries is minimal. We do not only investigate the cases $`\beta `$ = 1 (influence of the left boundary, Fig 6a) and $`\alpha `$ = 1 (influence of the right boundary, Fig 6b), but also $`\beta `$ = 1-$`\alpha `$ (Fig 6c). For the latter case there are similar conditions at both boundaries and the system can be compared at best with the corresponding system with periodic boundary conditions. Therefore it is no surprise that the correlation functions in Fig 6c are qualitatively the same as those for periodic boundary conditions . What is interesting, however, is that a classification into free flow, maximum current, and jamming cannot be done when short-range correlation functions are considered. Instead, due to Figs 6a-c three regimes can be distinguished from each other: (a) Free Flow: The free flow is characterized by anticorrelations around a (21) propagating peak at i = $`\text{v}_{max}`$(t-1) with a shoulder at i = $`\text{v}_{max}`$ t, that is, in free flow moving cars are surrounded by empty space (b) Coexistence Regime (JI + Maximum current): The coexistence of (25) free flow and jamming manifests itself in the double peak structure of the correlation function. The jamming causes a maximum at i = -1 according to the hindrance the back car feels in the jam. (c) Jamming II (JII, ”Superjamming”): The propagating peak disappears (28) as a consequence of the fact that in free flow moving cars do not exist any longer. As it is obvious from the previous sections the transition from JI to JII is not a phase transition and does not change the behaviour of the system in an essential way. In correspondence with the critical injection (extinction) rate $`\alpha _{c2}`$ ($`\beta _{c2}`$) for the JI-JII transition is defined by the vanishing of the propagating peak and takes place at $`\beta _{c2}`$ $``$ 0.65 (see also Fig 1b). Unfortunately, an exact value for the Free Flow-Jamming transition can neither be given. As a consequence of the randomized oscillations in the density at the beginning of the system it is not possible to determine the critical injection (extinction) rate at which the influence of the right boundary reaches the left boundary. This is a significant difference to the deterministic case where the oscillations of the free flow phase form a well-defined pattern due to the lack of randomization. The transition from Jamming I to Jamming II takes place at $`\beta `$ $``$ 0.65 for all $`\text{v}_{max}`$ $`>`$ 1. In other words: When the maximum velocity is varied Free Flow, JII and coexistence phase (JI+maximum current) keep constant and only the ratio between maximum current phase and JI changes (Figs 1a,b). ## V Conclusions The nondeterministic Nagel-Schreckenberg model depends on the randomization, the maximum velocity, and the boundary conditions: The buffer effect observed for the deterministic case p = 0 and $`\text{v}_{max}`$ $``$ 3 is strongly weakened with increasing randomization probability. For $`\text{v}_{max}`$ = 2 there are no buffers and therefore the corresponding phase diagram is similar to the case $`\text{v}_{max}`$ = 1 for all p values. For $`\text{v}_{max}`$ $``$ 3 and p $`>`$ $`\text{p}_c`$ ($`\text{p}_c`$ = 0.1172 $`\pm `$ 0.008 for $`\text{v}_{max}`$ = 5) the buffer effect completely vanishes since the development of jamming waves is no longer suppressed. As a consequence, a maximum current phase occurs for p $`>`$ $`\text{p}_c`$ and the free flow (jamming) phase can be divided into two regimes AI and AII (BI and BII) similarly to the case of $`\text{v}_{max}`$ = 1,2. Another analogy to the case $`\text{v}_{max}`$ = 1 is that the free flow/jamming (free flow/maximum current and jamming/maximum current) transition is of first (second) order. There are, however, essential differences between systems with $`\text{v}_{max}`$ = 1 and $`\text{v}_{max}`$ $`>`$ 1: In the maximum current phase the density profiles decay algebraically with an exponent $`\gamma `$ = $`\frac{2}{3}`$ for $`\text{v}_{max}`$ $``$ 2 whereas $`\gamma `$ = $`\frac{1}{2}`$ was obtained in the ASEP. This indicates that systems with $`\text{v}_{max}`$ $`>`$ 1 and $`\text{v}_{max}`$ = 1 belong to different universality classes. Another difference to the ASEP is the existence of oscillations at the beginning of the system due to the repulsion of the cars. The comparison of systems with periodic and with open boundary conditions suggests that there are mainly three differences. First of all, the transition from free flow to jamming for systems with open boundaries is sharp and there is no maximum current phase in the case of periodic boundary conditions. Moreover, the dependence on the maximum velocity is more complex for systems with open boundary conditions due to the occurrence of the buffers. However, there are common features, too: Measurements of the short-range correlation function show that – as for corresponding systems with periodic boundary conditions – three regimes can be distinguished from each other: (a) Free Flow: cars do not hinder each other (b) Maximum Current and Jamming I: coexistence of freely moving and jammed cars (c) Jamming II: cars are jammed. ## VI Acknowledgments This work was supported by the State of North Rhine-Westphalia and by OTKA(T029985). ## FIGURE CAPTIONS Fig 1a: Phase diagram for $`\text{v}_{max}`$ = 2 (p = 0.5, continuous line: first-order or continuous phase transition, dotted line: border between AI/AII and BI/BII, broken line: border between the JI and the JII regime). Although there is no symmetry along the $`\alpha `$ = $`\beta `$ \- line the phase diagram shows strong similarities to the $`\text{v}_{max}`$ = 1 case. Fig 1b: Phase diagram for $`\text{v}_{max}`$ $``$ 5 (p = 0.5, continuous line: first-order or continuous phase transition, dotted line: border between AI/AII and BI/BII, broken line: border between the JI and the JII regime). Phase diagrams for $`\text{v}_{max}`$ = 3,4 are qualitatively the same. Due to the buffer effect the border between the free flow and the jamming phase shows a slight bending. Fig 1c: Phase diagram for $`\text{v}_{max}`$ = 5 and p = 0, 0.1172, 0.25, and 0.5. For randomization probabilities p $`>`$ $`\text{p}_c`$ with $`\text{p}_c`$ = 0.1172 $`\pm `$ 0.008 the buffer effect vanishes and the free flow regime becomes smaller with increasing randomization probability. The borders of the maximum current regime for $`\alpha `$ $`>`$ 0.6, $`\beta `$ $`>`$ 0.92, and p = 0.25 ($`\alpha `$ $`>`$ 0.35, $`\beta `$ $`>`$ 0.89, and p = 0.5) are represented by broken lines. The phase diagram for maximum velocities $`\text{v}_{max}`$ $`>`$ 5 is very similar. Fig 2: Average density in the middle of the system for $`\text{v}_{max}`$ = 5 (p = 0.5). The first-order phase transition from freely moving to jammed traffic can be clearly seen. Moreover, there is a jump in the derivative of $`\rho `$(i,$`\frac{\text{L}}{\text{2}}`$) at the maximum current/free flow and the maximum current/jamming border what is a hint at a continuous phase transition. Fig 3a: Current q for $`\text{v}_{max}`$ = 2, 3, $`\mathrm{}`$, 10 and $`\beta `$ = 1 (p = 0.5). At $`\alpha `$ $``$ 0.35 there is a slight maximum for $`\text{v}_{max}`$ $``$ 5 which is explained in Fig 4a. Fig 3b: Current q for $`\text{v}_{max}`$ = 2, 3, $`\mathrm{}`$, 10 and $`\alpha `$ = 1 (p = 0.5). For $`\text{v}_{max}`$ $``$ 5 the curves are identical. Fig 4a: Current q for p = 0, 0.125, $`\mathrm{}`$, 0.875 and $`\beta `$ = 1 ($`\text{v}_{max}`$ = 10). The maximum at $`\alpha `$ $``$ 0.81 for p = 0 moves towards smaller injection rates with increasing randomization probabilities and for p $`>`$ 0.125 a continuous phase transition is observed. Fig 4b: Current q for p = 0, 0.125, $`\mathrm{}`$, 0.875 and $`\alpha `$ = 1 ($`\text{v}_{max}`$ = 5). For p $`>`$ 0, the buffer effect observed for the deterministic case vanishes and the first-order phase transition goes over into a continuous phase transition. Fig 5a: Density profiles for $`\beta `$ = 1 (p = 0.5, $`\text{v}_{max}`$ = 5). A typical feature of density profiles in the free flow regime are oscillations resulting from the hindrance the cars feel at the beginning of the system from each other. These oscillations die out for higher system sites due to randomization. The curves decay algebraically in the maximum current regime as it is already known from the ASEP. Fig 5b: Density profiles for $`\beta `$ = 0.7 ($`\text{v}_{max}`$ = 5, p = 0.5). The phase transition at $`\alpha `$ = 0.278 is of first order characterized by a linear density profile at the critical injection rate. The curve for $`\alpha `$ = 0.4 ($`\alpha `$ = 0.3 or $`\alpha `$ = 0.29) is typical for a density profile in the BII (BI) regime. Fig 5c: Density profiles for $`\alpha `$ = 1 (p = 0.5, $`\text{v}_{max}`$ = 5). The results from the $`\text{v}_{max}`$ = 1 case are recovered: In the maximum current regime ($`\beta `$ $`>`$ 0.89) the density profiles decay algebraically and in the BII-jamming regime they are described by an enhanced exponential function. Fig 6a: Correlation functions for $`\beta `$ = 1 (p = 0.5, $`\text{v}_{max}`$ = 5). The free flow regime is characterized by a propagating peak with anticorrelations around it. In the coexistence regime both the jamming and the propagating peak can be observed. Fig 6b: Correlation functions for $`\alpha `$ = 1 (p = 0.5, $`\text{v}_{max}`$ = 5). The curves in the maximum current and in the JI-jamming regime behave similarly. At $`\beta `$ $``$ 0.75 the propagating peak vanishes and the system is in the JII-Jamming ( = ”Superjamming”) regime. Fig 6c: Correlation functions for $`\beta `$ = 1-$`\alpha `$ (p = 0.5, $`\text{v}_{max}`$ = 5). As for the injection rate $`\alpha `$ = 1, the propagating peak vanishes at $`\beta `$ = 1-$`\alpha `$ $``$ 0.75. Moreover, there is a striking similarity to corresponding correlation functions in the case of periodic boundary conditions.
warning/0006/quant-ph0006042.html	ar5iv	text	# On quantum logic operations based on photon-exchange interactions in an ensemble of non-interacting atoms ## I Introduction In some recent papers Franson et al., suggested that exchange interactions of two photons in a macroscopic ensemble of identical, non-interacting atoms could lead to large conditional phase shifts. In contrast to “conventional” nonlinear optics which requires scattering of both photons from the same atom, exchange interactions are present even when the two photons interact with different atoms. This makes them much more likely to occur in a dense medium. The large magnitude of the predicted conditional phase shifts would make such systems very attractive for quantum logical operation. However, whether or not exchange interactions are capable of generating entanglement between photons has been subject of some debate . In view of the claimed potential advantages, the requirements and limitations of the proposed schemes need to be examined. In the present note I want to discuss a special type of exchange interactions. In particular I will analyze the possibility to entangle photon states through interactions in an ensemble of atoms under the conditions considered in . Namely: (i) All processes are unitary, i.e. losses are negligible; (ii) The atomic system returns to the same state as before the interaction; (iii) The “conventional” nonlinear response of the atoms is assumed to be negligible; (iv) It is assumed that there are no atom–atom interactions, except those through the quantized radiation modes under consideration. Conditions (i) and (ii) enshure that the pair of qubits, represented by the photons undergoes an effective unitary evolution and is asymptotically disentangled from the atoms and the environment. It will be shown in the following that in a system that fulfills conditions (i-iv) entanglement between a pair of photons in distinguishable modes can not be generated. Any initially factorizable state will evolve into a factorizable state. ## II Model and effective time-evolution operator Let me consider the interaction of the quantized radiation field with a large number of identical atoms in dipole and rotating-wave approximation as proposed in . In addition to the photon field, the atoms may be coupled to some external classical fields to allow for manipulations of the states after or during the interaction with the photons. The Hamiltonian of the system has the following general form $$H=H_{\mathrm{field}}+H_{\mathrm{atom}}(t)+V,$$ (1) where $`H_{\mathrm{field}}`$ is the free Hamiltonian of the quantized photon field and $`H_{\mathrm{atom}}(t)`$ is the free Hamiltonian of the atoms including the interaction with the (time-dependent) external, classical fields. For simplicity it is assumed that each mode of the photon field couples only to one atomic transition. It is however straight forward to lift this restriction. The interaction operator has thus the following general structure $`V=\mathrm{}{\displaystyle \underset{k}{}}g_k{\displaystyle \underset{j=1}{\overset{N}{}}}\left[\widehat{\sigma }_{j,k}^{}\widehat{a}_kf_k(\stackrel{}{r}_j)+\widehat{\sigma }_{j,k}\widehat{a}_k^{}f_k^{}(\stackrel{}{r}_j)\right].`$ (2) Here $`\widehat{a}_k`$ and $`\widehat{a}_k^{}`$ are annihilation and creation operators of the photon field. $`k`$ is a mode index and $`f_k(\stackrel{}{r})`$ is the associated mode function. $`f_k`$ is not restricted to plane waves but could also represent e.g. localized wave packets, distinguishable by their arrival time. The modes are assumed to be orthogonal, such that $`[a_k,a_k^{}^{}]=\delta _{kk^{}}`$. $`\widehat{\sigma }_{j,k}`$ denotes a flip operator of atom $`j`$ corresponding to the transition coupled to the mode $`k`$ with coupling strength $`g_k`$. (Introducing flip operators for different $`k`$-values takes into account that the individual modes of the quantized field may be coupled to different dipole transitions.) It is assumed that initially ($`t=t_0`$) all atoms are in their ground states, i.e. the total initial state vector has the form $$\|\psi (t_0)=\|\varphi (t_0)\|g,$$ (3) where $`\|\varphi (t_0)`$ is the initial field state and $`\|g`$ the collective ground state of the atoms. The Schrödinger-equation for the state vector in the interaction picture can formally be solved by $`\|\psi (t)`$ $`=`$ $`𝖳\mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _{t_0}^t}dt^{}V(t^{})\right\}\|\psi (t_0),`$ (4) where $`𝖳`$ is the time ordering operator. It is clear that photon-atom interactions in general entangle both sub-systems. This is however not of interest here. The question I want to address is, whether the interaction can generate an entangled state of the photons given that the atomic system returns to its initial ground state at some time $`t_1`$. Thus we require $$\|\psi (t_1)\|\varphi (t_1)\|g.$$ (5) In this case the atomic and photonic components of $`\|\psi (t_1)`$ factorize and the photonic part is given by $`\|\varphi (t_1)=g\|𝖳\mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _{t_0}^{t_1}}dt^{}V(t^{})\right\}\|g\|\varphi (t_0)=𝖲(t_1,t_0)\|\varphi (t_0).`$ (6) The operator $`S`$ describes the conditional evolution of the photon field when the atomic system returns to its ground state. In order to calculate the action of $`S`$, we make use of a generalization of the cumulant generation function for a classical statistical variable $`X`$ $`\mathrm{exp}\left\{sX\right\}_X=\mathrm{exp}\left\{{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{s^m}{m!}}X^m\right\}.`$ (7) Here $`X^m`$ denotes the $`m`$th order cumulant, i.e. $`X=X`$, $`XY=XYXY`$ etc. Applying eq.(7) to $`𝖲`$ yields $`𝖲(t_1,t_0)`$ $`=`$ $`𝖳\mathrm{exp}\{{\displaystyle }\mathrm{d1}{\displaystyle }\mathrm{d2}\widehat{a}_{k_1}^{}(\tau _1)𝒫(1;2)\widehat{a}_{k_2}(\tau _2)`$ (9) $`+{\displaystyle }\mathrm{d1}{\displaystyle }\mathrm{d2}{\displaystyle }\mathrm{d3}{\displaystyle }\mathrm{d4}\widehat{a}_{k_1}^{}(\tau _1)\widehat{a}_{k_2}^{}(\tau _2)𝒫^{(2)}(1,2;3,4)\widehat{a}_{k_3}(\tau _3)\widehat{a}_{k_4}(\tau _4)+\mathrm{}\}`$ where $`\mathrm{d1}`$ stands for integration over time $`\tau _1`$ and summation over the mode index $`k_1`$. It was assumed here for simplicity that the average dipole moment of the atoms vanishes. $`𝒫(1,2)={\displaystyle \underset{j}{}}𝒫^j(1,2),`$ (10) where $`𝒫^j(1,2)=g_k^2f_{k_1}^{}(\stackrel{}{r}_j)f_{k_2}(\stackrel{}{r}_j)𝖳\widehat{\sigma }_{jk_1}^{}(\tau _1)\widehat{\sigma }_{jk_2}(\tau _2)`$ (11) describes the linear response of the $`j`$th atom to the quantized radiation field. The higher-order terms $`𝒫^{(n)}`$ characterize the “conventional” nonlinear response. The scattering of two photons off the same atom is for example determined by $`𝒫^{(2)}`$. It should be emphasized here, that cumulants containing operators of different atoms vanish, since it was assumed that atom–atom correlations can be built up only by the quantized radiation field. As a consequence each term $`𝒫^{(n)}`$ scales only linearly with the number of atoms $`N`$. Thus “conventional” nonlinear interactions of increasing order require increasing photon densities or large coupling constants $`g_k`$. Franson et al. argued in that a nonlinear phase shift between two photons could emerge even if the “conventional” nonlinear couplings, characterized by the higher-order cumulants in eq.(9), are negligible. Such phase shifts should arize from exchange interactions resulting from to the symmetrization requirements imposed by the bosonic nature of the photons. Let me therefore consider in the following the case were all higher-order cumulants are neglected. In this situation $`𝖲`$ reduces to: $`𝖲𝖳\mathrm{exp}\left\{{\displaystyle \mathrm{d1}\mathrm{d2}\widehat{a}_{k_1}^{}(\tau _1)𝒫(1,2)\widehat{a}_{k_2}(\tau _2)}\right\}.`$ (12) It should be emphasized that although the evolution operator (12) is bilinear in the photon operators, it takes fully into account any exchange interaction. The implicit summation over mode indices accounts for processes where photon 1 is seen by atom $`A`$ and photon 2 by atom $`B`$ as well as the case where photon 1 is seen by atom $`B`$ and photon 2 by atom $`A`$. It will now be shown that the conditional evolution $`t_0t_1`$ described by $`𝖲`$ cannot generate entanglement. I.e. any initially factorizable state will evolve into a factorizable state after the interaction. ## III State evolution In order to discuss the evolution of photons described by $`𝖲`$ in (12), I consider the case of the field initially being in a factorizable two-mode state with at most one photon in each mode. $`\|\varphi (t_0)=\|\varphi _1\|\varphi _2\|\{0_k\}`$ with $`\|\varphi _1=\left(\alpha _1+\beta _1\widehat{a}_{k_1}^{}\right)\|0_1,\|\varphi _2=\left(\alpha _2+\beta _2\widehat{a}_{k_2}^{}\right)\|0_2.`$ (13) Here $`\|0_1,\|0_2`$ are the vacuum states of modes $`k_1`$ and $`k_2`$ and $`\|\{0_k\}`$ is the vacuum state of all other modes. I proceed with discussing the evolution of the individual components of $`\|\varphi (t_0)`$. The vacuum component remains of course unaffected and it is sufficient to consider $`\|\chi _1(t_1)`$ $`=`$ $`𝖲(t_1,t_0)\|\chi _1(t_0)=𝖲(t_1,t_0)\widehat{a}_{k_1}^{}(t_0)\|0,`$ (15) $`\mathrm{and}`$ $`\|\chi _{1,2}(t_1)`$ $`=`$ $`𝖲(t_1,t_0)\|\chi _{1,2}(t_0)=𝖲(t_1,t_0)\widehat{a}_{k_1}^{}(t_0)a_{k_2}^{}(t_0)\|0.`$ (16) To formally calculate these expressions we make use of Wick’s theorem, which states that a time-ordered operator expression can be replaced by the sum of all normally ordered expressions with all possible “contractions”. Contractions refer here to a replacement of any operator pairs $`\widehat{a}_k^{}^{}(\tau ^{})`$ and $`\widehat{a}_{k^{\prime \prime }}(\tau ^{\prime \prime })`$ by the $`𝖳`$-ordered propagator $$𝒟(1,2)=0\|𝖳\widehat{a}_k^{}^{}(\tau _1)\widehat{a}_{k^{\prime \prime }}(\tau _2)\|0.$$ (17) We first note that since $`t_0`$ is the smallest time, the creation operators $`\widehat{a}_{k_1}^{}(t_0)`$ and $`\widehat{a}_{k_2}^{}(t_0)`$ in eqs.(15) and (16) can be included in the $`𝖳`$-ordering. Since $`𝖲\widehat{a}_{k_1}^{}(t_0)`$ and $`𝖲\widehat{a}_{k_1}^{}(t_0)\widehat{a}_{k_2}^{}(t_0)`$ respectively act on the vacuum state, out of all normally ordered expressions only those survive which have no photon annihilation operator left. Now $`𝖲\widehat{a}_{k_1}^{}(t_0)`$ can be expanded into a power series and Wick’s theorem applied to each term. This leads to the following perturbation series $`\|\chi _1(t_1)`$ $`=`$ $`\{\widehat{a}_{k_1}^{}(t_0)+{\displaystyle }\mathrm{d1}{\displaystyle }\mathrm{d2}𝒟(0,1)[𝒫(1,2)+{\displaystyle }\mathrm{d3}{\displaystyle }\mathrm{d4}𝒫(1,3)𝒟(3,4)𝒫(4,2)`$ (19) $`+{\displaystyle }\mathrm{d3}{\displaystyle }\mathrm{d4}{\displaystyle }\mathrm{d5}{\displaystyle }\mathrm{d6}𝒫(1,3)𝒟(3,4)𝒫(4,5)𝒟(5,6)𝒫(6,2)+\mathrm{}]\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)\}\|0,`$ where $`\|0`$ denotes the vacuum of all field modes. The first term results from contractions of photon operators within $`𝖲`$. The other terms arise from all possible contractions of $`\widehat{a}_{k_1}^{}`$ with operators from $`𝖲`$. Eq.(19) can be given the compact form $`\|\chi _1(t_1)`$ $`=`$ $`\left[\widehat{a}_{k_1}^{}(t_0)+{\displaystyle \mathrm{d1}\mathrm{d2}𝒟(0,1)\mathrm{\Pi }(1,2)\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)}\right]\|0`$ (20) where $`\mathrm{\Pi }(1,2)`$ is the solution to the linear integral equation (Dyson equation) $`\mathrm{\Pi }(1,2)=𝒫(1,2)+{\displaystyle \mathrm{d2}\mathrm{d3}𝒫(1,3)𝒟(3,4)\mathrm{\Pi }(4,2)}.`$ (21) In fact one easily verifies that an interactive solution of this equations generates the whole perturbation series of (19). That the quantum evolution can formally be solved in such a simple way is not surprising since the system is linear. Eq.(21) describes nothing else than multiple scattering of the incoming photon at the atoms with all nonlinearities being absent. In a diagrammatic language, the Dyson equation (21) corresponds to a sum of chain-like diagrams without branching or merging. In a similar way as above one can proceed with $`𝖲\widehat{a}_{k_1}^{}\widehat{a}_{k_2}^{}`$. In this case contractions only within $`𝖲`$ generate a term proportional to the product $`\widehat{a}_{k_1}^{}\widehat{a}_{k_2}^{}`$ similar to the first term in eq.(19). Then two series of terms emerge where either $`\widehat{a}_{k_1}^{}`$ or $`\widehat{a}_{k_2}^{}`$ is contracted with operators from $`𝖲`$. These leads to expressions identical to the higher-order terms in (19) multiplied with either $`\widehat{a}_{k_1}^{}`$ or $`\widehat{a}_{k_2}^{}`$. Finally there is a series of terms resulting of contractions of both $`\widehat{a}_{k_1}^{}`$ and $`\widehat{a}_{k_2}^{}`$ with operators from $`𝖲`$. This yields $`\|\chi _{12}(t_1)=\{\widehat{a}_{k_1}^{}(t_0)\widehat{a}_{k_2}^{}(t_0)+`$ (22) $`+{\displaystyle \mathrm{d1}\mathrm{d2}𝒟(0^{},1)\left[𝒫(1,2)+\mathrm{d3}\mathrm{d4}𝒫(1,3)𝒟(3,4)𝒫(4,2)+\mathrm{}\right]\widehat{a}_{k_2}^{}(t_0)\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)}+`$ (23) $`+{\displaystyle \mathrm{d1}\mathrm{d2}𝒟(0^{\prime \prime },1)\left[𝒫(1,2)+\mathrm{d3}\mathrm{d4}𝒫(1,3)𝒟(3,4)𝒫(4,2)+\mathrm{}\right]\widehat{a}_{k_1}^{}(t_0)\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)}+`$ (24) $`+{\displaystyle }\mathrm{d1}{\displaystyle }\mathrm{d2}𝒟(0^{},1)[𝒫(1,2)+{\displaystyle }\mathrm{d3}{\displaystyle }\mathrm{d4}𝒫(1,3)𝒟(3,4)𝒫(4,2)+\mathrm{}]\times `$ (25) $`\times {\displaystyle }\mathrm{d}\stackrel{~}{1}{\displaystyle }\mathrm{d}\stackrel{~}{2}𝒟(0^{\prime \prime },\stackrel{~}{1})[𝒫(\stackrel{~}{1},\stackrel{~}{2})+{\displaystyle }\mathrm{d}\stackrel{~}{3}{\displaystyle }\mathrm{d}\stackrel{~}{4}𝒫(\stackrel{~}{1},\stackrel{~}{3})𝒟(\stackrel{~}{3},\stackrel{~}{4})𝒫(\stackrel{~}{4},\stackrel{~}{2})+\mathrm{}]\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)\widehat{a}_{\stackrel{~}{k}^{\prime \prime }}^{}(\stackrel{~}{\tau }_2)\}\|0.`$ (26) Here $`0^{}`$ and $`0^{\prime \prime }`$ stand for $`\{t_0,k_1\}`$ and $`\{t_0,k_2\}`$ respectively. This expression can again be brought into a compact form $`\|\chi _{12}(t_1)=\widehat{a}_{k_1}^{}(t_0)\widehat{a}_{k_2}^{}(t_0)\|0`$ (27) $`+{\displaystyle \mathrm{d1}\mathrm{d2}𝒟(0^{},1)\mathrm{\Pi }(1,2)\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)\widehat{a}_{k_2}^{}(t_0)\|0}+{\displaystyle \mathrm{d1}\mathrm{d2}𝒟(0^{\prime \prime },1)\mathrm{\Pi }(1,2)\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)\widehat{a}_{k_1}^{}(t_0)\|0}`$ (28) $`+{\displaystyle \mathrm{d1}\mathrm{d2}𝒟(0^{},1)\mathrm{\Pi }(1,2)d\stackrel{~}{1}d\stackrel{~}{2}𝒟(0^{\prime \prime },\stackrel{~}{1})\mathrm{\Pi }(\stackrel{~}{1},\stackrel{~}{2})\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)\widehat{a}_{\stackrel{~}{k}^{\prime \prime }}^{}(\stackrel{~}{\tau }_2)\|0}.`$ (29) One immediately recognizes that $`\|\chi _{12}(t_1)`$ can be written as $`\|\chi _{12}(t_1)`$ $`=`$ $`\left[\widehat{a}_{k_1}^{}(t_0)+{\displaystyle \mathrm{d1}\mathrm{d2}𝒟(0^{},1)\mathrm{\Pi }(1,2)\widehat{a}_{k^{\prime \prime }}^{}(\tau _2)}\right]`$ (31) $`\left[\widehat{a}_{k_2}^{}(t_0)+{\displaystyle d\stackrel{~}{1}d\stackrel{~}{2}𝒟(0^{\prime \prime },\stackrel{~}{1})\mathrm{\Pi }(\stackrel{~}{1},\stackrel{~}{2})\widehat{a}_{\stackrel{~}{k}^{\prime \prime }}^{}(\stackrel{~}{\tau }_2)}\right]\|0`$ The evolution of $`\|\varphi `$ from $`t_0`$ to $`t_1`$ is hence given by $`\|\varphi (t_0)=\left(\alpha _1+\beta _1\widehat{a}_{k_1}^{}\right)\left(\alpha _2+\beta _2\widehat{a}_{k_2}^{}\right)\|0`$ (33) $``$ $`\|\varphi (t_1)=`$ $`[\alpha _1+\beta _1(\widehat{a}_{k_1}^{}(t_0)+{\displaystyle }\mathrm{d1}{\displaystyle }\mathrm{d2}𝒟(0^{},1)\mathrm{\Pi }(1,2)\widehat{a}_{k^{\prime \prime }}^{}(\tau _2))]`$ (35) $`[\alpha _2+\beta _2(\widehat{a}_{k_2}^{}(t_0)+{\displaystyle }\mathrm{d}\stackrel{~}{1}{\displaystyle }\mathrm{d}\stackrel{~}{2}𝒟(0^{\prime \prime },\stackrel{~}{1})\mathrm{\Pi }(\stackrel{~}{1},\stackrel{~}{2})\widehat{a}_{\stackrel{~}{k}^{\prime \prime }}^{}(\stackrel{~}{\tau }_2))]\|0.`$ Thus if the process starts with a factorizable state with photons in distinguishable modes, i.e. if $`(\alpha _1+\beta _1\widehat{a}_{k_1})\|0`$ is orthogonal to $`(\alpha _2+\beta _2\widehat{a}_{k_2})\|0`$ and if the process generates photons in distinguishable modes, i.e. if $`[\alpha _1+\beta _1(\widehat{a}_{k_1}^{}+{\displaystyle }{\displaystyle }𝒟\mathrm{\Pi }\widehat{a}_{k^{\prime \prime }}^{})]\|0\mathrm{and}[\alpha _2+\beta _2(\widehat{a}_{k_2}^{}+{\displaystyle }{\displaystyle }𝒟\mathrm{\Pi }\widehat{a}_{\stackrel{~}{k}^{\prime \prime }}^{})]\|0`$ (36) are orthogonal, then the generated state vector remains factorizable. ## IV Conclusion In the present note I have shown it is not possible to generate entanglement between photons using solely exchange interactions in a large ensemble of atoms, if the atoms are left in the same quantum state after the interaction as they were initially. From a diagrammatic point of view entanglement between photons can not be generated if all possible diagrams are chain-like. To produce entanglement non-trivially connected diagrams are needed, as emerge for example from nonlinear atomic responses or from atom–atom interactions due to e.g. dipole-dipole or collisional interactions. ## acknowledgment I would like to thank Mikhail Lukin and Tomas Opatrny for stimulating discussions on the subject and the Institute for Atomic and Molecular Physics for the hospitality. The financial support of the Deutsche Forschungsgemeinschaft is highly appreciated.
warning/0006/astro-ph0006207.html	ar5iv	text	# 1 Photometry of the host galaxies The Multiband Photometry of the GRB Host Galaxies V. V. Sokolov, T. A. Fatkhullin, V. N. Komarova Special Astrophysical Observatory of R.A.S., Karachai-Cherkessia, Nizhnij Arkhyz, 369167 Russia Email: sokolov@sao.ru, timur@sao.ru, vkom@sao.ru We present photometric multiband spectral energy distributions for the GRB host galaxies: GRB 971214, GRB 970508, GRB 980613, GRB 980703 and GRB 990123 obtained with the 6-m telescope of SAO RAS. Using SEDs for the starburst galaxies, we made estimates of K-correction values and estimated the absolute magnitudes of the GRB host galaxies within the range of cosmological parameters. The comparison of the broad band spectra of these galaxies with the spectra of galaxies of different morphological types (Connolly et al, 1995, AJ, 110, 1071) shows that the GRB host galaxies are best fitted by the spectral properties of S2-S5 averaged SEDs of starburst galaxies. Observations and data reduction The observations of the GRB host galaxies were performed using primary focus CCD photometer of the 6m telescope of SAO RAS during July-August 1998 for the GRB 971214, GRB 970508, GRB 980613 and GRB 980703 host galaxies, July 1999 for the GRB 990123 host galaxy and March-April 2000 for the GRB 991208 and GRB 000301C host galaxies. It was carried out with standard (Johnson-Kron-Cousins) photometric $`BVR_cI_c`$ system. Using the Landolt (Landolt, 1992) standard field, the photometric calibrations were performed. Using the Galactic extinction curve from Cardelli, Clayton & Mathis 1989, we obtained values of the foreground extinction. Tables 1 and 2 present the dereddened magnitudes and fluxes of the host galaxies instead of the uncorrected for Galactic extinction magnitudes of the GRB 980703 and GRB 990123 host galaxies are presented in Sokolov et al. 2000. Cosmological models Here we use three Friedman cosmological models: $$H_0=60\text{ km s}\text{-1}\text{ Mpc}\text{-1}\text{}\mathrm{\Omega }_m=1\text{}\mathrm{\Omega }_\mathrm{\Lambda }=0\text{ (A)}$$ $$H_0=60\text{ km s}\text{-1}\text{ Mpc}\text{-1}\text{}\mathrm{\Omega }_m=0\text{}\mathrm{\Omega }_\mathrm{\Lambda }=0\text{ (B)}$$ $$H_0=60\text{ km s}\text{-1}\text{ Mpc}\text{-1}\text{}\mathrm{\Omega }_m=0\text{}\mathrm{\Omega }_\mathrm{\Lambda }=1\text{ (C)}$$ For these models the relation $`\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_k=1`$ is valid, where $`\mathrm{\Omega }_m=\rho _08\pi G/3H_0^2`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=\mathrm{\Lambda }c^2/3H_0^2`$, and $`\mathrm{\Omega }_k=kc^2/R_0^2H_0^2`$. Here $`\rho `$, $`\mathrm{\Lambda }`$, $`k`$, and $`R`$ are density, cosmological constant, curvature constant, and radius of curvature, respectively, and “0” denotes the present epoch. Comparision $`BVR_cI_c`$ photometry with spectra It would be interestingly to compare our $`BVR_cI_c`$ broadband photometry to spectra of the host galaxies. Figures 1 and 2 present this comparision. Comparision with local starburst galaxies To compare our broad band spectra we have used the S1, S2, S3, S4, S5, S6 averaged spectral energy distrubutions (SEDs) for starburst galaxies from Calzetti et al. (1994). The spectra of starburst were grouped according to increasing values of the color excess $`E(BV)`$: from S1, with $`E(BV)=0.05`$ to S6, with $`E(BV)=0.7`$. It should be noted that this SEDs are not observed but are the templates that have been constructed using real observed starburst SEDs up to a redshift of $`0.03`$ (Connolly et al. 1995). The fluxes of starburst SEDs have been convolved with sensitivity functions of the $`BVR_cI_c`$ filters (sentivity functions have been used from Bessel, 1990) and the derived values was compared to our observed fluxes. For each SED the $`\chi ^2`$ was calculated. The values of the $`\chi ^2`$ was calculated in follow way: $$\chi ^2=\underset{i}{}\left(\frac{f_{host,i}kf_{template,i}}{\sigma _{f_{host,i}}}\right)^2$$ Here $`i`$ denote the filters ($`BVR_cI_c`$), $`f_{host,i}`$ is the flux of the GRB host galaxies in the filter $`i`$, $`f_{template,i}`$ is the convolved with filter $`i`$ flux of the template SED on effective wavelength of filter $`i`$, $`\sigma _{f_{host,i}}`$ is the error of flux of the GRB host galaxy in filter $`i`$ and $`k`$ is the normalization coefficient. The definition of group of the starburst allows us to compare our broadband spectra to SEDs of real galaxies. Table 3 from Calzetti et al. (1994) was used to select real galaxies. Figures 3, 4, 5 and 6 present comparision of the $`BVR_cI_c`$ photometry with average starburst and real galaxy SEDs. K-correction and absolute magnitudes According to the our definitions of spectral types of host galaxies we can to estimate the K-correction. Using the definition of the K-correction (see Oke & Sandage, 1968): $$K_i=2.5\mathrm{log}\left\{\left(1+z\right)\frac{F\left(\lambda \right)S_i\left(\lambda \right)𝑑\lambda }{F\left(\lambda /\left(1+z\right)\right)S_i\left(\lambda \right)𝑑\lambda }\right\}$$ (1) and SEDs from Connolly et al. (1995) estimates of K-correction for the $`B`$ band are: $`K_B=0.44`$ for the GRB 970508 host galaxy, $`K_B=0.01`$ for the GRB 980703 host galaxy, $`K_B=0.13`$ for the GRB 990123 host galaxy and $`K_B=0.46`$ for the GRB 991208. This estimates allow us to derive the absolute magnitudes of the host galaxies. The absolute magnitude $`M_i`$ in filter $`i`$ of the source can be calculated from magnitude-redshift relation: $$M_i=m_iK_i\left(z\right)5\mathrm{log}\left(R_{lum}/Mpc\right)25$$ (2) where $`m_i`$ is the observed magnitude in filter $`i`$, $`K_i(z)`$ is the K-correction at redshift $`z`$, $`R_{lum}`$ is the luminosity distance. Using the $`B`$ magnitudes from Table 1 and the K-correction given above for cosmological models (A), (B), (C) we yield: $`M_{B_{rest}}=18.08,18.53,19.09`$ for the GRB 970508 host galaxy, $`M_{B_{rest}}=20.60,21.12,21.73`$ for the GRB 980703 host galaxy, $`M_{B_{rest}}=20.20,21.02,21.82`$ for the GRB 990123 host galaxy and $`M_{B_{rest}}=18.29,18.68,19.18`$ for the GRB 991208 host galaxy. In the case of the GRB 970508 ($`z=0.835`$) and GRB 991208 (z=0.7063) host galaxies the $`I_c`$ band roughly correspond to the $`B`$ band in rest frame. This allow us to calculate directly from Eq. (1) the value of the K-correction for $`B`$-magnitude, replacing with $`2.5\mathrm{log}(F_{B_{rest}}/F_{\lambda _{B/(1+z)}})2.5\mathrm{log}(F_{I_{obs}}/F_{B_{obs}})`$. We derived: $`K_B=0.51\pm 0.32`$ and $`K_B=0.21\pm 0.31`$ for the GRB 970508 and GRB 991208 host galaxies respectively. Then absolute $`B`$-magnitudes of the host galaxies are: $`M_{B_{rest}}=18.15,18.60,19.16`$ for the GRB 970508 host galaxy and $`M_{B_{rest}}=18.04,18.43,18.93`$ for the GRB 991208 host galaxy. Conclusions We presented the multiband photometry of the seven GRB host galaxies. It should be noted that observations were carried out with one instrument and in one photometric system. To solve the progenitors problem we obviously need the statistics of photometric and spectroscopic properties of the host galaxies. Our photometry of the GRB 970508, GRB 980703, GRB 990123 and GRB 991208 hosts has shown that the broadband spectra are best fitted by starburst spectral energy distributions. Moreover, there is evidence showing that GRBs spatially coincident with a bright star-forming region (Bloom et al. 1999). Probably, GRBs are associated with young stellar population what can be an evidence of SNe—GRBs connection.
warning/0006/astro-ph0006321.html	ar5iv	text	# 1 Trajectories in the 𝑇/𝑆 – (𝑛-1) plane. Squares indicate the initial choices for 𝑇/𝑆 and (𝑛-1); circles indicate the values 50 e-folds before the end of inflation. A trajectory ends when 𝑇/𝑆 and/or \|𝑛-1\| become large; most of inflation occurs when 𝑇/𝑆 and \|𝑛-1\| are small. The upper left panel shows a complete trajectory, with ticks indicating e-folds before the end of inflation (from the circle, 50,49,⋯,1). The other three panels show trajectories in more detail. Note how 𝑇/𝑆 and (𝑛-1) outside the attractor region are “pulled in” (the attractors are shown as broken lines and the boundary of the excluded region is a solid curve). ## Acknowledgments. We thank Andrew Liddle and David Spergel for valuable discussions and comments. This work was supported by the DoE (at Chicago and Fermilab) and by the NASA (at Fermilab by grant NAG 5-7092).
warning/0006/astro-ph0006340.html	ar5iv	text	# What do the UV Spectra of Narrow-line Seyfert 1 Galaxies tell us about their BLR? ## 1 Introduction Narrow-line Seyfert 1 galaxies (NLSy1), first suggested as a distinct class of AGN by Osterbrock and Pogge (1985), are characterized by Balmer lines whose FWHM is smaller than typical Seyfert 1 galaxies i.e. 500 $`<`$ FWHM $`<`$ 2000 km s<sup>-1</sup>, slightly broader than the forbidden lines. On the other hand they are clearly different from Seyfert 2 galaxies since the ratio of $`[`$OIII\] $`\lambda `$5007 to H$`\beta `$ is $`<`$3, i.e. below the limiting value found by Shuder and Osterbrock (1981) to discriminate between Seyfert 1 and Seyfert 2 galaxies. In NLSy1 strong Fe II optical multiplets and higher ionization iron lines (e.g.$`[`$FeVII\]6087Å and $`[`$FeX\] 6375Å) are often present. These are usually seen in Seyfert 1 and not in Seyfert 2 galaxies. Many NLSy1s have an unusually strong big blue bump (BBB) which, when compared to typical Seyfert 1 and QSO BBBs, is shifted towards higher energies, sometimes even out of the optical/UV range (at least one object actually peaks in the soft X-ray band: RE J1034+396, Puchnarewicz et al. 1995). Its high frequency tail is clearly seen in soft X-rays and these objects have generally steeper soft X-ray continua than is “typical” for Seyfert 1 galaxies (Boller, Brandt & Fink 1996), meaning that they have a stronger soft-X-ray excess over the hard X-ray power law. The intrinsic hard-X-ray continua of NLSy1s are also generally steeper (Brandt, Mathur, Elvis 1997) than in typical Seyfert 1s. The NLSy1s are usually only weakly absorbed in the soft X-rays (Boller, Brandt & Fink 1996) and in many cases both the UV flux and soft X-ray flux are strongly variable. NLSy1 objects are generally radio-quiet and their radio powers are typical of those found in other Seyfert galaxies (Ulvestad, Antonucci & Goodrich 1995). There is no widely adopted view on the basic reason why the continua of NLSy1 galaxies are different from classical Sy1. The two most probable explanations of the stronger big blue bumps in these objects are pole-on orientation (Puchnarewicz, Mason, Córdova 1994, Wilkes 1998), and higher accretion rate relative to the mass of the central object (e.g. Boller, Brandt, Fink 1996; Wandel 1997, Czerny, Witt & Życki 1997, Pounds, Done & Osborne 1995). Steeper hard X-ray spectra and strong permitted FeII lines are possibly a secondary effect of the atypical shape of the soft-X-ray continuum (Pounds, Done & Osborne 1995, Brandt, Mathur & Elvis 1997; Wilkes, Elvis and McHardy 1987, Shastri et al 1993). Wilkes et al. (1999), studying a sample of low redshift quasars and Sy1s, and their relations between optical/UV emission lines and the continuum, found that the four NLSy1 in their sample show smaller equivalent widths of CIII\] and CIV lines than typical AGN (EW(CIV) $`<`$ 40Å for NLSy1, while 30Å $`<`$ EW(CIV) $`<`$ 200Å for other AGN). In this paper we investigate in detail the UV line properties of a sample of NLSy1 objects to determine whether the weakness of the carbon lines is typical of these objects, constituting an additional property which distinguishes them from “normal” Sy1 galaxies. We investigate the physical conditions of the BLR which may explain these systematic differences. We also discuss the possibility that these objects have luminosities close to their Eddington luminosity. ## 2 UV line measurements ### 2.1 The sample From the currently known set of NLSy1 (Boller, Brandt & Fink 1996, Greiner et al. 1996, Puchnarewicz et al. 1992, Puchnarewicz et al. 1994, Brandt, Fabian & Pounds 1996, Grupe et al. 1996 , Moran et al. 1996, Brandt - private communication, Wilkes et al. 1999) we have defined a subset of 11 objects (Table 1) for which UV spectra are available either from the HST (5 objects) or IUE archives. The IUE spectra were taken from Lanzetta et al. (1993), and the reduced HST spectra from Dobrzycki (private communication, see also Bechtold et al. 2000). ### 2.2 Line parameters We have measured the EW (Table 2), line ratios (Table 3) and line widths (FWHM; Table 4) of all prominent UV lines: Ly$`\alpha `$$`\lambda `$1216, CIV $`\lambda `$1549, CIII\] $`\lambda `$1909, SiIII\] $`\lambda `$1892, AlIII $`\lambda `$1857, SiIV+OIV\] $`\lambda `$1400 blend, and MgII $`\lambda `$2798. The EW and FWHM of IUE spectra were measured using the splot task in IRAF: the EW by fitting a linear continuum to the data and integrating across the observed emission line (keystroke ‘e‘), the FWHM by measuring the width at half the flux in the line peak above the continuum. The same procedure was applied when the line parameters in the HST spectra were measured, although a different program findsl (provided by Aldcroft, Bechtold & Elvis 1994), specially written to handle the HST data, was used. The line parameters presented in Tables 2,3,4 have been corrected for absorption: for weak absorption by using a linear fit across the absorption line, for strong absorption (as in PG 1351+560 and PG 1411+442) by assuming a symmetric emission line profile and reflecting the unabsorbed wing about the peak. As has been noted by Vestergaard & Wilkes (2000), CIII\] is blended with FeIII UV34 $`\lambda `$1914 line, which should be taken into account especially when the CIII\] line is weak. In three spectra (IZw1, PG 1211+143, Mrk 478), where the FeIII UV34 was clearly visible, we subtracted this line (modeled as a Gaussian centered at $`\lambda `$=1914Å rest frame) from the CIII\] blend. ### 2.3 Comparison of NLSy1 with “normal” AGN In this section we compare the UV line properties of NLSy1 with Seyfert 1 galaxies and quasars. Fig. 1a shows the EW of Ly$`\alpha `$, CIV and MgII of our NLSy1 sample (shaded areas) compared to the sample of Seyfert 1 galaxies (dotted line) from Wu at al. (1983) (their sample includes three NLSy1: IZw1, Mrk 478, IIZw136 which we excluded here) and low redshift quasars from Wilkes et al. (1999), Corbin & Boroson (1996), and radio-loud quasars from Baldwin, Wampler & Gaskell (1989), combined and denoted by a dashed line. The EW of CIV and MgII lines are significantly smaller in NLSy1 than in the broad line Seyfert 1 galaxies and quasars. The K-S test yielded a 0.001 chance that the EW of CIV and MgII in NLSy1s and Sy1s are drawn from the same population. For Ly$`\alpha `$ the chance was $`<`$ 0.02. When compared to QSO the significance remained strong for CIV ($`p<0.01`$) and MgII ($`p<0.025`$), while for Ly$`\alpha `$ the distributions are similar ($`p>0.5`$). The smaller EW of the carbon and MgII lines cannot be due to a simple continuum increase, as this would effect the EW of all lines equally, while EW(Ly$`\alpha `$) is not significantly smaller. As can be seen from Table 3 the CIV/Ly$`\alpha `$ (mean 0.25$`\pm `$0.09), CIII\]/Ly$`\alpha `$ (mean 0.05$`\pm `$0.05) and the MgII/Ly$`\alpha `$ (mean 0.05$`\pm `$0.03) ratios are smaller compared to those typically observed in Seyfert 1 galaxies (Wu et al. 1983 give observed ranges: 0.35-2.01, 0.03-0.39, 0.07-0.63 respectively) and quasars (observed range: 0.3-1.04, 0.15-0.3, 0.15-0.35). To show this more clearly Figure 2 shows the CIII\]/Ly$`\alpha `$ vs CIV/Ly$`\alpha `$ line ratios for NLSy1 in our sample (denoted as filled squares) with the Seyfert 1 sample (denoted as circles) and quasars from Laor et al. (1995), Christiani & Vio (1990), Wilkes et al. (1999) and narrow line quasars from Baldwin et al. (1988). As the lines in NLSy1 are narrow, we can clearly resolve the components of the CIII\]+SiIII\]+AlIII blend (especially in the HST data, in the IUE data the S/N is often too low). From Table 3 it is also clear that the SiIII\] line in most of the NLSy1 is very strong compared to the CIII\] line. Also the SiIV+OIV\] blend is strong compared to CIV (mean SiIV+OIV\]/CIV ratio in NLSy1 is 0.49$`\pm `$0.26, larger than the mean ratio of 0.3 in quasars from Francis et al. 1991). However, the SiIV+OIV\]/Ly$`\alpha `$ ratio is in the range of normal AGN, indicating that the large SiIV+OIV\]/CIV ratio is due to weaker CIV emission. The AlIII doublet in NLSy1 is rather strong (equivalent width $``$ few Å - see Table 2). Although broader than H$`\beta `$ the UV lines in NLSy1 are narrow compared to other AGN. In Table 4 and Fig. 1b we compare our objects with samples of low redshift quasars from Corbin & Boroson (1996) and Wilkes et al. (1999), and with a radio-loud sample from Baldwin, Wampler & Gaskell (1989). ## 3 Discussion We will now investigate what the line strengths and line ratios tell us about the physical properties in the BLR clouds of NLSy1 (Section 3.1). Then we will discuss the continuum properties of NLSy1 (Section 3.2) and investigate what they indicate about their central engine. In conclusion we show how the deduced differences between the central engines of NLSy1 and “normal” AGN can explain their different, observed emission line spectra. ### 3.1 Physical properties of the BLR clouds in NLSy1 More than ten years ago Gaskell (1985) noticed that Seyfert 1 galaxies with narrow H$`\beta `$ lines of FWHM $`<`$ 1600 km s<sup>-1</sup> show lower H$`\beta `$ equivalent widths than typical Seyfert 1. He interpreted this finding as a result of collisional destruction of H$`\beta `$ in the higher density BLR clouds in these objects. Although we do not study the optical spectra of NLSy1 in this paper, we will now investigate whether the UV spectra lead to a similar conclusion. Rees et al. (1989) calculated the line intensities for different BLR cloud densities at constant column density ($`10^{23}`$ cm<sup>-2</sup>) and ionization parameter ($`U=10^2`$). They found that optically thick lines such as hydrogen and carbon lines have a fairly constant intensity up to a certain density, above which these lines become thermalized and their intensity drops considerably. For hydrogen lines, CIV and MgII this critical density is $`10^{10}`$cm<sup>-3</sup>. For the semi-forbidden lines CIII\] and SiIII\] it is around $`5\times 10^9`$ cm<sup>-3</sup> and $`10^{11}`$ cm<sup>-3</sup> respectively. The Ly$`\alpha `$ and CIV lines are usually strong coolants at densities smaller than these critical values, but as the density increases and these lines become thermalized, other high-excitation lines such as CIII $`\lambda `$977 and AlIII $`\lambda `$1857 take over the cooling. In the previous section we showed that the UV spectra of NLSy1s, when compared to “normal” Seyfert 1 and QSO galaxies, show weaker carbon and MgII lines. Although the wavelength of our spectra does not cover the range of CIII $`\lambda 977`$, the AlIII $`\lambda `$1857 doublet is clearly seen (where the S/N is high enough) and is especially strong in IZw1 (see Table 2). All these line properties suggest that in NLSy1 objects the BLR clouds have higher densities than the BLR clouds in “normal” AGN. We will estimate how much higher by studying the line ratios in the following section. a) The line ratios The CIII\] and CIV to Ly$`\alpha `$ ratios are often used as a density indicator. This is because the carbon lines are collisionally excited (hence sensitive to density), while Ly$`\alpha `$ is not (note however that Mathur et al. 1994 showed that for large ionizing parameters, $`U>0.1`$, where $`U=\frac{_{1Ryd}^{\mathrm{}}\frac{L_\nu }{h\nu }}{4\pi r^2cn_H}`$, CIII\] ceases to be a density indicator). These line ratios are also a sensitive function of the ionization parameter U, therefore we investigate the relation of these line ratios to density and U. We have calculated line ratios using the photoionization code CLOUDY (version 80.07, for reference see Ferland 1991). First as an input ionizing continuum we took a standard AGN continuum (table agn; Mathews & Ferland 1987). With this continuum, the observed lines were only reproduced with higher densities in the CIII\] than the CIV emitting clouds. This requires a steep increase in cloud density with radius, which is contrary to expectations and seems unrealistic. Then we used the spectral energy distribution (SED) of the NLSy1 PG 1211+143 as the ionizing continuum. Although no detailed study of the SEDs of NLSy1s has been made and is beyond the scope of this paper, PG 1211+143 is typical of those studied to date with $`\alpha _x=2.13\pm 0.22`$, where $`F_\nu =\nu ^\alpha `$ (Wang, Brinkmann & Bergeron 1996; where typical NLSy1 slopes are in the range 1.5 to 3.5 - see Boller, Brandt & Fink 1996) and $`\alpha _{io}=0.90`$ (i.e. slope measured between 1 $`\mu `$m and 2500Å, where typical NLSy1 values are 0.4 to 2.8 - see Lawrence et al. 1997). The IR to hard X-ray SED of this object was taken from Elvis et al. (1994) and is reproduced here in Fig. 3. The SED was linearly interpolated between the observational points in the optical/UV region. The EUV continuum was determined by a linear interpolation between the lowest energy point in the X-ray range and the highest in the UV, providing a conservative (i.e. low) estimate of the number of EUV photons. We investigated a range of cloud densities ($`n(H)=10^810^{13}`$ cm<sup>-3</sup>, well within the range of the applicability of the photoionization code CLOUDY<sup>1</sup><sup>1</sup>1“The hydrogen and atoms and ions of helium are treated in the code as 10-level atoms. The treatment of the heavy elements is not as complete as hydrogen and helium, but a 3-body recombination is included as a general recombination process. \[…\] The physical high-density limit is set by the approximate treatment of the three-body recombination-collisional ionization ($``$ $`10^{13}`$ cm<sup>-3</sup>) for the heavy elements and the approximate treatment of line transfer”. - see Ferland (1991).) and ionization parameters ($`U=10^310^1`$, where $`U=10^2`$ is the value for the “standard” BLR - Davidson & Netzer 1979). The metal abundances were assumed to be solar and the cloud column density 10<sup>23</sup> cm<sup>-2</sup>. The calculated line ratios are plotted in Fig. 4, where the observed line ratios for our NLSy1 are denoted by horizontal lines. The CIV/Ly$`\alpha `$ ratio (Fig. 4a) depends very strongly (more than any other line ratio) on the value of the ionization parameter as well as on the density of the BLR clouds. Small values of the ionization parameter ($`U=10^3`$) are clearly favored by our data, for which densities of the order of 10<sup>11</sup> cm<sup>-3</sup> to 10<sup>12</sup> cm<sup>-3</sup> are needed to produce the low observed CIV/Ly$`\alpha `$ ratios. For the same, small ionization parameter the observed CIII\]/Ly$`\alpha `$ ratios indicate densities between $`>10^9`$ cm<sup>-3</sup> and 10<sup>11</sup> cm<sup>-3</sup> (see Fig. 4b), which are smaller than the cloud densities inferred from the CIV/Ly$`\alpha `$ ratios. This strongly suggests that the CIII\] and CIV lines are formed in different clouds, which implies a stratified BLR. This is as expected from the results of reverberation mapping (e.g. NGC 5548 Korista et al. 1995, NGC 7469 Wanders et al. 1997, Fairall 9 Rodrigues-Pascual et al. 1997) which show that the CIV and Ly$`\alpha `$ line fluxes vary with a smaller time delay, relative to the UV continuum, than the CIII\] lines, indicating that CIV, Ly$`\alpha `$ emitting clouds lie nearer to the central engine than the CIII\] emitting clouds. The density of the CIV and Ly$`\alpha `$ emitting line region is typically estimated (e.g. Peterson et al. 1985) to be $``$ 10<sup>11</sup> cm<sup>-3</sup>, while the CIII\] region is $``$ 10<sup>9.5</sup> cm<sup>-3</sup>. Thus the density of the CIV, Ly$`\alpha `$ emitting clouds in our NLSy1s is comparable or somewhat ($`<`$ 10 times) larger than in normal AGN. To further constrain the density of the CIII\] emitting clouds, we need to investigate its ratio to a line which is formed in the same clouds for example SiIII\]. For $`U=10^3`$ the density of the emitting gas inferred from the SiIII\]/CIII\] line ratio is between 10<sup>9.5</sup> cm<sup>-3</sup> and 10<sup>10.5</sup> cm<sup>-3</sup> (see Fig. 4c, we omit here PG 1411+442, which is a BAL QSO and has $`n(H)10^9`$ cm<sup>-3</sup>). The ratio of SiIII\]/CIII\] is larger than the typical value of $`0.3\pm 0.1`$ seen in quasars (Laor et al. 1995) for all objects (except for BAL QSO PG 1411+442; see Table 3 and Fig. 4c) This high ratio is probably the result of the suppression of CIII\] while SiIII\] remains strong, due to the smaller critical density for CIII\] ($`5\times 10^9`$ cf. $`10^{11}`$ cm<sup>-3</sup> for SiIII\], see Section 3.1a). In Fig. 4d we present the observed and calculated SiIV+OIV\]/Ly$`\alpha `$ ratios. For $`U=10^3`$ the clouds have density of the order of 10<sup>11</sup> cm<sup>-3</sup> to 10<sup>12</sup> cm<sup>-3</sup>, similar to the range for CIV and Ly$`\alpha `$ emitting clouds. The MgII/Ly$`\alpha `$ ratio, on the other hand, is very small, and cannot be reproduced by clouds with small ionization parameter even when the column density, $`N_H`$, is varied over the range $`10^{22}`$ \- $`10^{24}`$ cm<sup>-2</sup>. Only larger than standard ionization parameters, in the range $`10^1`$ to $`10^2`$, can produce the observed ratios (Fig. 4e). This may indicate that the MgII lines do not form in the same region as the other BLR lines. We will return to this problem in the next section. To summarize the results of this subsection we conclude that the unusual UV line ratios in the NLSy1 objects can be explained if the BLR clouds have 10 times lower ionizing parameters ($`\mathrm{log}U3`$) and a few times ($`<`$10) higher densities (n(H)$``$ 10<sup>11-12</sup> cm<sup>-3</sup> for Ly$`\alpha `$, CIV, SiIV emitting clouds and 10<sup>9.5-10.5</sup> cm<sup>-3</sup> for CIII\], SiIII\] emitting clouds) than normal AGN. The BLR is clearly stratified with CIV, Ly$`\alpha `$, SiIV producing clouds lying closer to the central engine, while CIII\] and SiIII\] emitting clouds lying predominantly further out. The MgII emission cannot be produced by the same cloud population, suggesting that these lines form in a different region. Although clouds with a wide range of properties are likely to exist in the broad-line region, it was shown by Baldwin et al. (1995) and Korista et al. (1996), that each emission line is most efficiently produced in gas with the optimum parameters for that line. These are the so called locally optimally emitting clouds or LOCs. Thus our modeling derives the parameters of the LOCs for each line so that the line fluxes and ratios provide a good approximation to a detailed multi-zone model of the BLR (see Baldwin et al 1995, Korista et al. 1996 and our Table 3), which is beyond the scope of this paper. b) The weak MgII problem The MgII line is surrounded on both sides by FeII emission. As the FeII emission in NLSy1s is usually very strong (Boller, Brandt & Fink 1996), it is possible that the wings of MgII disappear in the stronger iron bumps. This effect could lead to an underestimation of the MgII emission of up to a factor of two in the strongest FeII sources (such as IZw1 - see Vestergaard and Wilkes 2000). However, even if our measurements were underestimating the MgII emission by such a large factor, the real MgII/Ly$`\alpha `$ ratio would still be much smaller than that observed in quasars or in the lower end of the range for Seyfert 1s (see Table 3). The ionization parameter U inferred from the observed line ratio would be $`10^110^2`$ (see Fig. 4e), still larger than that inferred from the other emission lines. Thus we conclude that the MgII emitting clouds have a different value of the ionization parameter, and are formed in a physically different region of the BLR (consistent with the different time lags shown by Ly$`\alpha `$, CIV and MgII lines in NGC 5548). Photoionization models predict that MgII is either formed in a partially ionized zone (PIZ) of the BLR clouds or in a low ionization region (LIL) separate from the high ionization region (HIL) where the Ly$`\alpha `$, CIV, CIII\] lines are formed (Collin-Souffrin et al. 1988). If the MgII line is formed in a PIZ it is possible that the stronger BBB in NLSy1s will push the ionization front further back into the cloud, resulting in a smaller PIZ and weaker MgII emission, than for objects with a “normal” BBB. However, the weaker MgII emission is not consistent with the stronger FeII optical emission observed in NLSy1, which in photoionization models is predicted to be formed in the same region (PIZ, Krolik & Kallman 1988, LIL, Collin-Souffrin et al. 1988). This inconsistency suggests that the FeII emission is instead generated in a different region from MgII. The observations of line variability in NGC 5548 (Sergeev et al. 1997) showed that the FeII optical multipets have a very long time lag of several hundred days, while the MgII has a 30-50 day time lag, also implying that these lines are formed in different regions. The FeII lines may be produced in the outer regions of the accretion disk as suggested by Dumont & Collin-Souffrin (1990) or in a separate, mechanically heated region closely related to the compact radio source as in Collin-Souffrin, Hameury & Jolly (1988) (hence the observed anti-correlation of FeII emission and radio flux). ### 3.2 High luminosity to the Eddington luminosity ratio It has been suggested by a number of authors that NLSy1 galaxies as a class have systematically higher ratios of their luminosity to the Eddington luminosity, i.e. they have systematically lower masses in a given luminosity range than Sy1 galaxies and QSOs (Pounds, Done & Osborne 1995, Wandel 1997). This suggestion was made based on the analogy with the Galactic black hole candidates. We will address this suggestion now. #### 3.2.1 Continuum properties One of the current explanations of the soft-X-ray excess in AGN is reprocessing of the hard X-rays by partially ionized, optically thick matter, probably in the accretion disk. The model describes well the soft X-ray continuum of low-luminosity, flat $`\alpha _{softX}`$ Seyfert galaxies, but has problems with fitting the steepest $`\alpha _{softX}`$ spectra (see Fiore, Matt & Nicastro 1997), which characterize NLSy1. The steep $`\alpha _{softX}`$ can instead be explained by emission from the innermost part of an accretion disk which is then Comptonized by an optically thin, hot corona surrounding the disk (Czerny & Elvis 1987; Laor et al. 1997). Theoretical models which can explain both the presence of the BBB and the hard X-ray emission are based either on radial or horizontal stratification between the hot optically thin and cold, optically thick accretion flow (Wandel & Urry 1991, Shapiro, Lightman & Eardley 1976, for a review see Wandel & Liang 1991). In this paper we use the model of an accretion disk corona (ADC) by Witt, Czerny & Życki (1997), where the corona itself accretes and generates energy through viscosity, and the division of the flow into optically thin and optically thick regions results from the cooling instability discussed by Krolik, McKee and Tarter (1981). Such a model is able to predict the fraction of the energy generated in the corona instead of adopting this quantity as a free parameter. The model is fully defined by 3 parameters: the mass of the central black hole ($`M_{bh}`$), the accretion rate or the ratio of the luminosity to the Eddington luminosity ($`L/L_{Edd}`$) and the viscosity parameter ($`\alpha _{vis}`$, assumed to be the same in both the disk and the corona). The model predicts a systematic change in the opt/UV/X-ray spectral energy distribution due to a change in $`L/L_{Edd}`$. A larger ratio results in a more pronounced BBB, which is shifted towards higher energies (resulting in stronger soft-X-ray emission and hence steeper soft X-ray slopes). We have determined continuum properties predicted by this model over a large range of $`L/L_{Edd}`$ (0.001 to 0.7), $`\alpha _{vis}`$ (0.02 to 0.4) and black hole masses ($`10^6`$ to $`10^{10}M_{}`$). We then compared the observed continua of our NLSy1 with the UV luminosity at 2500Å, and the soft and hard X-ray slopes ($`\alpha _{softX}`$ and $`\alpha _{hardX}`$) predicted by the model. Table 5 shows the observed $`\alpha _{softX}`$ (from ROSAT) and $`\alpha _{hardX}`$ (from ASCA) slopes for each object, while Table 6 gives the best fitted model parameters for each object. Our model was able to reproduce the steep soft and hard X-ray slopes within the observed uncertainties for most of the NLSy1. However for two objects (IZw1, PKS 0558$``$504) we did not succeed in fitting both the soft and hard-X-ray slopes simultaneously. This may be due to the way we treat Comptonization in our model (see Janiuk & Czerny 1999 for further details). In Fig. 5a,b we show how the X-ray slopes change with the model parameters. Each curve represents one value of $`L/L_{Edd}`$ and $`\alpha _{vis}`$ and a full range of black hole masses, where smaller $`M_{bh}`$ lie at smaller 2500Å luminosities. We see clearly that only the large ratios of $`L/L_{Edd}`$ can give the steep, observed soft X-ray slopes. As has been shown by Czerny, Witt & Życki (1997), quasars radiate usually at $``$ 0.01-0.2 of their Eddington luminosity, while Seyfert galaxies radiate at $``$ 0.001-0.3. Our NLSy1 (where we use the same ADC model as Czerny, Witt & Życki 1997 to fit the parameters of the central engine) radiate at $`L/L_{Edd}`$ $``$ 0.27-0.58, much larger than the typical AGN. The masses of the central black hole calculated from the model ($`10^8M_{}`$ to $`10^9M_{}`$) for our objects are of the same order as masses found in typical Seyfert 1 galaxies, but the bolometric luminosities are larger, and comparable to those of QSOs (see Wilkes et al. 1999, Table 12 for comparison). This is deduced from the stronger, higher energy BBB and places the NLSy1 in a transition zone between the Sy1s and QSOs i.e. among Sy1s with larger luminosities or QSOs with lower masses. We note that, while the absolute numbers we deduce depend upon the particular ADC model used, the general trends do not. #### 3.2.2 Density and the radius of the BLR The structure and the dynamics of the BLR is complex, as suggested by variability studies in the case of Seyfert galaxies. However, we can analize the scaling properties of the whole BLR of an object with the properties of the central source, including the shape of the X-ray continuum. The BLR gas Compton heated by the ionizing continuum will form (in any geometry) two phases: a cool phase with $`T_c10^4`$ (the BLR clouds) and a hot phase with $`T_h10^8`$ (the intercloud medium, see Krolik and Kallman 1988, Czerny & Dumont 1998, Wandel & Liang 1991), when in equilibrium. The precise values of these temperatures depend on the shape of the continuum. In the context of the two-phase model, we will now investigate how the properties of the BLR change due to the steeper X-ray continuum of a NLSy1. We use the ionization parameter of Krolik, McKee & Tarter (1981): $$\mathrm{\Xi }=\frac{2.3F_{ion}}{cp}=\frac{2.3F_{ion}}{\frac{ck\rho _cT_c}{\mu H}}$$ (1) where $`p`$ is the total pressure, $`\rho _c`$ and $`T_c`$ the density and temperature of the cold phase, and $`F_{ion}`$ is the flux above 1 Ryd determined by the ionizing luminosity of the central source $`L_{ion}`$ and the current radius $`r`$ (where effects of geometry have been neglected): $$F_{ion}=L_{ion}/4\pi r^2.$$ (2) The two phases coexist at a value of the ionization parameter, $`\mathrm{\Xi }_h`$, which scales with the hot phase temperature, $`T_h`$ in the following way (Begelman et al. 1983): $$\mathrm{\Xi }_h=0.65\left(\frac{T_h}{10^8}\right)^{3/2}$$ (3) The BLR is most probably radially extended. For the purpose of exploring the various dependencies, we determine a representative radius for the BLR. Note that this is a scaling factor rather than the specific radius at which a particular emission line is generated. If the cloud number density profile is flatter than $`r^2`$, then most of the emission would come from the outer radii of the BLR. As in the case of the Inverse Compton heated coronae discussed by Begelman et al. (1983), a nearly hydrostatic corona will exist up to a radius where the temperature of the hot medium is equal to the “escape” temperature (i.e. the virial temperature). At larger radii the corona is heated to temperatures exceeding the escape temperature, becomes unstable and forms an outflowing wind. We therefore identify the outer edge of the BLR, $`r_{BLR}`$ with the radius where the hot medium temperature is equal to the virial temperature $$kT_h=\frac{GM_{bh}m_H}{r_{BLR}}$$ (4) The size of the BLR expressed in units of the Schwarzschild radius, $`R_{Schw}`$ is then given by: $$r_{BLR}/R_{Schw}=\frac{m_Hc^2}{2kT_h},$$ (5) so a lower value of the hot medium temperature in NLSy1 galaxies is consistent with larger values of $`r_{BLR}/R_{Schw}`$ and, consequently, lower values of the typical velocities. A similar conclusion, that the BLR radius is larger in NLSy1s, was reached by Wandel (1997) who assumed that the representative radius of the BLR is determined by the requirement to have a standard value of the ionization parameter. He then showed that the size of the BLR region is dependent not only on the luminosity of the central source, but also on the soft X-ray spectral slope. A steeper (softer) X-ray spectrum has a stronger ionizing power and hence, for a constant ionizing parameter, the BLR clouds are at larger distances from the central source, have smaller velocity dispersions and as a result form narrower emission lines. Laor et al. (1995) also reach a similar conclusion but in their picture the narrow lines in NLSy1 result purely from the lower black hole mass. In our scenario the lower black hole mass and the shape of the SED (i.e. the steeper soft-X-rays which decrease $`T_h`$) combine to produce the narrow lines. Combining equations (1)-(4) we estimate the cloud density: $$\rho _c\frac{L}{M_{bh}^2}\times \frac{T_h^{7/2}}{T_c}$$ (6) or using logarithms: $$\mathrm{log}\rho _c\mathrm{log}L2\mathrm{log}M_{bh}+7/2\mathrm{log}T_h\mathrm{log}T_c$$ (7) As has been argued in Section 3.2.1, NLSy1s have bolometric luminosities comparable to QSOs, although their central black holes have lower masses. The median value of a black hole in quasars is $``$ $`10^{10}M_{}`$ (see Czerny, Witt & Życki 1997 calculations), while the median black hole mass in NLSy1s (as inferred from our calculations, using the same ADC model - see Table 6) is $`10^{8.26}M_{}`$ i.e. $``$ 55 times lower. Let us assume that a typical quasar SED is composed of a power law and an accretion disk spectrum peaking at 10 eV ($`\mathrm{log}\nu =15.38`$, 1240Å), while a typical NLSy1 SED has a power law and a disk peaking at 80 eV (however note that the most extreme NLSy1 RE J1034$`+`$396 has its peak at 120eV \- see Puchnarewicz et al. 1995). Krolik and Kallman (1988) calculated the Compton temperatures of the hot phase for these SEDs, normalizing both to have the same total ionizing energy. The 10eV bump spectrum gave Compton temperatures $``$ 3.0 $`\times 10^7K`$ while the 80eV bump gave a lower temperature $``$ 8.0 $`\times 10^6K`$. At the same time the temperature of the cool phase increased by a factor $``$ 3.0 (0.5 in logarithm see Krolik & Kallman 1988 Fig. 2). Hence the Compton temperature of the hot phase in NLSy1 and QSOs differs by: $`\mathrm{log}T_{h,NLSy1}\mathrm{log}T_{h,QSO}=0.57`$ and the temperature of the cold phase is larger by: $`\mathrm{log}T_{c,NLSy1}\mathrm{log}T_{c,QSO}=0.5`$. Substituting the above values into equation (7) implies that $`\mathrm{log}\rho _{c,NLSy1}\mathrm{log}\rho _{c,QSO}1`$ i.e the densities of the BLR should be higher by a factor of 10 in NLSy1 than in typical QSOs with redder BBB. The larger BLR radii and larger by a factor 10 densities obtained from our modeling (as being due to hotter BBBs) are consistent with the narrow lines and line ratios observed in NLSy1s. Thus we conclude that the unusually hot and strong BBB in NLSy1s can naturally produce their observed UV spectra. ## 4 NLSy1 vs. BAL QSOs It has been suggested (e.g. Leighly et al. 1997, Lawrence et al. 1997) that there may exist a connection between NLSy1 and BAL QSOs. Both these classes have strong FeII $`\lambda `$4570 and AlIII $`\lambda `$1857 emission and weak CIV $`\lambda `$1549 and \[OIII$`]\lambda `$5007. Their continua are red in the optical and strong in the IR; additionally both classes are mostly radio-quiet. Leighly et al. (1997) reported evidence for relativistic outflows in three NLSy1. Observationally there are also many differences. NLSy1 are strong soft-X-ray emitters, while BAL QSOs are weak, possibly due to X-ray absorption (Mathur, Elvis & Singh 1995). BAL QSOs are thought to be seen more edge-on, at viewing angles skimming the edge of the dusty torus (Turnshek et al. 1996, Aldcroft, Elvis & Bechtold 1993). NLSy1s, on the other hand, are probably viewed more face-on, as they show low absorption from the torus (Boller, Brandt & Fink 1996) and some even show beaming in their radio spectra (e.g. PKS 0558-504, Remillard et al. 1991). In high resolution HST spectra NLSy1 show absorption features which are much weaker than in BAL QSOs (see Table 2). However this is expected since 50% of Seyfert 1s show absorption features (Crenshaw et al. 1995). The optical spectra of some BAL QSOs may resemble the spectra of NLSy1, showing narrow H$`\beta `$ with FWHM $`<`$ 2000 kms<sup>-1</sup> (that is why the two BAL QSOs: PG 1351+640 and PG 1411+442 were initially chosen to be in our sample), but this only cautions us that basing classifications on optical spectra alone is potentially misleading. ## 5 Conclusions In this paper we have studied the UV emission line properties of a class of extreme opt/UV/X-ray AGN: the narrow line Seyfert 1 galaxies. We found 11 NLSy1s that had been observed in the UV by either HST or IUE. We have shown that in comparison with “normal” broader line AGN, the equivalent widths of CIV and MgII are significantly smaller (NLSy1 have EW(CIV)$`<`$60 and EW(MgII)$`<`$20, normal AGN have EW(CIV)$`<`$210 and EW(MgII)$`<`$120), the EW of AlIII larger (few Å), and the UV line widths are narrower (although not as narrow as the optical H$`\beta `$ line). Also the CIII\]/Ly$`\alpha `$, CIV/Ly$`\alpha `$ and MgII/Ly$`\alpha `$ line ratios are smaller, while those of SiIII\]/CIII\], SiIV+OIV\]/CIV lines are larger. Photoionization models predict that these line ratios are formed in material with densities higher, by a factor few (less than 10) than standard BLR cloud densities, and with the ionization parameter lower by a factor 10. These parameters however predict higher MgII/Ly$`\alpha `$ ratio, in contradiction to the lower ratios observed requiring that MgII be produced in a separate region. We have fitted the SEDs of our NLSy1s to the Witt, Czerny & Życki (1997) model of an accretion disk with a Compton cooled corona and found that NLSy1s radiate at $`0.27<L/L_{Edd}<0.58`$, much larger than the typical AGN ($`L/L_{Edd}<`$ 0.3). The masses of the central black holes calculated from the model are, in our objects, of the order of masses found in typical Seyfert 1 galaxies ($`10^8M_{}`$) but the bolometric luminosities ($`\nu L_\nu 10^{46}`$ erg s<sup>-1</sup>) are larger and comparable to those of QSOs. Krolik & Kallman (1988) predict that steeper soft-X-ray BBBs, such as these of NLSy1s, change the equilibrium of the two-phase cloud-intercloud medium, decreasing the temperature of the hot intercloud medium (which we assume to be the corona above the accretion disk) and increasing the temperature of the cool BLR clouds. We show that this change in equilibrium increases the density of the BLR clouds resulting in a change of the observed line intensities and ratios consistent with these in NLSy1s. In addition the resulting decrease in $`T_h`$, causes an increase in the radius of the BLR, a correspondingly lower velocity dispersion and narrower lines as observed in NLSy1s. The NLSy1s lie at the extreme end of the Boroson and Green eigenvector 1 (Boroson & Green 1992), which was then found (Brandt & Boller 1998) to link the soft X-ray properties with the optical properties i.e. the FeII/H$`\beta `$ and $`[`$OIII\] strengths and H$`\beta `$ line width. We have found that the NLSy1s have very weak CIV and CIII\] lines, and narrow UV lines extending the set of parameters linked to eigenvector 1. The large BLR cloud densities, deduced from these characteristic UV line ratios, are probably due to the steep soft X-ray SEDs, which are in turn, the result of larger $`L/L_{Edd}`$ ratios (as inferred from the Witt, Czerny & Życki 1997 ADC model). In this scenario a larger $`L/L_{Edd}`$ is the physical parameter driving the Boroson & Green eigenvector 1. We are grateful to Niel Brandt for helping us to obtain a complete list of known NLSy1 and their X-ray slopes, Adam Dobrzycki for providing us with the HST data and Ken Lanzetta for the IUE Atlas of AGN spectra. We wish to thank Martin Gaskell, Martin Elvis, Marianne Vestergaard, Kirk Korista and Suzy Collin-Souffrin for valuable discussions and thank the anonymous referee for comments that improved the manuscript. JK greatfully acknowledges the support of a Smithsonian pre-doctoral fellowship at the Harvard-Smithsonian Center for Astrophysics and grant no. 2P03D018.16 of the Polish State Committee for Scientific Research (JK and BCz). BJW acknowledges NASA contract NAS8-39073 (Chandra X-ray Center) and SM a NASA grant NAG5-3249 (LTSA)
warning/0006/hep-lat0006030.html	ar5iv	text	# Exact local fermionic zero modes ## I Introduction Lattice fermions in interaction with background gauge fields have a tendency to produce localized fermionic approximate zero modes. Since the zero modes are localized, the influence of gauge fields far away from the center of the zero mode is small and should be well approximated by a mean-field, uniform configuration. One would guess then that what causes the zero mode is a certain local structure in the gauge configuration. In this paper we show how one can easily test whether a given local gauge configuration will have an exact zero mode at finite or infinite lattice size. We also show that one can get approximate zero modes in this way, using the variational principle. Inverting the logic, we show how one can find local gauge structures that bind fermionic (approximate) zero modes to them. There are two main classes of fermionic lattice zero modes: relatively large ones that should scale and have a physical effect in the continuum limit and very local ones, whose size shrinks to zero on continuum scales and therefore are lattice artifacts. The lattice artifact zero modes are a serious impediment to practical simulations in the quenched approximation and to the implementation of exact chirality on the lattice using the overlap Dirac operator . A better understanding of the fermionic zero modes is important in both cases: we want to see their effects clearly when these are genuine continuum effects and would like to suppress them when they are lattice artifacts. There is a difference between the roles the approximate zero modes which are lattice artifacts play in the traditional context and in the overlap context. In a traditional quenched QCD calculation the lattice artifacts can become an unsurmountable problem because their effect on fermion propagators can be spuriously large and there is no natural prescription for how to handle these non-universal effects. In the context of the overlap however, the problem is less severe because there is a well defined natural procedure to deal with these lattice artifacts. The problem is now only of a numerical nature, since the evaluation of the sign function for numerically tiny arguments is computationally costly. Strictly speaking, exact zero modes are not a major difficulty in either framework. Still, when trying to understand the origin of the approximate fermionic zero modes, it is best to start analytical work from exact ones. ## II The basic idea Let $`S`$ denote a cluster comprising of a finite local collection of sites $`s`$ and let the gauge field configuration have $`U_\mu (x)1`$ on all links outside the cluster. Let $`V_S`$ be the subspace of the space $`V`$ of fermion fields $`\psi `$ with $`\psi (x)=0`$ for $`xS`$. Let $`D_f`$ be the appropriate sparse matrix realization of the free lattice Dirac operator with uniform link variables $`U_\mu (x)1`$. Let us assume that $`G_f=D_f^1`$ exists. The total Dirac operator, $`D`$, maps $`V`$ into $`V`$ and is given by $$D=L+D_f$$ (1) $`L`$ maps the entire $`V`$ into $`V_S`$. Consider now the operator $$R=1+LG_f$$ (2) It naturally defines a restricted operator $`R_S`$ from $`V_S`$ to $`V_S`$. The equation $`D\psi =0`$ has a solution in $`V`$ if and only if there is a $`\varphi V_S`$ such that $`R_S\varphi =0`$. (At infinite volume this solution will be normalizable if the decay of $`G_f`$ is fast enough.) The main point is that $`R_S`$ is a small matrix if the cluster is small. The proof of the above assertion is easy: If $`\psi `$ is a zero mode of $`D`$ define $`\varphi =L\psi `$; clearly, $`\varphi V_S`$ and $`L(1+G_fL)\psi =0`$. If $`\varphi V_S`$ obeys $`R_S\varphi =0`$, define $`\psi =G_f\varphi `$. Since $`\varphi =LG_f\varphi `$, we also have $`L\psi =\varphi `$ implying $`D\psi =0`$. The above observations can be easily extended to the case that $`R_S^{}R_S`$ has a small eigenvalue $`\eta `$. Then, one gets a variational upper bound for the lowest eigenvalue of $`D^{}D`$. If there are several linearly independent small eigenvalue eigenstates of $`R_S`$ then we have bounds on several of the lowest eigenvalues of $`D^{}D`$. For a single state $`\varphi V_S`$ satisfying $$R_S^{}R_S\varphi =\eta \varphi $$ (3) we make the variational ansatz $$\psi =G_f\varphi $$ (4) and find $$\lambda _{\mathrm{min}}(D^{}D)\eta \frac{\varphi ^{}\varphi }{\varphi ^{}\frac{1}{D_f^{}D_f}\varphi }$$ (5) ## III “Fluxon” configuration with staggered fermions This is an example that was partially analyzed before . The gauge configuration consists of making all links that go out into the positive directions of a fixed site be $`1`$. All other links are set to unity. This configuration was called a “fluxon” in the past. It is a local minimum of the single plaquette action with gauge groups $`U(1)`$ or $`SU(2)`$ and with both a fundamental and an adjoint term. For $`SU(2)`$ it is known that there is a crossover along the Wilson line which is related to an end point of a transition line in the extended fundamental-adjoint plane . The transition across that line can be argued to be related to a condensation of fluxons. At weak coupling there is a finite and calculable density of local fluxons. We first consider naive fermions in a single fluxon background. In $`d`$ dimensions we shall find, at infinite volume, several degenerate normalizable fermionic zero modes. The decay of the zero modes is power-like, since the fermionic theory is massless. In $`d`$ dimensions the $`\gamma _\mu `$ matrices are $`2^{d_2}\times 2^{d_2}`$, where $`d_2`$ is the integer part of $`\frac{d}{2}`$. A simple calculation produces the following $`(d+1)2^{d_2}\times (d+1)2^{d_2}`$ $`R_S`$ matrix: $$R_S=\left(\begin{array}{cc}1\frac{1}{d}\gamma \gamma ^T& 0\\ 0& 0\end{array}\right)$$ (6) Here we introduced a $`2^{d_2}d\times 2^{d_2}`$ matrix $`\gamma `$ and the $`2^{d_2}\times 2^{d_2}d`$ matrix $`\gamma ^T`$: $$\gamma =\left(\begin{array}{c}\gamma _1\\ \gamma _2\\ \mathrm{}\\ \gamma _d\end{array}\right),\gamma ^T=\left(\begin{array}{cccc}\gamma _1& \gamma _2& \mathrm{}& \gamma _d\end{array}\right)$$ (7) The most evident zero mode (normalizable for $`d3`$) is given by choosing $$\varphi =\left(\begin{array}{c}0\\ 0\\ \mathrm{}\\ \chi \end{array}\right)$$ (8) Each entry above has $`2^{d_2}`$ components. We then find $`2^{d_2}`$ fold degenerate zero modes given by (up to normalization factors): $$\psi _0(x)=\frac{d^dp}{(2\pi )^d}\frac{\mathrm{sin}px}{_\mu \gamma _\mu \mathrm{sin}p_\mu }\chi $$ (9) This class of zero modes has been found previously in . Our more detailed study here shows that there are more zero modes: Choose $`\varphi ^{}`$ $$\varphi ^{}=\left(\begin{array}{c}\gamma _1\chi ^{}\\ \gamma _2\chi ^{}\\ \mathrm{}\\ \gamma _d\chi ^{}\\ 0\end{array}\right)$$ (10) It is easy to check that $$(1\frac{1}{d}\gamma \gamma ^T)\varphi ^{}=0$$ (11) The new zero modes of $`D`$ are $$\psi _0^{}(x)=i\delta _{x,0}\chi ^{}+\frac{1}{2}\frac{d^dp}{(2\pi )^d}\mathrm{sin}px\frac{1}{_\mu \mathrm{sin}^2p_\mu }\left[\underset{\mu \nu }{}\gamma _\mu \gamma _\nu \mathrm{sin}(p_\mu p_\nu )+\underset{\mu }{}\mathrm{sin}2p_\mu \right]\chi ^{}$$ (12) In four dimensions for example, we have found 8 zero modes. There are no other zero modes. Each set of four zero modes maps into itself under the action of $`\gamma _5`$. One can easily decompose each set of 4 zero modes into 2 of positive chirality and 2 of negative chirality. One naive lattice fermion should represent 16 continuum fermions. The fluxon configuration background breaks the continuum flavor symmetry because there are half as many zero modes as expected continuum fermions. Consider now staggered fermions; they are treated by block diagonalizing the naive fermion system. When we decouple the four dimensional naive fermions into four staggered fermions, each will have two zero modes. In the continuum limit each staggered fermion should reproduce four flavors. With only two zero modes for a staggered fermion it is again clear that the fluxons are nonperturbative contributors to lattice flavor symmetry breaking. ## IV “Fluxon” configuration with Wilson fermions The zero modes we found above were localized only by power decays. To get light quarks with Wilson fermions the mass parameter has to be given negative values. The free propagator is now massive, decaying exponentially. We use the same technique to look for negative mass values where fluxons can produce exact, exponentially localized, fermionic zero modes. We work in dimension $`d`$ larger or equal to 2. The matrix $`R_S`$ now becomes: $$R_S=1+\left(\begin{array}{cccc}(1+\gamma _1)G_f(0,\widehat{1})& \mathrm{}& (1+\gamma _1)G_f(0,\widehat{d})& (1+\gamma _1)G_f(0,0)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ (1+\gamma _d)G_f(0,\widehat{1})& \mathrm{}& (1+\gamma _d)G_f(0,\widehat{d})& (1+\gamma _d)G_f(0,0)\\ _\mu (1\gamma _\mu )G_f(\widehat{\mu },\widehat{1}))& \mathrm{}& _\mu (1\gamma _\mu )G_f(\widehat{\mu },\widehat{d})& _\mu (1\gamma _\mu )G_f(\widehat{\mu },0)\end{array}\right)$$ (13) To search for zero modes we calculate the determinant of $`R_S`$. It inherits reality from the determinant of $`D`$. By working out explicitly the cases $`d=2,3,4,5,6`$ we arrive at a formula for general $`d`$. We have not derived the formula for arbitrary $`d`$, so we present it as a conjecture. $$detR_S=\left[1+\frac{d1}{d}[x(m+d)1]^2+2[x(m+d)1]+d(d1)xy\right]^{2^{d_2}}$$ (14) The parameters $`x`$ and $`y`$ are defined by fermionic propagators: $`x={\displaystyle \frac{d^dp}{(2\pi )^d}\frac{b(p)}{b^2(p)+s^2(p)}}`$ , $`y={\displaystyle \frac{d^dp}{(2\pi )^d}\frac{\mathrm{cos}p_1[2\mathrm{sin}^2p_2b(p)\mathrm{cos}p_2]}{b^2(p)+s^2(p)}},`$ (15) $`b(p)=m+d{\displaystyle \underset{\mu }{}}\mathrm{cos}p_\mu `$ , $`s^2(p)={\displaystyle \underset{\mu }{}}\mathrm{sin}^2p_\mu .`$ (16) When there are zero modes, they will be $`2^{d_2}`$-fold degenerate (at least for $`d=2,3,4,5,6`$). It is a simple matter now to search for an $`m`$ for which there are zero modes. We start from $`m=0`$ and search going towards negative values. The first mass values which gives zero modes is listed below. As far as simulations go, only the region $`2<m<0`$ is of interest. For $`d=4`$ the first zero is at $`m=2.87982`$, for $`d=3`$ the first zero is at $`m=1.940553`$, and for $`d=2`$ the first zero is at $`m=.90096`$. So, in four dimensions the fluxon is not of great importance, but in two dimensions it is. Note that the fermions do not react to the fluxon as they would to an instanton in the continuum as reflected by the zero-mode degeneracy we find. ## V A small two dimensional instanton We now turn to analyze a small two dimensional instanton. It is described by four plaquettes making up a square, with each plaquette carrying $`\frac{\pi }{2}`$ units of angular flux. The calculations are now more involved, but the bottom line is that we find a non-degenerate zero at $`m=0.39182`$, which certainly is of interest. For this value of $`m`$ the free propagator decays quite fast so the zero mode is quite local. Effects of boundary conditions can be easily investigated by appropriately adjusting $`D_f`$ and the gauge background. They are more sizable with periodic boundary conditions than with antiperiodic boundary conditions. We have also looked at the dependence on the toron coordinates (zero Fourier modes of the gauge fields $`A_\mu (x)`$ where the link variables are $`U_\mu (x)=e^{iA_\mu (x)})`$. Boundary conditions are implemented by choosing specific toron configurations. Finite volume calculations require the replacement of the momentum integrals by the appropriate sums. ## VI Finding gauge configurations that bind zero modes Choosing a cluster of sites of certain shape we now view the link variables residing on the links in the cluster as unknowns. For each such gauge field configuration we numerically compute the smallest eigenvalue of $`R_S^{}R_S`$, $`\lambda _{\mathrm{min}}(R_S^{}R_S)`$, using Wilson fermions at some negative value of the mass parameter $`m`$. Varying the gauge background, and proceeding by steepest descent we can find local minima of $`\lambda _{\mathrm{min}}(R_S^{}R_S)`$ as a function of the link variables. Often, we end up finding the minimum to be zero and then the corresponding gauge field configuration binds a localized exact zero mode. This search for a gauge configuration can be repeated for different starting gauge configurations and different values of the mass parameter. To speed the numerical procedure up we took the free propagators appearing in $`R_S`$ on a finite periodic lattice of size $`8^d`$. Restricting ourselves to $`U(1)`$ fields, in 3D, on a cluster comprising of four sites that can support a fluxon, we find zero modes only for $`m<1.9`$. The minimizing configuration is not a fluxon and it depends upon the gauge configuration we start our steepest descent search from. Nevertheless, always, the Wilson plaquette action of the minimizing configuration is close to the one of a fluxon. As long as we restrict our gauge configurations to having nontrivial links only along the links contained in the cluster, different link configurations are gauge inequivalent. Thus, there are many gauge inequivalent configurations that can produce fermionic zero modes. In four dimensions the situation is similar for the fluxon type cluster but now we find zero modes only for $`m<2.5`$, which is outside of the range of immediate interest. Our next step is to see what happens when we increase the cluster size. We consider now larger clusters that make up a cube in 3D or a hypercube in 4D. Keeping the gauge fields still in $`U(1)`$ we find that now we can get zero modes at masses closer to zero. In 3D zero modes are occurring for $`m<1.1`$ and in 4D for $`m<1.7`$. Thus, a moderate increase in cluster size increases the chances to produce zero modes at Wilson masses relevant to $`QCD`$ simulations. Finally we also look at what happens when one goes from the abelian $`U(1)`$ gauge group to the non-abelian $`SU(2)`$ gauge group. We find that making the non-abelian gauge group always allows certain gauge configuration to bind zero modes for Wilson masses closer to zero than in the abelian case. Thus, increasing the size of the group also increases the chances to produce zero modes at practically relevant values of the Wilson mass parameter $`m`$. The results of our numerical search are summarized in Fig. 1. We plot there on logarithmic scale the minimum of the smallest eigenvalue of $`R_S^{}R_S`$ as a function of the Wilson mass parameter $`m`$, for four different clusters. The 3-link and 4-link clusters in 3D and 4D mentioned in the figure contain 4 sites and 5 sites respectively and can accommodate the fluxon configurations we studied analytically before. The other two clusters we plot results for are the cube in 3D and the hypercube in 4D. For example, we see in the figure that we can get zero modes for $`m<1.1`$ in 4D if we consider a cluster that makes up a hypercube. Note the sudden decrease in $`\lambda _{\mathrm{min}}(R_S^{}R_S)`$ as a function of $`m`$; it reflects a major change in the minimizing gauge configuration. Thus, for this type of cluster, $`\lambda _{\mathrm{min}}(R_S^{}R_S)`$ has a non-analytic dependence on the gauge background at some value of $`m`$ located inside the range of interest $`2<m<0`$. Whatever the ultimate role of this effect is in a full quenched simulation of QCD is, it cannot be helpful to the hope of recovering a limit that can be described by a local effective Lagrangian of some kind. ## VII How to avoid zero modes We found that a cluster of SU(2) gauge fields as small as a hypercube can give rise to zero modes for Wilson fermions at masses close enough to zero to be relevant to practical simulations. These simulations can be with traditional fermions, or with overlap fermions. In both cases the presence of fermionic zero modes is a problem for quenched simulations; this problem has a practical solution in the overlap case, but they are still costly to handle because they need to be individually identified and projected out . Our construction directly and explicitly shows that the number of fermionic zero modes is proportional to the lattice volume, because the wave functions are explicitly known to decay exponentially, and the clusters are small and can be well separated. The rate of decay of the zero mode wave functions is given by the distance of $`m`$ from the nearest even non-negative integer. The pure gauge action for the clusters is exponentially suppressed and it is possible that the density of the associated zero modes be so small that in a practical simulation volume there would be a small chance to even produce one such mode. On the other hand, we know that small modes in numbers proportional to the lattice volume do occur for quenched gauge field configurations. To be sure, we have not established that the particular clusters we investigate here are responsible for all the small modes one sees in typical QCD simulations. But, in principle, we have established that there is a finite density of fermionic zero modes, and that this density is exponentially suppressed in the gauge coupling constant, consistent with numerical findings . Clearly, in four dimensions, a crucial question is how this suppression relates to the standard asymptotic freedom formula for a density of dimension four. We hope to come back to this point in the future. It would be nice to be able to avoid these unphysical zero modes. One can do that by changing the pure gauge action so that the density of the zero modes be further reduced. This has been argued to work to some extent . Another way would be, in accordance with the exact bound , to make the local field strength close to zero. This could be achieved, for example, by replacing all links by APE smeared ones . However, this decouples the fermions from higher momentum modes of the gauge field and makes it somewhat unclear whether we are still dealing with a single scale problem, rather than a two scale problem, where the fermions couple to gauge fields only at the lower scale. Another possibility, suggested in , is to APE smear only the links that enter into the Wilson mass term, but not those that are coupled to the Dirac $`\gamma `$-matrices. In practice, implementing this last alternative would roughly quadruple the time it takes to calculate the action of the Wilson Dirac operator on a fermion field, because the projector structure of the $`1\pm \gamma _\mu `$ terms is lost and the fermions now interact with two kinds of gauge fields.We thank Urs Heller for correcting our original erroneous statement at this point. With our methods it is easy to obtain a rough estimate of the effect on the zero modes this latter possibility will have. We simply replace the links entering the Wilson mass term by unity, leaving the links entering the other terms in the Wilson Dirac operator free to change. This clearly is inconsistent in terms of the entire gauge background, but is very easy to implement. We find that this replacement of the Wilson mass term chases the fermionic zero modes away from negative mass values too close to zero. This is an encouraging sign, and the idea deserves further investigation. ## VIII Conclusions In this paper we have devised a simple way to look for approximate and exact fermionic zero modes on the lattice and for the gauge configurations that bind them. The fermionic zero modes are bound to local gauge configurations and are exponentially localized. In a real quenched simulation they will come at a finite density making the quenched approximation in principle unusable when traditional fermions and a traditional gauge action are employed and one wants realistically light quarks to emerge in the continuum. Our method makes it obvious that such fermion zero modes will occur with a finite probability per unit Euclidean four volume. We hope to use this technique in the reverse, namely, as a method to efficiently search for ways to ameliorate the practical problems posed by fermionic almost zero modes in the overlap context. ###### Acknowledgements. This research was supported in part by the DOE under grant # DE-FG05-96ER40559. F. B. thanks P. Sodano for discussions at earlier stages of the project and gratefully acknowledges the Fondazione Angelo Della Riccia for financial support.

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