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# Interpreting the optical data of the Hubble Deep Field South: colors, morphological number counts and photometric redshifts
## 1 Introduction
The Hubble Deep Field South (HDF-S) was observed in October 1998 by the Hubble Space Telescope (HST). It is the southern counterpart of the Hubble Deep Field North (HDF-N) and shares its characteristics of depth and spatial resolution. The HDF-S is a four arcmin<sup>2</sup> survey located at RA 22$`h`$ 32$`m`$ 56$`s`$, DEC -60 33’ 02”, observed during 150 orbits with the Wide Field Planetary Camera 2 (WFPC2). HDF-S images cover a wavelength range from the ultraviolet to near-infrared with 4 broad-band filters: F300W, F450W, F606W, F814W. Details about observations and data reduction may be found in the Hubble Deep Field South web page (http://www.stsci.edu/ ftp/science/hdfsouth/hdfs.html).
The main goal of these surveys, HDF-N and HDF-S, is to study faint galaxies, in particular to investigate the nature of the faint blue galaxies excess and the evolution of galaxies (e.g. Ellis 1997).
Faint galaxy counts have been of paramount importance to show that galaxies do evolve with redshift, even if the overall scenario and the physical processes that led to evolution are still debated. An excess of observed galaxies with respect to a no-evolution model begins to appear at B$`2122`$ (e.g. Maddox et al. 1990; Jones et al. 1991; Metcalfe et al. 1991) and continues to rise up to B$`>28`$ (Metcalfe et al. 1995) exceeding by a factor 4 to 10, depending on the choice of $`q_0`$ (0-0.5) and on the LF (see e.g. Koo & Kron 1992; Ellis 1997). This excess has actually undoubtly proved that galaxies evolve. However, the analysis of the optical data have not offered a unequivocal interpretation: the steepening of the faint end luminosity function (LF) of galaxies at high redshifts (Lilly et al. 1995; Ellis et al. 1996) suggests that either the typical luminosity of galaxies was brighter by $`1`$ magnitude at $`z0.7`$ or that there were 2-3 times more galaxies at this epoch. Most likely both luminosity and number density evolution are acting: luminosity evolution seems to play an important role at least on disk galaxies as shown by the results obtained by Schade et al. (1995) and Lilly et al. (1998) on the CFRS-LDSS sample. On the other hand also mergers of galaxies seem to significantly contribute to the evolution of both the luminosity function and luminosity density of the Universe at least down to $`z1`$ (Le Fèvre et al. 2000). Analyzing high resolution ground based imaging of blue galaxies at redshift $`z<0.7`$, Colless et al. (1994) actually noticed that most of them have a close companion at a projected distance of about 10 kpc. This result suggests the presence of merging at $`z<0.5`$, which could increase the star formation, explaining the blue colors of the sources.
Disentangling the effects of luminosity and number density evolution is fundamental in understanding the nature of field galaxy formation and evolution, allowing a direct comparison with models. For instance, semi-analytical models of galaxy formation in hierarchical clustering scenarios, assuming that galaxies assemble through successive merger events from smaller sub-units, make concrete predictions about the evolution of merger rate, morphological mix and relative redshift distribution (e.g. Baugh et al. 1996; Kauffmann 1996; Kauffmann & Charlot 1998). In this scenario the number density of galaxies should increase with redshift, their sizes should decrease and the fraction of high redshift galaxies ($`z>2`$) in an IR-selected sample should be very low. Moreover the morphological mix should change in favour of late type and irregular galaxies which thus would dominate both galaxy counts at faint magnitudes and the high redshift tail of the redshift distribution, where a lack of elliptical galaxies should be observed. Most of these predictions have been the object of detailed analysis of data resulting in conflicting results. For instance, some works claim a deficit of ellipticals at $`z>1`$ (Zepf 1997; Franceschini et al. 1998; Barger et al. 1999) and a very low fraction ($`5\%`$) of $`z>2`$ galaxies in a K$`<21`$ selected sample (Fontana et al. 1999). On the contrary other authors find a constant comoving density of ellipticals down to $`z2`$ (e.g. Benítez et al. 1999; Broadhurst & Bowens 2000) and a strong clustering signal of EROs in K-selected samples (Daddi et al. 2000). HST Medium Deep Survey (MDS), on the other hand, indicates that a large fraction of galaxies at moderate redshift ($`z0.5`$) are actually intrinsically faint and exhibit peculiar morphologies suggestive of merging (Griffiths et al. 1994; Driver et al. 1995; Glazebrook et al. 1995; Neuschaefer et al. 1997). The results derived from the morphological analysis on the HDF-N data (van den Bergh et al. 1996; Abraham et al. 1996; Driver et al. 1998) seem to confirm the high fraction of apparent irregular/peculiar galaxies with respect to the local universe which may be the cause of the faint blue galaxies excess (Ellis 1997).
Williams et al. (1996) stress also that HDF-N galaxies show a strong bluing color trend: while red galaxies are compatible with an elliptical galaxy colors, most of blue sources are bluer than an irregular galaxy, and these colors might be accounted for only considering luminosity evolution.
Thus it seems that the overall evolutionary scenario is still missed and the understanding of the contribution of the different populations of galaxies to this scenario rather confused.
In this work we present an analysis of number counts, colors and morphological distributions of HDF-S galaxies in an attempt to reveal and better constraint the evolving populations of galaxies. The plan of the paper is as follows: in Section 2 we summarize the steps leading to the creation of the catalogue, the computation of optical galaxy counts is explained in Section 3, in Section 4 the technique used to determine galaxy colors. In Section 5 we discuss our results, both counts and colors, while in Section 6 we present and discuss morphological number counts and in Section 7, by adding the information given by photometric redshifts and number counts models, we disentangle the contribution of different galaxy populations and discuss the interpretation of redshift distributions. Finally, in Section 8 we summarize our results and express our conclusions.
## 2 Image Analysis
The samples here used are not extracted from the public HDF-S catalogs available on the Web. On the contrary they have been independently created beginning from the optical HST images in order to have a control of the detection procedure and of magnitude estimate. A detailed description of the procedure adopted and the selection criteria used to obtain the catalogs in the different bands, as well as magnitude estimates, can be found in Volonteri et al. (2000). Here we briefly describe the steps related to the derivation of number counts.
The object detection was performed across the area observed in optical bands (WFPC2 field) using SExtractor (Bertin and Arnouts 1996). Images were first smoothed with a Gaussian function having FWHM equal to the one measured on the images ($``$ 0.2 arcsec). A detection threshold of 1 sigma per pixel and a minimum detection area equal to the seeing disk ($`0.02`$ arcsec$`{}_{}{}^{2}=13`$ pixel) were adopted to peak up objects. This threshold corresponds to a minimum signal-to-noise ratio of $`S/N_{WF}=1.34`$ and $`S/N_{PC}=0.67`$ for the faintest sources detectable on the WF area and on the PC area respectively.
In Table 1 we report for each filter the zero-point (AB magnitude, Oke 1974), the sky RMS estimated by SExtractor and the corresponding 5$`\sigma `$ limiting magnitude for a point source.
“Pseudo-total” magnitudes were estimated using the method of Djorgovski et al. (1995) and Smail et al. (1995) on the basis of the following steps:
* the SExtractor isophotal corrected magnitude has been assigned to *large* sources, i.e. those sources having an isophotal diameter $`D_{iso}>\theta _1`$ and to those sources flagged by SExtractor as “blended”;
* the aperture corrected magnitude (estimated within $`\theta _1`$ and then corrected to $`\theta _2`$, being $`\theta _1<\theta _2`$) has been assigned to *small* sources with $`D_{iso}<\theta _1`$;
$`\theta _1`$ has been defined as the minimum apparent diameter of a galaxy having an effective diameter $`r_e=10`$ kpc. Hereafter we use a $`\mathrm{\Lambda }=0`$ cosmology, with $`q_0=0.5`$ and $`H_0=50`$ km s<sup>-1</sup>Mpc<sup>-1</sup> unless differently specified. With this choice $`\theta _1=1.2`$ arcsec. $`\theta _2`$ is the diameter corresponding to the area for which 90$`\%`$ of the bright sources have a smaller isophotal area (see Volonteri et al. 2000 for a complete description of procedure, tests and results). This procedure was applied independently for each band.
A measure of the detection reliability is necessary in order to evaluate the number of spurious sources included in the sample. We treated this problem statistically, in the hypothesis that noise is symmetrical with respect to the mean sky value. Operationally we first created for each filter a noise frame by reversing the original images to reveal the negative fluctuations and to make negative (i.e. undetectable) real sources (Saracco et al. 1999). Then we run SExtractor with the same detection parameter set used to search for sources in the original images detecting, by definition, only spurious sources. In Figure 1 the magnitude distribution of spurious sources obtained on the WF area and the PC area in the F606W band are shown respectively. Applying a S/N=5 cut off is sufficient to reduce the spurious contamination to a negligible fraction (4$`\%`$) on the WF area, after removing the edges of the images with lower sensitivity. On the contrary such a cut off is not able to reduce spurious detections to a reasonable level on the PC area, contamination being more than 35$`\%`$. Thus due to such a large number of spurious sources which would be introduced by the PC data, we restricted the selection of sources to the central WF area only corresponding to 4.38 arcmin<sup>2</sup>. On this area 450, 1153, 1694 and 1416 sources were selected accordingly to the above criteria in the F300W, F450W, F606W and F814W band respectively.
## 3 Optical Galaxy Counts
To derive galaxy counts, we first removed stars from the sample by using the SExtractor star$`/`$galaxy classifier. We defined as stars those sources brighter than F814W$`{}_{AB}{}^{}=22`$ and having a value of the SExtractor “stellarity” index larger than 0.9. This choice will tend to underestimate stars both at faint magnitudes where no classification is considered, and at bright magnitudes where some fuzzy stars could be misclassified as galaxies. On the other hand this will ensure that our galaxy sample is not biased against compact galaxies. The star “cleaning” procedure has classified and removed 14 stars with F814W$`{}_{AB}{}^{}<22`$ in agreement with the number of stars found in the HDF-N by Mendez et al. (1998) to this depth and in excess by a factor of two with respect to the prediction of the galaxy model of Bahcall & Soneira (1981).
Completeness correction for faint undetected sources strongly depends on the source apparent spatial structure and on their magnitude. The high resolution of HDF-S images allows the detection of sub-galactic structures, such as HII regions. Moreover the light observed from a high fraction of the galaxies in the HDF-S is emitted in the UV and F450W pass-bands, so that the observed morphology is strongly affected by star formation episodes. These features imply that “typical” profiles of galaxies are not able to well describe the shapes of a lot of galaxies in the HDF-S. Thus, in order to reproduce the manifold of shapes which characterizes sources in the HDF-S, following Saracco et al. (1999), we generated a set of simulated frames by directly dimming the original frames by various factors while keeping constant the RMS. This procedure has allowed us to avoid any assumption on the source profile providing an artificial fair dimmed sample in a real background noise even if it cannot take into account differences due to evolution. We thus define the correction factor $`\overline{c}`$ as the mean number of dimmed galaxies which should enter the fainter magnitude bin over the mean number of detected ones. In Table. 2 we report the raw counts $`n_r`$, the completeness correction factor $`\overline{c}`$, the counts per square degree corrected for incompleteness $`N`$ and their errors $`\sigma _N`$.
Colley et al. (1996) suggested that galaxies in the HDF-N may suffer from a wrong selection. High redshift galaxies on optical images have a lumpy appearance: first the redshift moves the ultraviolet rest-frame light into the optical, so galaxies are observed in UV rest-frame, where star-forming regions are more prominent; second the fraction of irregular galaxies is higher than locally (van den Bergh et al. 1996, Abraham et al. 1996) and a large number of galaxies may display asymmetry and multiple structure. We treated this feature analyzing our sample in F814W and F450W. The F814W-band catalogue should suffer less from the effects described above, being selected in the reddest filter, the vice versa is true for the F450W-band catalogue (but see Volonteri et al. 2000 for a more detailed analysis). About 20-30$`\%`$ of sources in F450W-band catalogue have separation $`<1`$ arcsec. We therefore analyzed these sources, by cross-correlating F814W and F450W catalogues. In the F450W-band catalogue we selected pairs with a separation $`<1`$ arcsec which were not included in the F814W-band catalogue. These objects were single sources split in the F450W-band (with F450W$`{}_{AB}{}^{}`$ 27-29), corresponding to a single detection in the F814W-band. We then used SExtractor on the F450W frame, after choosing a higher deblend\_mincont$`=0.1`$.
87 sources, with $`21<`$F450W$`{}_{AB}{}^{}<26`$, corresponding to about 7$`\%`$ of the whole sample, were then considered as single galaxies. In Figure 2 counts obtained with the uncorrected catalogue are shown as stars, while corrected counts are shown as empty circles. In the faintest bins the correction is within the error, accounting for less than 10$`\%`$ of the counts.
Errors were obtained by quadratically summing the Poissonian contribution $`\sigma _{n_r}=\sqrt{n_r}`$ of raw counts, the contribution due to clustering fluctuation
$$\sigma _\omega \omega (\theta )^{1/2}N$$
(1)
being $`\omega (\theta )=A_\omega \theta ^{(\gamma 1)}`$ the angular correlation function, and the uncertainty on the correction factor
$$\sigma _{\overline{c}}=\sqrt{1/k(c_i\overline{c})^2}$$
(2)
where $`k`$ is the number of frames dimmed for each dimming factor. Assuming that the amplitude $`A_\omega `$ evolves with magnitude on the basis of the relation log$`A_\omega =0.3\times mag+const`$ (Brainerd et al. 1994; Roche et al. 1996) we derived
$$\mathrm{log}A_\omega (F300W_{AB})=4.4380.3F300W_{AB}$$
$$\mathrm{log}A_\omega (F450W_{AB})=4.3140.3F450W_{AB}$$
$$\mathrm{log}A_\omega (F606W_{AB})=4.4380.3F606W_{AB}$$
$$\mathrm{log}A_\omega (F814W_{AB})=4.1610.3F814W_{AB}.$$
The number counts here derived in the F300W and F450W bands and in the F606W and F814W bands are shown in Figures 3 and 4 together with those from the literature. The relation between counts and magnitude may be written as
$$\frac{d\mathrm{log}N}{dm}=\{\begin{array}{cc}\gamma _1m+c_1\hfill & mm^{}\hfill \\ \gamma _2m+c_2\hfill & m>m^{}\hfill \end{array}$$
with $`\gamma _1>\gamma _2`$. The knee $`m^{}`$ is at B$``$25 (Lilly et al. 1991, Metcalfe et al. 1995), and the value of $`\gamma _1`$ varies between 0.4 and 0.6 according to the band and analogously for $`\gamma _2`$ it varies between 0.2 and 0.5. We estimated both $`\gamma _1`$ and $`\gamma _2`$ in F450W<sub>AB</sub>, F606W<sub>AB</sub> and F814W<sub>AB</sub>, obtaining for $`\gamma _1`$: $`\gamma _{1,F450W_{AB}}0.4\pm 0.1`$, $`\gamma _{1,F606W_{AB}}0.34\pm 0.1`$ and $`\gamma _{1,F814W_{AB}}0.62\pm 0.1`$, and for $`\gamma _2`$: $`\gamma _{2,F450W_{AB}}0.19\pm 0.01`$, $`\gamma _{2,F606W_{AB}}0.19\pm 0.1`$ and $`\gamma _{2,F814W_{AB}}0.19\pm 0.1`$.
The most notable feature appears in the F300W-band counts: these are described by a slope $`\gamma _{F300W_{AB}}=0.47\pm 0.05`$ which is much steeper than the value $`\gamma _{F300W_{AB}}0.15`$ derived by Williams et al. (1996) and Pozzetti et al. (1998) on the HDF-N data in the same magnitude range. Such a steep slope, which agrees with the findings of Hogg et al. ($`\gamma _U0.5`$; 1997) and Fontana et al. ($`\gamma _U0.49`$; 1999), does not depend on a possible over-estimate of the incompleteness in the faintest magnitude bins: the same slope is actually described by counts at magnitudes U$`{}_{300}{}^{}<26`$ where the sample is 100$`\%`$ complete. Moreover the F300W-band counts do not show evidence of any turnover or flattening down to F300W<sub>AB</sub>=27, contrary to what is claimed by Pozzetti et al. (1998). We will discuss this issue in Section 5.1. The amplitude and slope of the number counts in the whole range of magnitudes in the other optical bands are in a good agreement with those previously derived by other authors. We estimated $`\gamma _{F450W_{AB}}0.35\pm 0.02`$, $`\gamma _{F606W_{AB}}0.28\pm 0.01`$ and $`\gamma _{F814W_{AB}}0.28\pm 0.01`$ for the F450W<sub>AB</sub>, F606W<sub>AB</sub> and F814W<sub>AB</sub> counts respectively.
## 4 Optical Galaxy Colors
In order to measure colors unbiased with respect to the selection band, we created a *metaimage* by summing all four frames, after normalizing them to the same rms sky noise (see also Volonteri et al. 2000). We run SExtractor in the so-called double image mode: detection and isophote boundaries were measured on the combined image, while isophotal magnitudes were measured on F300W, F450W, F606W, F814W images individually. Using this procedure (Moustakas et al. 1997) both object detection and isophote determination are based on the summed image, and isophotes are not biased towards any of the bands.
Then we cross-correlated the catalogue obtained from the combined image with F814W and F300W samples selected according to our criteria (see Section 2): we ended up by having two sample selections. Starting from each sample, we assigned a lower limit in magnitude to sources undetected in any of other bands. The limiting magnitude is the 5$`\sigma `$ isophotal magnitude within the isophote measured in the combined image. In Figures 5, 6, 7, 8 the color-magnitude diagrams of the two samples are shown together with the median locus. The error bars are the standard deviation from the mean of the values in each bin. The most notable feature is the initial trend towards blue colors followed by a flattening in the last bins.
We then compared our colors with colors measured from spectra by Coleman, Wu & Weedman (1980, CWW hereafter) after convolution with throughput curves for the WFPC2 filters. In each magnitude bin, we estimated the fraction of sources with (F450W$``$F606W)<sub>AB</sub> bluer than an irregular galaxy at $`z=0`$ (i.e. F450W$`{}_{AB}{}^{}`$ F606W$`{}_{AB}{}^{}<`$0.18). As shown in Figure 9, at F814W$`{}_{AB}{}^{}>27`$ about $`50\%`$ of sources have (F450W$``$F606W)<sub>AB</sub> bluer than a typical local irregular. For the F300W-selected catalogue we also estimated the percentage of galaxies with (F300W$``$F450W)<sub>AB</sub> bluer than a typical irregular: at F300W<sub>AB</sub>=26 this percentage is almost $`80\%`$.
We will discuss these features in Section 5.2.
## 5 Discussion
### 5.1 Number Counts
Our results show a decreasing slope at redder wavelengths for the whole sample, in good agreement with other works. The very faint end is flatter than 0.2 in F450W<sub>AB</sub>, F606W<sub>AB</sub>, F814W<sub>AB</sub>, while in F300W<sub>AB</sub> the slope is steeper. The difference with Pozzetti et al. F300W<sub>AB</sub> counts is likely due to a different selection of the catalogue: they selected the sample in a red band (on a summed F606W+F814W image) to derive F300W-band counts and used the isophotal magnitude instead of a *pseudo-total* magnitude. They assume that galaxies fainter than F606W$`{}_{AB}{}^{}<26`$ have approximately a (F300W$``$F606W)<sub>AB</sub> color not bluer than the brighter galaxies, not biasing the red-selected sample against very blue objects.
In Figure 10 we show the (F300W$``$F606W)<sub>AB</sub> color as a function of F300W<sub>AB</sub> isophotal magnitude (as Pozzetti et al.) for our F606W-selected sample, within the limiting F606W<sub>AB</sub> (i.e. $`30.2`$) and for the F300W-selected one (Figure 11). The colors for the F606W-selected sample are obtained as described in Section 4. In order to obtain a sample with characteristics similar to Pozzetti et al. sample, we considered for F300W<sub>AB</sub> magnitude the flux within the isophote in the *metaimage*.
The solid oblique line represents the (F300W $``$F606W)<sub>AB</sub> color of a F606W$`{}_{AB}{}^{}=29`$ galaxy (the F606W-band limiting magnitude of Pozzetti et al.): all galaxies in the hatched area would be not counted since they are fainter than F606W$`{}_{AB}{}^{}>29`$, the vertical line is the F300W<sub>AB</sub>=29.5 limit. The F606W<sub>AB</sub> limiting magnitude biases the sample against blue objects, that is the last bins are affected by censored data (see next Section), thus making the (F300W$``$F606W)<sub>AB</sub> color appear redder. Considering the bluing trend observed until F300W<sub>AB</sub>=29, and the fact that the F606W<sub>AB</sub> limiting magnitude biases the sample against blue objects, this implies that F300W-band counts recovered from a red-selected sample are affected by this color bias and consequently by incompleteness in the last bins. The flattening found by other Authors may be mainly caused by the influence of such incompleteness (noise being equally important for faint galaxies, be them at low or high redshift) and not to the “crossing” of the Lyman break.
### 5.2 Colors
As explained in Section 4 we selected two samples: the F814W-selected one should trace the global features of the HDF-S sample and the F300W-selected sample evidencing star-forming galaxies at moderate $`z`$.
The F814W-selected catalogue has by far the biggest number of sources, so that the number of galaxies with lower limits in the other bands is rather high. We applied a linear fit to the color-magnitude relation by making use of “survival analysis” (Avni et al. 1980, Isobe et al. 1986), in order to avoid uncertainties introduced by the heavy influence of censored data. We applied a linear regression based on Kaplan-Meier residuals. This analysis has suggested an almost flat relation for (F450W$``$F814W)<sub>AB</sub> vs F814W<sub>AB</sub>, (F450W$``$F606W)<sub>AB</sub> vs F814W<sub>AB</sub> and (F606W$``$F814W)<sub>AB</sub> vs F814W<sub>AB</sub> for F814W$`{}_{AB}{}^{}>24`$, while (F300W$``$F450W)<sub>AB</sub> vs F814W<sub>AB</sub> did not reach convergence. Smail et al. (1995) noticed a similar tendency on their sample limited at R=27. After an initial bluing the median V$``$R color gets redder, V$``$I gets flat, while R$``$I keeps on following a bluing trend.
This trend may suggest the presence of a flat spectrum population, whose colors seem to saturate but whose contribution is more and more important at faint magnitudes, as confirmed by the rising fraction of galaxies with (F450W$``$F606W)<sub>AB</sub> bluer than a typical irregular galaxy. Broadly speaking galaxy colors are dominated by blue sources, as noted by Williams et al. (1996)in the HDF-N, that is the median color is always bluer than typical local samples. The F300W-selected sample shows a very blue (F300W$``$F450W)<sub>AB</sub> color, though the relation (F300W$``$F450W)<sub>AB</sub> vs F300W<sub>AB</sub> seems almost flat. These two features suggest a considerable contribution by flat spectrum sources. Comparing the median (F300W $``$F450W)<sub>AB</sub> color with Bruzual & Charlot (1993) and CWW predicted colors, the main contribution may be due to sources with $`z>0.5`$, while (F450W$``$F814W)<sub>AB</sub>, (F606W$``$F814W)<sub>AB</sub> and (F450W$``$F606W)<sub>AB</sub> are compatible with those of local irregular galaxies.
## 6 Morphological Number Counts
We performed a morphological analysis of the F814W-selected sample, limited to F814W<sub>AB</sub>=25 (Volonteri et al. 2000). We chose a quantitative classification of galaxies morphology, following Abraham et al. (1994, 1996). For each galaxy we measured an asymmetry index ($`A`$) and a concentration index ($`C_{abr}`$). The former is determined by rotating the galaxy by 180 and subtracting the resulting image from the original one. The asymmetry index is given by the sum of absolute values of the pixels in the residual image, normalized by the sum of the absolute value of the pixels in the original image and corrected for the intrinsic asymmetry of the background. The concentration index is given by the ratio of fluxes in two isophotes, based on the analysis of light profiles.
We determined number counts by splitting the sample in three bins: E/S0, spirals and irregular/peculiar/ interacting according to our quantitative classification. In every bin the error is estimated as explained in Section 3, considering only contributions due to Poissonian and clustering fluctuation, since according to our simulations the sample at F814W<sub>AB</sub>=25 is complete.
In Figure 12 the decrease in differential counts for early-type galaxies at F814W$`{}_{AB}{}^{}>`$22.5, and the steeper slope for spiral galaxies and irregular/peculiar/interacting ones is clear. They are described by a slope $`\gamma _{irr}=0.43\pm 0.05`$, $`\gamma _{spiral}=0.37\pm 0.05`$, we discuss this feature further in Section 7.
Abraham et al. (1996) and Driver et al. (1998) studied the morphological number counts for all galaxies with F814W$`{}_{AB}{}^{}<`$26.0 detected in the HDF-N. Our results are in agreement with theirs: number counts are dominated by late type galaxies and early type sources show a flat curve. Results for late type galaxies are compatible with a strong evolution in number or with a non negligible fraction of sources being at moderate redshift, as suggested also by the analysis of colors. This result fits in Marzke et al. (1997) findings about the LF of galaxies as a function of their color: the slope of the LF faint end is steeper at low luminosities for the bluest galaxies ($`M_B^{}19.4`$, $`\alpha <1.7`$ for galaxies bluer than B$``$V=0.4). Actually by integrating this LF, within $`z=0.5`$, the number density of galaxies bluer than an Irregular galaxy is compatible with the number found in our sample. The flattening of the curve for early type sources may be due to a decrease in density at higher redshift of elliptical galaxies, but it is also expected by the $`\alpha >1`$ slope of the LF of early type galaxies (e.g. Marzke et al. 1998). Number counts models nevertheless show that an $`\alpha >1`$ LF does not account for such a steep decrease (Figure 12).
Brinchmann et al. (1998), Im et al. (1999) and Driver et al. (1998) studied the redshift distribution for different morphological types in deep surveys. Brinchmann et al. (1998) analyzed 341 galaxies at $`z<1`$ from the Canada-France Redshift Survey (CFRS, Lilly et al. 1995, limited at F814W<sub>AB</sub>=22.5), and Autofib/Low-Dispersion Survey Spectrograph (LDSS, Ellis et al. 1996, limited at $`b_j=24`$). They find the distribution of elliptical galaxies peaked at $`z0.5`$, while spiral galaxies have a shallow distribution and the excess of peculiar sources begins at $`z0.4`$. According to their conclusions, the number of regular galaxies is compatible with a passive evolution model, while the excess of peculiar ones suggests an active luminosity evolution and-or number evolution. Im et al. (1999) for a sample of 464 galaxies limited at I$`=21.5`$ do agree with Brinchmann et al. (1998). Driver et al. (1998) analyzed HDF-N sources, with photometric redshifts, and underline the excess both in spiral galaxies at $`z=1.5`$ and irregulars at $`z>1`$, suggesting number evolution and passive evolution for dwarf LSBG at $`z<1`$ giving the density of irregular galaxies at low redshift.
## 7 Constraining the Redshift of Galaxies from Photometric Data
In a deep survey, such as the HDF-S, the contribution of low redshift foreground sources should be disentangled from the contribution of high redshift ones. We first tried to split the two populations with very simple considerations, regarding galaxy colors and counts.
We compared colors of HDF-S galaxies with those predicted by CWW spectra and by synthetic spectra of Bruzual & Charlot (1993, BC hereafter). While BC account for evolution, the CWW spectra take into account $`Kcorrections`$. In Figures 13-14 apparent colors from CWW and BC are shown. We limited our analysis at $`z<2`$ to avoid relying too much on the UV part of spectral energy distributions which has a strong influence at higher redshift.
By comparison with colors predicted by CWW or BC spectra, the observed very blue galaxies (med(F450W$`{}_{AB}{}^{}`$ F606W<sub>AB</sub>) = 0.06, med(F606W$`{}_{AB}{}^{}`$ F814W<sub>AB</sub>) = 0.24) are compatible with two very different ranges in redshift: $`z<0.2`$ or $`z>1.5`$ (and perhaps beyond our analysis limit). The slope of morphological number counts is in good agreement with this hypothesis, assuming that most of the faint late type galaxies are at moderate redshift. In the hypothesis of the moderate redshift population ($`z0.2)`$ we computed number counts as a function of optical colors, which should reflect rest-frame colors for this range in z. We divided the F814W-selected sample in a “red” subsample, composed by sources redder than the median colour of the sample (F450W$`{}_{AB}{}^{}`$ F606W<sub>AB</sub>=0.45) and a “blue” subsample of the remaining sources. We limited our analysis at F814W<sub>AB</sub>=26, where the catalogue is complete, according to our simulations.
The curve for blue sources is much steeper than the red one, as shown in Figure 15. We estimated $`\gamma _{blue}0.49\pm 0.01`$. This value is therefore compatible with the previous hypothesis, that a non negligible fraction of the blue galaxies may have a moderate redshift, contributing with an almost Euclidean slope to the counts or with number evolution, with very faint galaxies merging at high redshift. However we cannot rule out the hypothesis that such a steep slope can be due to earlier-type evolving galaxies which move into the blue sample at high redshift. In such a case the steep counts described by the faint blue sample would not require a significant rate of merging to be justified.
In order to check the low-redshift solution we compared the characteristics of those galaxies to the predictions of a local luminosity function. The brightest of the very blue sources would have $`M_B18`$ if placed at $`z=0.2`$ (assuming $`H_0`$=50 km s<sup>-1</sup> Mpc<sup>-1</sup> and a (F450W$``$F814W)<sub>AB</sub>=1).
Considering the local luminosity function (LF) of blue galaxies as estimated by Marzke et al. (1997), described by $`\alpha =1.8`$ and $`M_B^{}`$19.5 and integrating within $`z=0.2`$ and $`M_B=`$18, the number density found is compatible with all very blue galaxies (i.e. galaxies with (F450W$``$F606W)<sub>AB</sub> and (F606W$``$F814W)<sub>AB</sub> bluer than a local Irregular galaxy) being at such a low redshift. The first raw technique used to constrain the redshift of sources suggests that a high fraction of faint galaxy is composed by low redshift ($`z<0.2`$) sources, nevertheless the “high redshift solution” is not ruled out, since steep counts may be due also to merging and our color analysis is very naive.
We then made our analysis better by means of the photometric redshift technique. We computed photometric redshifts for the galaxies in the F814W catalogue by means of the public code hyperz (Bolzonella et al. 2000). The method basically consists in a comparison between the observed magnitudes and the photometry expected from a set of Spectral Energy Distributions (SEDs). The efficiency of the method relies on the presence of strong spectral features, in particular the $`4000`$ Å break and the Lyman break for galaxies at high redshift. To reduce the uncertainties estimating $`z_{phot}`$, the set of filters must be able to identify these characteristics, spanning a wide range of wavelength. The photometric redshift $`z_{phot}`$ is computed by a $`\chi ^2`$ minimization, considering the whole set of possibilities with different combinations of the involved parameters.
In this procedure we decided to use the SEDs built from the Bruzual & Charlot’s synthetic library (GISSEL98, Bruzual & Charlot 1993), rather than the observed mean local spectra by CWW, because tests on HDF-N indicate a slightly lower dispersion at $`z<2`$. In any case, CWW spectra lead to not very different conclusions. Hence, we selected 5 spectral types, characterized by different star formation rates, matching the observed colors of the morphological galaxy sequence; a subsample of ages allowed by the GISSEL library is also considered. Moreover, we took into account the possible presence of dust applying the reddening law by Calzetti et al. (2000) for different values of $`E_{BV}`$. The flux decrement produced by intervening neutral hydrogen is computed following the recipe by Madau (1995). Only spectra with solar metallicity are considered here, because the metallicity can be regarded as a secondary parameter.
The accuracy of photometric redshift estimate obtained from the considered set of filters can be studied by means of simulations on synthetic catalogue.
Some degeneracy can be found (Bolzonella et al., 2000), with a non negligible part of galaxies lying at $`z1`$$`2.4`$ incorrectly located at low redshift. These galaxies show frequently a probability function with several comparable peaks, at low and high $`z`$, due to the degeneracy in the parameter space. Near infrared data can in principle avoid these uncertainties. Nonetheless, also with the available set of filters, we can guess that the objects with $`z_{phot}1`$ really belong to this redshift, because high-$`z`$ objects are rarely misidentified as low-$`z`$ galaxies.
Photometric redshifts for our catalogue indicate that all very blue galaxies are at $`z_{phot}1`$ (med($`z_{phot}`$)=1.6). This feature, along with their steep slope (see Section 5), suggests a high merging rate for galaxies at $`z_{phot}1`$. Vice versa very red sources have a shallower redshift distribution (med($`z_{phot}`$)=0.6), though at F814W$`{}_{AB}{}^{}>25`$ most sources have a high redshift.
Similarly, at F814W$`{}_{AB}{}^{}<`$25, late type galaxies have a higher median redshift (med($`z_{phot}`$)=1.1) than early type ones (med($`z_{phot}`$)=0.6) and the redshift distribution (Figure 16) is similar to Driver et al. (1998) one, though the number of sources at $`z_{phot}>1.5`$ is lower. There is a point to be emphasized about redshift distributions and morphological types: the detection of galaxies at high redshift is seriously biased against elliptical galaxies because of $`Kcorrections`$ in the F814W-band. At $`z1`$, early type galaxies have $`Kcorrections`$ ($`1`$ mag) greater than late type $`Kcorrections`$, making them fainter. It means that a cut in apparent magnitude, such as F814W$`{}_{AB}{}^{}<`$25, biases the sample toward galaxies with lower $`Kcorrections`$. To skip this problem, we tried to select a volume-limited sample, based on our photometric redshifts. When selecting all galaxies with $`z_{phot}1`$, we found only 15$`\%`$ to be with F814W$`{}_{AB}{}^{}<`$25, that is with a reliable morphological classification according to our method. Vice versa it means that a cut off in apparent magnitude does not offer a fair sample for a redshift distribution of the different morphological types.
Probably also the apparent decrease in early-type galaxies number counts, as shown in Figure 12, is due to this bias.
We compared our number counts with different models, by using galaxy counts models by Gardner (1998). We adopted a flat cosmological model ($`q_0=0.5`$, $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup>) in all models. We used the luminosity function of Marzke et al. (1998) both considering a single Schechter function for all morphological types and three different Schechter fits for elliptical, spiral and irregular galaxies. Our counts are best fitted by considering three LFs, luminosity evolution and a moderate merging (Figure 17) consistent with our interpretations. The merging rate in Gardner (1998) follows Rocca-Volmerange & Guiderdoni (1990) with number evolution parameterized as $`\varphi ^{}(1+z)^\eta `$ in the Schechter fit to the LF, and in order to conserve the luminosity density, $`L^{}(1+z)^\eta `$; $`\eta `$ is a free parameter. Good fits to our data are obtained with $`\eta 0.5`$ in the case of three LFs for the cosmological parameters considered.
## 8 Summary and Conclusions
The HDF-S represents a unique opportunity for the study of faint galaxies up to now, both for its depth (F814W$`{}_{AB,lim}{}^{}29`$ for detection and F814W$`{}_{AB,lim}{}^{}27`$ for completeness) and spatial resolution (FWHM$``$0.2 arcsec). We presented here colors and number counts of HDF-S galaxies, along with number counts determined by splitting the sample considering the morphology and the colors of the galaxies. We also analyzed the photometric data to constrain the redshift of HDF-S galaxies, and determine the contribution of different redshift populations to the counts. The main results are the following:
* the number-counts relation has an increasing slope up to the limits of the survey in all four bands. The slope is steeper at shortest wavelengths;
* the number counts model which best fits our data is obtained by considering three different Schechter functions for elliptical, spiral and irregular galaxies (Marzke et al. 1998), luminosity evolution and a moderate merging ($`\eta <0.5`$);
* optical colors show that the sample contains a high fraction of galaxies bluer than local sources, for instance at F814W$`{}_{AB}{}^{}>27`$ about $`50\%`$ of sources have (F450W$``$F606W)<sub>AB</sub> bluer than a typical local irregular galaxy;
* after an initial blueing trend the color-magnitude relation gets flat, suggesting a color-saturation due to a strong contribution of a flat spectrum population in the faintest bins;
* morphological number counts ( F814W$`{}_{AB}{}^{}<25`$) are dominated by late type galaxies ($`\gamma _{irr}=0.43\pm 0.03`$, $`\gamma _{spir}=0.37\pm 0.03`$), while early type galaxies show a steep decrease in the faintest bin;
* photometric redshifts of our sample galaxies show that the galaxies contributing with a steep slope to the number counts have $`z_{phot}1`$, suggesting a moderate merging;
* morphological redshift distributions and number counts are biased against elliptical galaxies when using an apparent magnitude cut-off.
Our number counts are in good agreement with previous results, except for the F300W-band number counts in the HDF-N.
We do not see a flattening in F300W-band number counts, limiting our analysis at F300W$`{}_{AB}{}^{}<`$26.5. A steep F300W-slope may be explained both with a high fraction of low redshift galaxies or with merging or a mixture of them. In the former case, intrinsically faint, low redshift galaxies, imply a steep faint end slope of the LF in the F300W-band ($`\alpha 2`$) in the latter case we emphasize that if merging is the reason of the steep F300W-band counts, it would take place at $`z<2`$ since the Lyman break prevents us from detecting galaxies in the F300W-band at higher redshift. A way of detecting merging and hence to check the validity of the hierarchical model, could be in principle the study of the evolution with look-back time of some structural parameters of galaxies, such as their intrinsic size.
Morphological number counts are in agreement with previous results, as well as the morphological redshift distributions. The morphological classification is, however, possible only for bright galaxies. The cut in apparent magnitude biases the sample against early-type galaxies, due to their large $`Kcorrections`$. A correct test implying morphological classification should be done using K-band data (an extension of the test proposed by Kauffmann & Charlot, 1998) or a volume-limited sample, with morphology determined from spectra. Extracting a volume-limited sample from a magnitude limited survey does not fit the goal since at $`z1`$ almost 80$`\%`$ of the galaxies are too faint for a reliable morphological classification from imaging.
###### Acknowledgements.
We would like to thank Roser Pelló for the code hyperz, R. Abraham and J.P. Gardner for making available their software and A. Buzzoni, M. Massarotti for interesting and stimulating discussions.
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# Modal Logics of Topological Spaces (Extended abstract of doctoral dissertation)
## 1 Introduction
In this thesis we shall present two logical systems, $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$, for the purpose of reasoning about knowledge and effort. These logical systems will be interpreted in a spatial context and therefore, the abstract concepts of knowledge and effort will be defined by concrete mathematical concepts.
Our general framework consists of a set of possible worlds (situations, scenarios, consistent theories, etc.) A state of knowledge is a subset of this set and our knowledge consists of all facts common to the worlds belonging to this subset. This subset of possibilities can be thought as our view. Thus two knowers having distinct views can have different knowledge. This treatment of knowledge agrees with the traditional one (\[Hin62\], \[HM84\], \[PR85\], \[CM86\], \[FHV91\]) expressed in a variety of contexts (artificial intelligence, distributed processes, economics, etc.)
Our treatment is based on the following simple observation
> “a restriction of our view increases our knowledge.”
This is because a smaller set of possibilities implies a greater amount of common facts. Moreover, such a restriction can only be possible due to an increase of information. And such an information increase can happen with spending of time or computation resources. Here is where the notion of effort enters. A restriction of our view is dynamic (contrary to the view itself which is a state) and is accompanied by effort during which a greater amount of information becomes available to us (Pratt expresses a similar idea in the context of processes \[Pra92\].)
We make two important assumptions.
Our knowledge has a subject. We collect information for a specific purpose. Hence we are not considering arbitrary restrictions to our view but restrictions parameterized by possibilities contained in our view, i.e. neighborhoods of possibilities. After all, only one of these possibilities is our actual state. This crucial assumption enables us to express topological concepts and use a mathematical set-theoretic setting as semantics. Without such an assumption these ideas would have been expressed in the familiar theory of intuitionism ( \[Hey56\]\[Dum77\]\[TvD88\].) As Fitting points out in \[Fit69\]
> “Let $`𝒢,,`$ be a \[intuitionistic, propositional\] model. $`𝒢`$ is intended to be a collection of possible universes, or more properly, states of knowledge. Thus a particular $`\mathrm{\Gamma }`$ in $`𝒢`$ may be considered as a collection of (physical) facts known at a particular time. The relation $``$ represents (possible) time succession. That is, given two states of knowledge $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ of $`𝒢`$, to say $`\mathrm{\Gamma }\mathrm{\Delta }`$ is to say: if we now know $`\mathrm{\Gamma }`$, it is possible that later we will know $`\mathrm{\Delta }`$.”
Considering neighborhoods and, inevitably, points which parameterize neighborhoods, the important duality between the facts, which constitute our knowledge, and the possible worlds, where such facts hold, emerges.
The other assumption is that of indeterminacy. Each state of knowledge is closed under logical deduction. Thus an increase of knowledge can happen only by a piece of evidence or information given from outside. Our knowledge is external (a term used by Parikh to describe a similar idea in \[Par87b\].) This fact leads to indeterminacy (we do not know which kinds of information will be available to us, if at all) and resembles indeterminacy expressed in intuitionism through the notion of lawless sequence (see \[Kre58\]\[Tro77\]) where, not surprisingly, topological notions arise.
To illustrate better these simple but fundamental ideas we present the following examples:
* Suppose that a machine emits a stream of binary digits representing the output of a recursive function $`f`$. After time $`t_1`$ the machine emitted the stream $`111`$. The only information we have about the function being computed at this time on the basis of this (finite) observation is that
$$f(1)=f(2)=f(3)=1.$$
As far as our knowledge concerns $`f`$ is indistinguishable from the constant function $`\mathrm{𝟏}`$, where $`\mathrm{𝟏}(n)=1`$ for all $`n`$. After some additional time $`t_2`$, i.e. spending more time and resources, $`0`$ might appear and thus we could be able to distinguish $`f`$ from $`\mathrm{𝟏}`$. In any case, each binary stream will be an initial segment of $`f`$ and this initial segment is a neighborhood of $`f`$. In this way, we can acquire more knowledge for the function the machine computes. The space of finite binary streams is a structure which models computation. Moreover, this space comprises a topological space. The set of binary streams under the prefix ordering is an example of Alexandrov topology (see \[Vic89\].)
* A policeman measures the speed of passing cars by means of a device. The speed limit is $`80`$ km/h. The error in measurement which the device introduces is $`1`$ km/h. So if a car has a speed of $`79.5`$ km/h and his device measures $`79.2`$ km/h then he knows that the speed of the passing car lies in the interval $`(78.2,80.2)`$ but he does not know if the car exceeds the speed limit because not all values in this interval are more than $`80`$. However, measuring again and combining the two measurements or acquiring a more accurate device he has the possibility of knowing that a car with a speed of $`79.5`$ km/h does not exceed the speed limit. Note here that if the measurement is, indeed, an open interval of real line and the speed of a passing car is exactly $`80`$ km/h then he would never know if such a car exceeded the speed limit or not.
To express this framework we use two modalities $`𝖪`$ for knowledge and $`\mathrm{}`$ for effort. Moss and Parikh observed in \[MP92\] that if the formula
$$A\mathrm{}𝖪A$$
is valid, where $`A`$ is an atomic predicate and $`\mathrm{}`$ is the dual of the $`\mathrm{}`$, i.e. $`\mathrm{}\neg \mathrm{}\neg `$, then the set which $`A`$ represents is an open set of the topology where we interpret our systems. Under the reading of $`\mathrm{}`$ as “possible” and $`𝖪`$ as “is known”, the above formula says that
> “if $`A`$ is true then it is for $`A`$ possible to be known”,
i.e. $`A`$ is affirmative. Vickers defines similarly an affirmative assertion in \[Vic89\]
> “an assertion is affirmative iff it is true precisely in the circumstances when it can be affirmed.”
The validity of the dual formula
$$\mathrm{}𝖫AA,$$
where $`𝖫`$ is the dual of $`𝖪`$, i.e. $`𝖫\neg 𝖪\neg `$, expresses the fact that the set which $`A`$ represents is closed, and hence $`A`$ is refutative, meaning if it does not hold then it is possible to know that. The fact that affirmative and refutative assertions are represented by opens and closed subsets, respectively, should not come to us as a surprise. Affirmative assertions are closed under infinite disjunctions and refutative assertions are closed under infinite conjunctions. Smyth in \[Smy83\] observed first these properties in semi-decidable properties. Semi-decidable properties are those properties whose truth set is r.e. and are a particular kind of affirmative assertions. In fact, changing our power of affirming or computing we get another class of properties with a similar knowledge-theoretic character. For example, using polynomial algorithms affirmative assertions become polynomially semi-decidable properties. If an object has this property then it is possible to know it with a polynomial algorithm even though it is not true we know it now.
Does this framework suffers from the problem of logical omniscience? Only in part. Expressing effort we are able to bound the increase of knowledge depending on information (external knowledge.) Since the modality $`𝖪`$ which corresponds to knowledge is axiomatized by the normal modal logic of $`\mathrm{𝐒𝟓}`$, knowledge is closed under logical deduction. However, because of the strong computational character of this framework it does not seem unjustified to assume that in most cases (as in the binary streams example) a finite amount of data restricts our knowledge to a finite number of (relevant) formulae. Even without such an assumption we can incorporate the effort to deduce the knowledge of a property in the passage from one state of knowledge to the other.
We have made an effort to present our material somewhat independently. However, knowledge of basic modal logic, as in \[Che80\], \[HC84\], or \[Fit93\], is strongly recommended.
The language and semantics of our logical framework is presented in Chapter 2. In the same Chapter we present two systems: $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$. The former was introduced in \[MP92\] and was proven complete for arbitrary sets of subsets. It soon became evident that such sets of subsets should be combined, whenever it is possible, to yield a further increase of knowledge or we should assume a previous state of other states of knowledge where such states are a possible. Therefore the set of subsets should be closed under union and intersection. Moreover, topological notions expressed in $`\mathrm{𝐌𝐏}`$ make sense only in topological models. For this reason we introduce an extension of the set of axioms of $`\mathrm{𝐌𝐏}`$ and we call it $`\mathrm{𝐌𝐏}^{}`$. In Chapter 3, we study the topological models of $`\mathrm{𝐌𝐏}^{}`$ by semantical means. We are able to prove the reduction of the theory of topological models to models whose associated set of subsets is closed under finite union and intersection. Finding for each satisfiable formula a model of bounded size we prove decidability for $`\mathrm{𝐌𝐏}^{}`$. The results of this chapter will appear in \[Geo93\]. In Chapter 4, we prove that $`\mathrm{𝐌𝐏}^{}`$ is a complete system for topological models as well as topological models comprised by closed subsets. We also give necessary and sufficient conditions for turning a Kripke frame into such a topological model. In Chapter 5, we present the modal algebras of $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$ and some of their properties. Finally, in Chapter 6, we present some of our ideas towards future work.
## 2 Two Systems: $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$
In section 2.1 we shall present a language and semantics which appeared first in \[MP92\]. In section 2.2, we shall present the axiom system $`\mathrm{𝐌𝐏}`$, introduced and proven sound and complete with a class of models called subset spaces in \[MP92\], and the axiom system $`\mathrm{𝐌𝐏}^{}`$, introduced by us, which we shall prove sound and complete for, among other classes, the class of topological spaces.
### 2.1 Language and Semantics
We follow the notation of \[MP92\].
Our language is bimodal and propositional. Formally, we start with a countable set $`𝖠`$ of atomic formulae containing two distinguished elements $``$ and $``$. Then the language $``$ is the least set such that $`𝖠`$ and closed under the following rules:
$$\frac{\varphi ,\psi }{\varphi \psi }\frac{\varphi }{\neg \varphi ,\mathrm{}\varphi ,𝖪\varphi }$$
The above language can be interpreted inside any spatial context.
Definition 1 Let $`X`$ be a set and $`𝒪`$ a subset of the powerset of $`X`$, i.e. $`𝒪𝒫(X)`$ such that $`X𝒪`$. We call the pair $`X,𝒪`$ a subset space. A model is a triple $`X,𝒪,i`$, where $`X,𝒪`$ is a subset space and $`i`$ a map from $`𝖠`$ to $`𝒫(X)`$ with $`i()=X`$ and $`i()=\mathrm{}`$ called initial interpretation.
We denote the set $`\{(x,U):xX,U𝒪,\text{ and }xU\}X\times 𝒪`$ by $`X\dot{\times }𝒪`$. For each $`U𝒪`$ let $`U`$ be the set $`\{V:V𝒪\text{ and }VU\}`$ the lower closed set generated by $`U`$ in the partial order $`(𝒪,)`$, i.e. $`U=𝒫(U)𝒪`$.
Definition 2 The satisfaction relation $`_{}`$, where $``$ is the model $`X,𝒪,i`$, is a subset of $`(X\dot{\times }𝒪)\times `$ defined recursively by (we write $`x,U_{}\varphi `$ instead of $`((x,U),\varphi )_{}`$):
$$\begin{array}{cc}x,U_{}A\hfill & \text{iff}xi(A),\text{ where }A𝖠\hfill \\ x,U_{}\varphi \psi \hfill & \text{if}x,U_{}\varphi \text{ and }x,U_{}\psi \hfill \\ x,U_{}\neg \varphi \hfill & \text{if}x,U\vDash ̸_{}\varphi \hfill \\ x,U_{}𝖪\varphi \hfill & \text{if}\text{for all }yU,y,U_{}\varphi \hfill \\ x,U_{}\mathrm{}\varphi \hfill & \text{if}\text{for all }VU\text{ such that }xV,x,V_{}\varphi .\hfill \end{array}$$
If $`x,U_{}\varphi `$ for all $`(x,U)`$ belonging to $`X\dot{\times }𝒪`$ then $`\varphi `$ is valid in $``$, denoted by $`\varphi `$.
We abbreviate $`\neg \mathrm{}\neg \varphi `$ and $`\neg 𝖪\neg \varphi `$ by $`\mathrm{}\varphi `$ and $`𝖫\varphi `$ respectively. We have that
$$\begin{array}{cc}x,U_{}𝖫\varphi \hfill & \text{if there exists }yU\text{ such that }y,U_{}\varphi \hfill \\ x,U_{}\mathrm{}\varphi \hfill & \text{if there exists }V𝒪\text{ such that }VU,xV,\text{ and }x,V_{}\varphi .\hfill \end{array}$$
Many topological properties are expressible in this logical system in a natural way. For instance, in a model where the subset space is a topological space, $`i(A)`$ is open whenever $`A\mathrm{}𝖪A`$ is valid in this model. Similarly, $`i(A)`$ is nowhere dense whenever $`𝖫\mathrm{}𝖪\neg A`$ is valid (cf. \[MP92\].)
Example. Consider the set of real numbers $`𝐑`$ with the usual topology of open intervals. We define the following three predicates:
$`\mathrm{𝚙𝚒}`$ where $`i(\mathrm{𝚙𝚒})=\{\pi \}`$
$`𝙸_1`$ where $`i(𝙸_1)=(\mathrm{},\pi ]`$
$`𝙸_2`$ where $`i(𝙸_2)=(\pi ,+\mathrm{})`$
$`𝚀`$ where $`i(𝚀)=\{q:q\text{ is rational }\}.`$
There is no real number $`p`$ and open set $`U`$ such that $`p,U𝖪\mathrm{𝚙𝚒}`$ because that would imply $`p=\pi `$ and $`U=\{\pi \}`$ and there are no singletons which are open.
A point $`x`$ belongs to the closure of a set $`W`$ if every open $`U`$ that contains $`x`$ intersects $`W`$. Thus $`\pi `$ belongs to the closure of $`(\pi ,+\mathrm{})`$, i.e every open that contains $`\pi `$ has a point in $`(\pi ,+\mathrm{})`$. This means that for all $`U`$ such that $`\pi U`$, $`\pi ,U𝖫𝙸_2`$, therefore $`\pi ,𝐑\mathrm{}𝖫𝙸_2`$. Following the same reasoning $`\pi ,𝐑\mathrm{}𝙸_1`$, since $`\pi `$ belongs to the closure of $`(\mathrm{},\pi ]`$.
A point $`x`$ belongs to the boundary of a set $`W`$ whenever $`x`$ belong to the closure of $`W`$ and $`XW`$. By the above, $`\pi `$ belongs to the boundary of $`(\mathrm{},\pi ]`$ and $`\pi ,𝐑\mathrm{}(𝖫𝙸_1𝖫𝙸_2)`$.
A set $`W`$ is closed if it contains its closure. The interval $`i(𝙸_1)=(\mathrm{},\pi ]`$ is closed and this means that the formula $`\mathrm{}𝖫𝙸_1𝙸_1`$ is valid.
A set $`W`$ is dense if all opens contain a point of $`W`$. The set of rational numbers is dense which translates to the fact that the formula $`\mathrm{}𝖫𝚀`$ is valid. To exhibit the reasoning in this logic, suppose that the set of rational numbers was closed then both $`\mathrm{}𝖫𝚀`$ and $`\mathrm{}𝖫𝚀𝚀`$ would be valid. This implies that $`𝚀`$ would be valid which means that all reals would be rationals. Hence the set of rational numbers is not closed.
### 2.2 $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$
The axiom system $`\mathrm{𝐌𝐏}`$ consists of axiom schemes 1 through 10 and rules of table 1 (see page 1) and appeared first in \[MP92\].
The following was proved in \[MP92\].
###### Theorem 3
The axioms and rules of $`\mathrm{𝐌𝐏}`$ are sound and complete with respect to subset spaces.
We add the axioms 11 and 12 to form the system $`\mathrm{𝐌𝐏}^{}`$ for the purpose of axiomatizing spaces closed under union and intersection and, in particular, topological spaces.
A word about the axioms (most of the following facts can be found in any introductory book about modal logic, e.g. \[Che80\] or \[Gol87\].) The axiom 2 expresses the fact that the truth of atomic formulae is independent of the choice of subset and depends only on the choice of point. This is the first example of a class of formulae which we are going to call bi-persistent and their identification is one of the key steps to completeness. Axioms 3 through 5 and axioms 6 through 9 are used to axiomatize the normal modal logics S4 and S5 respectively. The former group of axioms expresses the fact that the passage from one subset to a restriction of it is done in a constructive way as actually happens to an increase of information or a spending of resources (the classical interpretation of necessity in intuitionistic logic is axiomatized in the same way). The latter group is generally used for axiomatizing logics of knowledge.
Axiom 10 expresses the fact that if a formula holds in arbitrary subsets is going to hold as well in the ones which are neighborhoods of a point. The converse is not sound.
Axiom 11 is a well-known formula which characterizes incestual frames, i.e. if two points $`\beta `$ and $`\gamma `$ in a frame can be accessed by a common point $`\alpha `$ then there is a point $`\delta `$ which can be accessed by both $`\beta `$ and $`\gamma `$. It appeared in the equivalent form (in \[MP92\])
$$\mathrm{}\mathrm{}\varphi \mathrm{}\mathrm{}\psi \mathrm{}\mathrm{}(\varphi \psi )$$
and was proved sound in subset spaces closed under (finite) intersection.
Obviously our attention is focused on axiom 12. It is sound in spaces closed under (finite) union and intersection as the following proposition shows.
###### Proposition 4
Axioms 1 through 12 are sound in the class of subset spaces closed under finite union and intersection.
Proof. Soundness for Axioms 1 through 11 is easy. For Axiom 12, suppose
$$x,U\mathrm{}(𝖪\varphi \psi )𝖫\mathrm{}(𝖪\varphi \chi ).$$
Since $`x,U\mathrm{}(𝖪\varphi \psi )`$, there exists $`U_xU`$ such that
$$x,U_x𝖪\varphi \psi $$
and, since $`x,U𝖫(\mathrm{}𝖪\varphi \mathrm{}\chi )`$, there exists $`yU`$ and $`U_yU`$ such that
$$y,U_y𝖪\varphi \chi .$$
We now have that $`U_xU_yU`$ (we assume closure under unions.) Thus
$$x,U_xU_y𝖪\mathrm{}\varphi ,y,U_xU_y𝖪\mathrm{}\varphi ,x,U_xU_y\mathrm{}\psi ,\text{and}y,U_xU_y\mathrm{}\chi .$$
Therefore,
$$x,U\mathrm{}(𝖪\mathrm{}\varphi \mathrm{}\psi 𝖫\mathrm{}\chi ).$$
With the help of axiom 12 we are able to prove the key lemma 33 which leads to the DNF Theorem (page 35.) and this is the only place where we actually use it. Any formula, sound in the class of subset spaces closed under finite union and intersection, which implies the formula (note the difference from axiom 12)
$$\mathrm{}(𝖪\varphi \psi )𝖫\mathrm{}(𝖪\varphi \chi )\mathrm{}(𝖪\varphi \psi 𝖫\chi )$$
where $`\mathrm{}\varphi \mathrm{}\varphi `$, $`\mathrm{}\psi \psi `$ and $`\chi \mathrm{}\chi `$ are theorems, can replace axiom 12.
## 3 A Semantical analysis of $`\mathrm{𝐌𝐏}^{}`$
### 3.1 Stability and Splittings
Suppose that $`X`$ is a set and $`𝒯`$ a topology on $`X`$. In the following we assume that we are working in the topological space $`(X,𝒯)`$. Our aim is to find a partition of $`𝒯`$, where a given formula $`\varphi `$ “retains its truth value” for each point throughout a member of this partition. We shall show that there exists a finite partition of this kind.
Definition 5 Given a finite family $`=\{U_1,\mathrm{},U_n\}`$ of opens, we define the remainder of (the principal ideal in $`(𝒯,)`$ generated by) $`U_k`$ by
$$\mathrm{𝖱𝖾𝗆}^{}U_k=U_k\underset{U_kU_i}{}U_i.$$
###### Proposition 6
In a finite set of opens $`=\{U_1,\mathrm{},U_n\}`$ closed under intersection, we have
$$\mathrm{𝖱𝖾𝗆}^{}U_i=U_i\underset{U_jU_i}{}U_j,$$
for $`i=1,\mathrm{},n`$.
We denote $`_{U_i}U_i`$ with $``$.
###### Proposition 7
If $`=\{U_1,\mathrm{},U_n\}`$ is a finite family of opens, closed under intersection, then
1. $`\mathrm{𝖱𝖾𝗆}^{}U_i\mathrm{𝖱𝖾𝗆}^{}U_j=\mathrm{}`$, for $`ij`$,
2. $`_{i=1}^n\mathrm{𝖱𝖾𝗆}^{}U_i=`$, i.e. $`\{\mathrm{𝖱𝖾𝗆}^{}U_i\}_{i=1}^n`$ is a partition of $``$. We call such an $``$ a finite splitting (of $``$),
3. if $`V_1,V_3\mathrm{𝖱𝖾𝗆}^{}U_i`$ and $`V_2`$ is an open such that $`V_1V_2V_3`$ then $`V_2\mathrm{𝖱𝖾𝗆}^{}U_i`$, i.e. $`\mathrm{𝖱𝖾𝗆}^{}U_i`$ is convex.
Every partition of a set induces an equivalence relation on this set. The members of the partition comprise the equivalence classes. Since a splitting induces a partition, we denote the equivalence relation induced by a splitting $``$ by $`_{}`$.
Definition 8 Given a set of open subsets $`𝒢`$, we define the relation $`_𝒢^{}`$ on $`𝒯`$ with $`V_1_𝒢^{}V_2`$ if and only if $`V_1UV_2U`$ for all $`U𝒢`$.
We have the following
###### Proposition 9
The relation $`_𝒢^{}`$ is an equivalence.
###### Proposition 10
Given a finite splitting $``$, $`_{}^{}=_{}`$ i.e. the remainders of $``$ are the equivalence classes of $`_{}^{}`$.
We state some useful facts about splittings.
###### Proposition 11
If $`𝒢`$ is a finite set of opens, then $`\mathrm{𝖢𝗅}(𝒢)`$, its closure under intersection, yields a finite splitting for $`𝒢`$.
The last proposition enables us to give yet another characterization of remainders: every family of points in a complete lattice closed under arbitrary joins comprises a closure system, i.e. a set of fixed points of a closure operator of the lattice (cf. \[GHK<sup>+</sup>80\].) Here, the lattice is the poset of the opens of the topological space. If we restrict ourselves to a finite number of fixed points then we just ask for a finite set of opens closed under intersection i.e. Proposition 11. Thus a closure operator in the lattice of the open subsets of a topological space induces an equivalence relation, two opens being equivalent if they have the same closure, and the equivalence classes of this relation are just the remainders of the open subsets which are fixed points of the closure operator. The maximum open in $`\mathrm{𝖱𝖾𝗆}^{}U`$, i.e. $`U`$, can be taken as the representative of the equivalence class which is the union of all open sets belonging to $`\mathrm{𝖱𝖾𝗆}^{}U`$.
We now introduce the notion of stability corresponding to what we mean by “a formula retains its truth value on a set of opens”.
Definition 12 If $`𝒢`$ is a set of opens then $`𝒢`$ is stable for $`\varphi `$, if for all $`x`$, either $`x,V\varphi `$ for all $`V𝒢`$, or $`x,V\neg \varphi `$ for all $`V𝒢`$, such that $`xV`$.
###### Proposition 13
If $`𝒢_1`$,$`𝒢_2`$ are sets of opens then
1. if $`𝒢_1𝒢_2`$ and $`𝒢_2`$ is stable for $`\varphi `$ then $`𝒢_1`$ is stable for $`\varphi `$ ,
2. if $`𝒢_1`$ is stable for $`\varphi `$ and $`𝒢`$ is stable for $`\chi `$ then $`𝒢_1𝒢_2`$ is stable for $`\varphi \chi `$.
Definition 14 A finite splitting $`=\{U_1,\mathrm{},U_n\}`$ is called a stable splitting for $`\varphi `$, if $`\mathrm{𝖱𝖾𝗆}^{}U_i`$ is stable for $`\varphi `$ for all $`U_i`$.
###### Proposition 15
If $`=\{U_1,\mathrm{},U_n\}`$ is a stable splitting for $`\varphi `$, so is
$$^{}=\mathrm{𝖢𝗅}(\{U_0,U_1,\mathrm{},U_n\}),$$
where $`U_0`$.
The above proposition tells us that if there is a finite stable splitting for a topology then there is a closure operator with finitely many fixed points whose associated equivalence classes are stable sets of open subsets.
Suppose that $`=X,𝒯,i`$ is a topological model for $``$. Let $`_{}`$ be a family of subsets of $`X`$ generated as follows: $`i(A)_{}`$ for all $`A𝖠`$, if $`S_{}`$ then $`XS_{}`$, if $`S,T_{}`$ then $`ST_{}`$, and if $`S_{}`$ then $`S^{}_{}`$ i.e. $`_{}`$ is the least set containing $`\{i(A)|A𝖠\}`$ and closed under complements, intersections and interiors. Let $`_{}^{}`$ be the set $`\{S^{}|S_{}\}`$. We have $`_{}^{}=_{}𝒯`$. The following is the main theorem of this section.
###### Theorem 16 (Partition Theorem)
Let $`=X,𝒯,i`$ be a topological model. Then there exists a a set $`\{^\psi \}_\psi `$ of finite stable splittings such that
1. $`^\psi _{}^{}`$ and $`X^\psi `$, for all $`\psi `$,
2. if $`U^\psi `$ then $`U^\psi =\{xU|x,U\psi \}_{}`$, and
3. if $`\varphi `$ is a subformula of $`\psi `$ then $`^\varphi ^\psi `$ and $`^\psi `$ is a finite stable splitting for $`\varphi `$,
where $`_{}`$, $`_{}^{}`$ as above.
Proof. By induction on the structure of the formula $`\psi `$. In each step we take care to refine the partition of the induction hypothesis. Rather long proof.
Theorem 16 gives us a great deal of intuition for topological models. It describes in detail the expressible part of the topolocical lattice for the completeness result as it appears in Chapter 4 and paves the road for the reduction of the theory of topological models to that of spatial lattices and the decidability result of this chapter.
### 3.2 Basis Model
Let $`𝒯`$ be a topology on a set $`X`$ and $``$ a basis for $`𝒯`$. We denote satisfaction in the models $`X,𝒯,i`$ and $`X,,i`$ by $`_𝒯`$ and $`_{}`$, respectively. In the following proposition we prove that each equivalence class under $`_{}`$ contains an element of a basis closed under finite unions.
###### Proposition 17
Let $`(X,𝒯)`$ be a topological space, and let $``$ be a basis for $`𝒯`$ closed under finite unions. Let $``$ be any finite subset of $`𝒯`$. Then for all $`V`$ and all $`xV`$, there is some $`U`$ with $`xUV`$ and $`U\mathrm{𝖱𝖾𝗆}^{}V`$.
###### Corollary 18
Let $`(X,𝒯)`$ be a topological space, $``$ a basis for $`𝒯`$ closed under finite unions, $`xX`$ and $`U`$. Then
$$x,U_𝒯\varphi x,U_{}\varphi .$$
We shall prove that a model based on a topological space $`𝒯`$ is equivalent to the one induced by any basis of $`𝒯`$ which is lattice. Observe that this enables us to reduce the theory of topological spaces to that of spatial lattices and, therefore, to answer the conjecture of \[MP92\] : a completeness theorem for subset spaces which are lattices will extend to the smaller class of topological spaces.
###### Theorem 19
Let $`(X,𝒯)`$ be a topological space and $``$ a basis for $`𝒯`$ closed under finite unions. Let $`_1=X,𝒯,i`$ and $`_2=X,,i`$ be the corresponding models. Then, for all $`\varphi `$,
$$_1\varphi _2\varphi .$$
### 3.3 Finite Satisfiability
###### Proposition 20
Let $`X,𝒯`$ be a subset space. Let $``$ be a finite stable splitting for a formula $`\varphi `$ and all its subformulae, and assume that $`X`$. Then for all $`U`$, all $`xU`$, and all subformulae $`\psi `$ of $`\varphi `$, $`x,U_𝒯\psi `$ iff $`x,U_{}\psi `$.
Constructing the quotient of $`𝒯`$ under $`_{}`$ is not adequate for generating a finite model because there may still be an infinite number of points. It turns out that we only need a finite number of them.
Let $`=X,𝒯,i`$ be a topological model, and define an equivalence relation $``$ on $`X`$ by $`xy`$ iff
(a) for all $`U𝒯`$, $`xU`$ iff $`yU`$, and
(b) for all atomic $`A`$, $`xi(A)`$ iff $`yi(A)`$.
Further, denote by $`x^{}`$ the equivalence class of $`x`$, and let $`X^{}=\{x^{}:xX\}`$. For every $`U𝒯`$ let $`U^{}=\{x^{}:xX\}`$, then $`𝒯^{}=\{U^{}:U𝒯\}`$ is a topology on $`X^{}`$. Define a map $`i^{}`$ from the atomic formulae to the powerset of $`X^{}`$ by $`i^{}(A)=\{x^{}:xi(A)\}`$. The entire model $``$ lifts to the model $`^{}=X^{},𝒯^{},i^{}`$ in a well-defined way.
###### Lemma 21
For all $`x`$, $`U`$, and $`\varphi `$,
$$x,U_{}\varphi \text{iff}x^{},U^{}_{^{}}\varphi .$$
Proof. By induction on $`\varphi `$.
###### Theorem 22
If $`\varphi `$ is satisfied in any topological space then $`\varphi `$ is satisfied in a finite topological space.
Observe that the finite topological space is a quotient of the initial one under two equivalences. The one equivalence is $`_{}^\varphi `$ on the open subsets of the topological space, where $`^\varphi `$ is the finite splitting corresponding to $`\varphi `$ and its cardinality is a function of the complexity of $`\varphi `$. The other equivalence is $`_X`$ on the points of the topological space and its number of equivalence classes is a function of the atomic formulae appearing in $`\varphi `$. The following simple example shows how a topology is formed with the quotient under these two equivalences
Example: Let $`X`$ be the interval $`[0,1)`$ of real line with the the set
$$𝒯=\{\mathrm{}\}\{[0,\frac{1}{2^n})|n=0,1,2,\mathrm{}\}$$
as topology. Suppose that we have only one atomic formula, call it $`A`$, such that $`i(A)=\{0\}`$. then it is easy to see that the model $`X,𝒯,i`$ is equivalent to the finite topological model $`X^{},𝒯^{},i^{}`$, where
$$\begin{array}{c}X^{}=\{x_1,x_2\},\\ 𝒯^{}=\{\mathrm{},\{x_1,x_2\}\},\text{ and}\\ i(A)=\{x_1\}.\end{array}$$
So the overall size of the (finite) topological space is bounded by a function of the complexity of $`\varphi `$. Thus if we want to test if a given formula is invalid we have a finite number of finite topological spaces where we have to test its validity. Thus we have the following
###### Theorem 23
The theory of topological spaces is decidable.
Observe that the last two results apply for lattices of subsets by Theorem 19.
## 4 Completeness for $`\mathrm{𝐌𝐏}^{}`$
Open subsets of a topological space were used in \[MP92\] and in the previous section to provide motivation, intuition and finally semantics for $`\mathrm{𝐌𝐏}^{}`$. But in this chapter we shall show that the canonical model of $`\mathrm{𝐌𝐏}^{}`$ is actually a set of subsets closed under arbitrary intersection and finite union, i.e. the closed subsets of a topological space. However, these results are not contrary to our intuition for the following reasons: the spatial character of this logic remains untouched. The fact that the canonical model is closed under arbitrary intersections implies strong completeness with the much wider class of sets of subsets closed under finite intersection and finite union. Now, the results of the previous section allow us to deduce strong completeness (in the sense that a consistent set of formulae is simultaneously satisfiable in some model) also for the class of sets of subsets closed under infinite union and finite intersection, i.e. the open subsets of a topological space.
### 4.1 Subset frames
As we noted in section 2.1, we are not interpreting formulae directly over a subset space but, rather in the pointed product $`X\dot{\times }𝒪`$. The pointed product can be turned in a set of possible worlds of a frame. We have only to indicate what the accessibility relations are.
Definition 24 Let $`(X,𝒪)`$ be a subset space. Its subset frame is the frame
$$X\dot{\times }𝒪,R_{\mathrm{}},R_𝖪,$$
where
$$(x,U)R_{\mathrm{}}(y,V)\text{if}U=V$$
and
$$(x,U)R_𝖪(y,V)\text{if}x=y\text{and}VU.$$
If $`𝒪`$ is a topology, intended as the closed subsets of a topological space, we shall call its subset frame closed topological frame.
Our aim is to prove the most important properties of such a frame. We propose the following conditions on a possible worlds frame $`=S,R_1,R_2`$ with two accessibility relations
1. $`R_1`$ is reflexive and transitive.
2. $`R_2`$ is an equivalence relation.
3. $`R_1R_2R_2R_1`$
4. (ending points) $``$ has ending points with respect $`R_1`$, i.e
> for all $`sS`$ there exists $`s_0S`$ such that for all $`s^{}S`$ if $`sR_1s^{}`$ then $`s^{}R_1s_0`$.
5. (extensionality condition) For all $`s,s^{}S`$, if there exists $`s_0S`$ such that $`sR_1s_0`$ and $`s^{}R_1s_0`$ and
> for all $`tS`$ such that $`tR_2s`$ there exist $`t^{},t_0S`$ such that $`t^{}R_2s^{}`$, $`tR_1t_0`$ and $`t^{}R_1t_0`$, and for all $`t^{}S`$ such that $`t^{}R_2s^{}`$ there exist $`t,t_0S`$ such that $`tR_2s`$, $`t^{}R_1t_0`$ and $`tR_1t_0`$,
then $`s=s^{}`$.
6. (union condition) For all $`s_1,s_2S`$,
> if there exists $`sS`$ such that $`sR_2R_1s_1`$ and $`sR_2R_1s_2`$, then there exists $`s^{}S`$ such that for all $`tS`$ with $`tR_2s^{}`$ then $`tR_1R_2s_1`$ or $`tR_1R_2s_2`$.
7. (intersection condition) For all $`\{s_i\}_{iI}S`$,
> if there exists $`sS`$ such that $`s_iR_1s`$ for all $`iI`$ then there exists $`s^{}S`$ such that for all $`\{t_i\}S`$ with $`t_iR_2s_i`$ and $`t_iR_1t_0`$ for all $`iI`$ and some $`t_0S`$ then $`t_iR_1R_2s^{}`$.
8. The frame $``$ is strongly generated in the sense that
> there exists $`sS`$ such that for all $`s^{}S`$, $`sR_2R_1s^{}`$.
We have the following observations to make about the above conditions. Conditions 1 to 6 and 8 are first order, while the intersection condition is not. The extensionality condition implies the following
> for all $`s,s^{}S`$ such that $`sR_1s_0`$ and $`s^{}R_2s_0`$ then $`s=s^{}`$
which implies that $`R_1R_2`$ is the identity in $`S`$. In view of the extensionality condition the relation $`R_1`$ is antisymmetric. So we can replace condition 1 with the condition that $`R_1`$ is a partial order.
Now, we have the following proposition
###### Proposition 25
If $`(X,𝒯)`$ is a topological space then its closed topological frame $`_𝒯`$ satisfies conditions 1 through 8.
The above proposition could lead to the consequence that topological models are possible worlds models in disguise. But the following theorem shows that this is not the case. There is a duality.
###### Theorem 26
Let $`=S,R_1,R_2`$ be a frame satisfying conditions 1 through 8. Then $``$ is isomorphic to a closed topological frame $`_𝒯`$.
Note that, in the above definitions, we could have used equally well the equivalence class of $`sS`$ under the equivalence induced by the symmetric closure of $`R_1`$ instead of the ending point of $`s`$ in $``$. The above proofs show that the crucial conditions are conditions 1 through 5 and if we are to strengthen or relax the union and intersection conditions we get accordingly different conditions in the lattice of the set of subsets of the space. The same holds for condition 8. We only used this condition to show that there exists a top element, i.e. the whole space, and satisfy the hypothesis of the union condition. If we do not assume this condition the union of two subsets will belong to the set of subsets if they have an upper bound in it. We state this case formally without a proof because we are going to use it later.
###### Proposition 27
1. Let $`(X,𝒪)`$ be a subset space closed under infinite intersections and if $`U,V𝒪`$ have an upper bound in $`𝒪`$ then $`UV𝒪`$. Then its frame $`_𝒪`$ satisfies conditions 1 through 7.
2. A frame $``$ satisfying conditions 1 through 7 is isomorphic to a frame $`_𝒪`$ where $`(X,𝒪)`$ as in (1).
### 4.2 On the proof theory of $`\mathrm{𝐌𝐏}^{}`$
We shall identify certain classes of formulae in $``$. This approach is motivated by the results of Chapter 3. In fact, these formulae express definable parts of the lattice of subsets (see section 3.1.)
Definition 28 Let $`^{}`$ be the set of formulae generated by the following rules:
$$𝖠^{}\frac{\varphi ,\psi ^{}}{\varphi \psi ^{}}\frac{\varphi ^{}}{\neg \varphi ,\mathrm{}𝖪\varphi ^{}}$$
Let $`^{\prime \prime }`$ be the set $`\{𝖪\varphi ,𝖫\varphi |\varphi ^{}\}`$.
Formulae in $`^{}`$ have the following properties
Definition 29 A formula $`\varphi `$ of $``$ is called persistent whenever $`\varphi \mathrm{}\varphi `$ is a theorem (see also \[MP92\].)
A formula $`\varphi `$ of $``$ is called anti-persistent whenever $`\neg \varphi `$ is persistent, i.e. $`\neg \varphi \mathrm{}\neg \varphi `$ (or, equivalently $`\mathrm{}\varphi \varphi `$) is a theorem.
A formula $`\varphi `$ of $``$ is called bi-persistent whenever $`\left(\varphi \mathrm{}\varphi \right)\left(\neg \varphi \mathrm{}\neg \varphi \right)`$ (or, equivalently $`\mathrm{}\varphi \mathrm{}\varphi `$) is a theorem.
Thus the truth of bi-persistent formulae depends only on the choice of the point of the space while the satisfaction of persistent formulae can change at most once in any model. We could go on and define a hierarchy of sets of formulae where each member of hierarchy contains all formulae which their satisfaction could change at most $`n`$ times in all models.
All the following derivations are in $`\mathrm{𝐌𝐏}^{}`$ (Axioms 1 through 12 — see table at page 1.)
###### Proposition 30
All formulae belonging to $`^{}`$ are bi-persistent.
Proof. We prove it by induction, i.e. bi-persistence is retained through the application of the formation rules of $`^{}`$.
A faster (semantical) proof would be “the initial assignment on atomic formulae extends to the wider class of $`^{}`$”! This implies that formulae in $`^{}`$ define subsets of the topological space.
Formulae in $`^{\prime \prime }`$ have similar properties as the following lemma show.
###### Lemma 31
If $`\varphi `$ is bi-persistent then $`𝖪\varphi `$ is persistent and $`𝖫\varphi `$ is anti-persistent.
We prove some theorems of $`\mathrm{𝐌𝐏}^{}`$ that we are going to use later.
###### Lemma 32
If $`\varphi `$ is bi-persistent then $`_{\mathrm{𝐌𝐏}}\mathrm{}(\varphi \psi )\mathrm{}\varphi \mathrm{}\psi .`$
The following is the key lemma to the DNF Theorem and generalizes Axiom 12
###### Lemma 33
For all $`n`$,
$$_{\mathrm{𝐌𝐏}^{}}\mathrm{}𝖪\varphi \underset{i=1}{\overset{n}{}}𝖫(\mathrm{}𝖪\varphi \psi _i)\mathrm{}(𝖪\varphi \underset{i=1}{\overset{n}{}}𝖫\psi _n),$$
where $`\varphi `$, $`\psi _i`$ are bi-persistent.
All formulae of $`^{}`$ can be expressed in terms of bi-persistent, persistent and antipersistent formulae by means of the following normal form.
Definition 34
1. $`\varphi `$ is in prime normal form (PNF) if it has the form
$$\psi 𝖪\psi ^{}\underset{i=1}{\overset{n}{}}𝖫\psi _i$$
where $`\psi ,\psi ^{},\psi _i^{}`$ and $`n`$ is finite.
2. $`\varphi `$ is in disjunctive normal form (DNF) if it has the form $`_{i=1}^m\varphi _i`$, where each $`\varphi _i`$ is in PNF and $`m`$ is finite.
We now give the formal analogue of the Partition Theorem.
###### Theorem 35 (DNF)
For every $`\varphi `$, there is (effectively) a $`\psi `$ in DNF such that
$$_{\mathrm{𝐌𝐏}^{}}\varphi \psi .$$
The DNF theorem is the most important property of $`\mathrm{𝐌𝐏}^{}`$. An immediate corollary is that, as far as $`\mathrm{𝐌𝐏}^{}`$ is concerned, we could have replaced the $`\mathrm{}`$ modality with $`\mathrm{}𝖪`$, since the formulae in normal form are defined using these two modalities. Almost all subsequent proof theoretic properties are immediate or implicit corollaries of the DNF Theorem.
We close this section with the following proposition, which together with Axiom 11 shows that $`\mathrm{}\mathrm{}`$ is equivalent to $`\mathrm{}\mathrm{}`$.
###### Proposition 36
For all $`\varphi `$, $`_{\mathrm{𝐌𝐏}^{}}\mathrm{}\mathrm{}\varphi \mathrm{}\mathrm{}\varphi `$
### 4.3 Canonical Model
The canonical model of $`\mathrm{𝐌𝐏}^{}`$ is the structure
$$𝒞=(S,\{R_{\mathrm{}},R_𝖪\},v),$$
where
$$\begin{array}{ccc}& S=\{s|s\text{ is }\mathrm{𝐌𝐏}^{}\text{-maximal consistent}\},\hfill & \\ & sR_{\mathrm{}}t\text{ iff }\{\varphi |\mathrm{}\varphi s\}t,\hfill & \\ & sR_𝖪t\text{ iff }\{\varphi |𝖪\varphi s\}t,\hfill & \\ & v(A)=\{sS|AS\},\hfill & \end{array}$$
along with the usual satisfaction relation (defined inductively):
$$\begin{array}{ccc}s_𝒞A\hfill & \text{iff}\hfill & sv(A)\hfill \\ s\vDash ̸_𝒞\hfill & & \\ s_𝒞\neg \varphi \hfill & \text{iff}\hfill & s\vDash ̸_𝒞\varphi \hfill \\ s_𝒞\varphi \psi \hfill & \text{iff}\hfill & s_𝒞\varphi \text{ and }s_𝒞\psi \hfill \\ s_𝒞\mathrm{}\varphi \hfill & \text{iff}\hfill & \text{for all }tS,sR_{\mathrm{}}t\text{ implies }t_𝒞\varphi \hfill \\ s_𝒞𝖪\varphi \hfill & \text{iff}\hfill & \text{for all }tS,sR_𝖪t\text{ implies }t_𝒞\varphi .\hfill \end{array}$$
We write $`𝒞\varphi `$, if $`s_𝒞\varphi `$ for all $`sS`$.
A canonical model exists for all consistent bimodal systems with the normal axiom scheme for each modality (as MP and $`\mathrm{𝐌𝐏}^{}`$.) We have the following well known theorems (see \[Che80\], or \[Gol87\].)
###### Theorem 37 (Truth Theorem)
For all $`sS`$ and $`\varphi `$,
$$s_𝒞\varphi \text{iff}\varphi s.$$
###### Theorem 38 (Completeness Theorem)
For all $`\varphi `$,
$$𝒞\varphi \text{iff}_{\mathrm{𝐌𝐏}^{}}\varphi .$$
We shall now prove some properties of the members of $`𝒞`$. The DNF theorem implies that every maximal consistent theory $`s`$ of $`\mathrm{𝐌𝐏}^{}`$ is determined by the formulae in $`^{}`$ and $`^{\prime \prime }`$ it contains, i.e. by $`s^{}`$ and $`s^{\prime \prime }`$. Moreover, the set $`\{𝖪\varphi ,𝖫\varphi |𝖪\varphi ,𝖫\varphi s\}`$ is determined by $`s^{\prime \prime }`$ alone (this is the $`𝖪`$-case of the DNF theorem.)
The following definition is useful
Definition 39 Let $`P^{}`$. We say $`P`$ is an $`^{}`$ theory if $`P`$ is consistent and for all $`\varphi ^{}`$ either $`\varphi P`$ or $`\neg \varphi P`$.
Let $`S^{\prime \prime }`$. We say $`S`$ is an $`^{\prime \prime }`$ theory if $`S`$ is consistent and for all $`\varphi ^{\prime \prime }`$ either $`\varphi S`$ or $`\neg \varphi S`$.
Hence, $`s^{}`$ is an $`^{}`$ theory and $`s^{\prime \prime }`$ is an $`^{\prime \prime }`$ theory.
What about going in the other direction? When does an $`^{}`$ theory and $`^{\prime \prime }`$ theory determine an $`\mathrm{𝐌𝐏}^{}`$ maximal consistent theory? When their union is consistent because in this case there is a unique maximal extension. To test consistency we have the following lemma.
###### Lemma 40
If $`P`$ and $`S`$ are an $`^{}`$ and $`^{\prime \prime }`$ theory respectively then $`PS`$ is consistent if and only if
$$\text{if}\varphi P\text{then}𝖫\varphi S.$$
It is expected that since $`^{}`$ and $`^{\prime \prime }`$ theories determine $`\mathrm{𝐌𝐏}^{}`$ maximal consistent sets they will determine their accessibility relations, as well.
###### Proposition 41
For all $`s,tS`$,
$$\begin{array}{ccc}a.sR_{\mathrm{}}t\hfill & \text{if and only if}\hfill & \text{i. }\varphi t\text{ if and only if }\varphi s,\text{ where }\varphi ^{},\hfill \\ & & \text{ii. if }𝖫\varphi t\text{ then }𝖫\varphi s\text{, where }\varphi ,\psi ^{}.\hfill \\ b.sR_𝖪t\hfill & \text{if and only if}\hfill & 𝖪\varphi t\text{ if and only if }𝖪\varphi s,\text{ where }\varphi ^{}.\hfill \end{array}$$
From the above proposition we have that for all $`s,tS`$, if $`sR_{\mathrm{}}t`$ then $`s^{}=t^{}`$ and if $`sR_𝖪t`$ then $`s^{\prime \prime }=t^{\prime \prime }`$.
We write $`R_𝖪R_{\mathrm{}}`$ for the composition of the relation $`R_𝖪`$ and $`R_{\mathrm{}}`$, i.e. if $`s,tS`$, we write $`sR_𝖪R_{\mathrm{}}t`$ if there exists $`rS`$ such that $`sR_𝖪r`$ and $`rR_{\mathrm{}}t`$. Similarly for $`R_{\mathrm{}}R_𝖪`$.
For the composite relation $`R_𝖪R_{\mathrm{}}`$ and $`R_{\mathrm{}}R_𝖪`$ we have the following corollary of proposition 41
###### Corollary 42
For all $`s,tS`$,
$$\begin{array}{ccc}a.sR_{\mathrm{}}R_𝖪t\hfill & \text{if and only if}\hfill & \text{i. if }\varphi s\text{ then }𝖫\varphi t,\text{ where }\varphi ^{},\hfill \\ & & \text{ii. if }𝖫\varphi t\text{ then }𝖫\varphi s\text{, where }\varphi ,\psi ^{}.\hfill \\ b.sR_𝖪R_{\mathrm{}}t\hfill & \text{if and only if}\hfill & \text{if }𝖫\varphi t\text{ then }𝖫\varphi s,\text{ where }\varphi ^{}.\hfill \end{array}$$
We shall now prove that the canonical model $`𝒞`$ of $`\mathrm{𝐌𝐏}^{}`$ satisfies the conditions of Section 4.1 on page 4.1.
We now have the following
###### Theorem 43
The canonical model $`𝒞`$ of $`\mathrm{𝐌𝐏}^{}`$ satisfies conditions 1 to 7 of Section 4.1 on page 4.1.
###### Corollary 44
The canonical frame of $`\mathrm{𝐌𝐏}^{}`$ is isomorphic to a subset frame $`_{𝒪_c}`$ where $`(X_c,𝒪_c)`$ is a subset space closed under infinite intersections and if $`U,V𝒪_c`$ have an upper bound in $`𝒪_c`$ then $`UV𝒪_c`$.
By the construction of Theorem 26, $`X_c`$ consists of the ending points of the members of the domain of the canonical model. We define the following initial assignment $`i_c`$
$$i(A)=\{s_0|As_0\}.$$
In this way the model $`=X_c,𝒪_c,i_c`$ is equivalent to the canonical model as a corollary of frame isomorphism.
###### Corollary 45
For all $`sS`$ and $`\varphi `$ we have
$$\varphi s\text{if and only if}s_0,U_s_{}\varphi .$$
###### Proposition 46
The frame of a generated submodel $`𝒞^t`$ is isomorphic to a closed topological frame.
Now as above we have the following corollary
###### Corollary 47
A submodel $`𝒞^t`$ is equivalent to a closed topological model.
It is a well known fact that a modal system is characterized by the class of generated frames of the canonical frame.
###### Proposition 48
The system $`\mathrm{𝐌𝐏}^{}`$ is (strongly) characterized by closed topological frames.
Since the axioms and rules of $`\mathrm{𝐌𝐏}^{}`$ are sound for the wider class of subset spaces with finite union and intersection, we also have the following.
###### Proposition 49
The system $`\mathrm{𝐌𝐏}^{}`$ is (strongly) characterized by subset frames closed under finite unions and intersections.
Now by Proposition 48 and 49, Corollary 18 and Theorem 19 of Chapter 3, where we proved the equivalence of a topological model with the model induced by a basis closed under finite unions, we have the following corollary
###### Corollary 50
The system $`\mathrm{𝐌𝐏}^{}`$ is (strongly) characterized by open topological frames as well as subset frames closed under infinite unions and intersections.
THe following disjunction property holds for $`\mathrm{𝐌𝐏}^{}`$
$$\text{if}_{\mathrm{𝐌𝐏}^{}}𝖪\varphi _1𝖪\varphi _2\mathrm{}𝖪\varphi _n\text{then}_{\mathrm{𝐌𝐏}^{}}\varphi _i,\text{for some }i,1in,$$
for $`\varphi _1,\varphi _2,\mathrm{},\varphi _n^{}`$. Note that the disjunction rule does not hold for $`\mathrm{𝐒𝟓}`$.
###### Proposition 51
$`\mathrm{𝐌𝐏}^{}`$ provides the above rule of disjunction.
We can similarly prove a stronger disjunction property, namely
$$\begin{array}{cc}\text{if}\hfill & _{\mathrm{𝐌𝐏}^{}}𝖪\varphi 𝖪\varphi _1𝖪\varphi _2\mathrm{}𝖪\varphi _n\hfill \\ \text{then}\hfill & _{\mathrm{𝐌𝐏}^{}}𝖪\varphi \varphi _i,\text{for some }i,1in,\hfill \end{array}$$
for $`\varphi ,\varphi _1,\varphi _2,\mathrm{},\varphi _n^{}`$.
Now we are able to prove the following
###### Theorem 52
The canonical model of $`\mathrm{𝐌𝐏}^{}`$ is strongly generated.
By Theorem 52 we complete the set of conditions of page 1 which turn the frame of the canonical model into a closed subset frame. To summarize, we have the following corollary (note that the canonical subset model is $`X_c,𝒪_c,i_c`$ of Corollary 44)
###### Corollary 53
The canonical subset model of $`\mathrm{𝐌𝐏}^{}`$ is a topological space.
## 5 The Algebras of $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$
In this section we shall give a more general type of semantics for $`\mathrm{𝐌𝐏}`$ and $`\mathrm{𝐌𝐏}^{}`$. Every modal logic can be interpreted in an algebraic framework. An algebraic model is nothing else but a valuation of the propositional variables in a class of appropriately chosen algebras. We shall also make connections with the previous chapters.
### 5.1 Fixed Monadic Algebras
Interior operators were introduced by McKinsey and Tarski \[MT44\].
Definition 54 An interior operator $`I`$ on a Boolean algebra $`=B,\mathrm{𝟎},\mathrm{𝟏},,`$ is an operator satisfying the conditions
$$\begin{array}{c}I(ab)=IaIb,\\ Iaa,\\ IIa=Ia,\\ I\mathrm{𝟏}=\mathrm{𝟏}.\end{array}$$
To each interior operator $`I`$ we associate its dual $`C=I`$, called closure operator.
Universal quantifiers were introduced by P. Halmos \[Hal56\].
Definition 55 A universal quantifier $``$ on a Boolean algebra $``$ is an operator satisfying the conditions
$$\begin{array}{c}(ab)=ab,\\ aa,\\ \mathrm{𝟏}=\mathrm{𝟏}.\end{array}$$
To each universal quantifier $``$ we associate its dual $`=`$, called existential quantifier.
Definition 56 Let $`I`$ be an interior operator on a Boolean algebra $``$. Let $`IB=\{a|aIa\}`$ and $`CB=\{a|Caa\}`$, i.e. the fixed points of $`I`$ and $`C`$ respectively. Let $`B^I=IBCB`$ then $`^I=B^I,\mathrm{𝟎},\mathrm{𝟏},,,`$ is a Boolean subalgebra of $``$.
Definition 57 A fixed monadic algebra (FMA) $``$ is a Boolean algebra with two operators $`I`$ and $``$ satisfying
$$IaIa.$$
A valuation $`v`$ on $``$ is a function from the formulae of MP to the elements of $`B`$ such that
$$\begin{array}{ccc}\hfill v(A)& \hfill & B^I,\text{ where }A\text{ is atomic},\hfill \\ \hfill v(\neg \varphi )& =\hfill & v(\varphi ),\hfill \\ \hfill v(\varphi \psi )& =\hfill & v(\varphi )v(\varphi ),\hfill \\ \hfill v(\varphi \psi )& =\hfill & v(\varphi )v(\varphi ),\hfill \\ \hfill v(\mathrm{}\varphi )& =\hfill & Iv(\varphi ),\hfill \\ \hfill v(𝖪\varphi )& =\hfill & v(\varphi ).\hfill \end{array}$$
An algebraic model of MP is a FMA $``$ with a valuation $`v`$ on it. We say $`\varphi `$ is valid in this model iff $`v(\varphi )=\mathrm{𝟏}`$ and valid in an FMA iff it is valid in all models based on this algebra. Finally, $`\varphi `$ is FMA-valid if it is valid in all FMA’s. The notion of validity can extend to a set of formulae.
Observe that the important part of the algebra is the smallest subalgebra containing $`B^I`$ and closed under the operators $`I`$ and $``$.
###### Theorem 58 (Soundness for FMA-validity)
If a formula $`\varphi `$ is a theorem of MP then $`\varphi `$ is FMA-valid.
###### Theorem 59 (Completeness for FMA-validity)
If $`\varphi `$ is FMA-valid then $`\varphi `$ is a theorem of $`\mathrm{𝐌𝐏}`$.
### 5.2 Generated Monadic Algebras
We shall now define the algebraic models of $`\mathrm{𝐌𝐏}^{}`$
Definition 60 A generated monadic algebra (GMA) $``$ is an FMA satisfying in addition
$$\begin{array}{c}CIa=ICa\\ C(ab)C(ac)C(CaCbCc).\end{array}$$
The concepts of algebraic model, validity, GMA-validity are defined as for FMA’s. We used the direct algebraic translation of $`\mathrm{𝐌𝐏}^{}`$ axioms but we could have defined it with a different presentation. Observe that we only need $`CIaICa`$ because the other direction is derivable (see Proposition 36.)
We now have the following
###### Theorem 61 (Algebraic completeness of $`\mathrm{𝐌𝐏}^{}`$)
A formula $`\varphi `$ is a theorem of $`\mathrm{𝐌𝐏}^{}`$ if and only if $`\varphi `$ is GMA-valid.
It is known that a modal algebra determines a (general) frame (see \[BS84\].) So, in our case, the canonical algebraic model of $`\mathrm{𝐌𝐏}^{}`$, i.e. its Lindenbaum algebra, must determine a closed topological model (actually its canonical frame.) We shall state only the interesting part of this correspondence: the bijection on the domains. The accessibility relations are defined in the usual way.
###### Theorem 62
There is a bijection between the set of the ultrafilters of the canonical algebra of $`\mathrm{𝐌𝐏}^{}`$ and the pointed product $`X\dot{\times }𝒯`$, where $`(X,𝒯)`$ is the canonical topology of $`\mathrm{𝐌𝐏}^{}`$.
The general theory of modal logic provides for yet another construction. A frame determines a modal algebra. In case of the canonical frame, the modal algebra determined must be isomorphic to the canonical modal algebra. In our case, this algebra (which must be a GMA) has a nice representation. It is the algebra of partitions of the topological lattice as it appeared in Section 3.1.
## 6 Further Directions
There are several further directions
1. Due to the indeterminacy assumption (see Introduction) $`\mathrm{𝐌𝐏}^{}`$ can be a “core” logical system for reasoning about computation with approximation or uncertainty.
2. A discrete version of our epistemic framework can arise in scientific experiments or tests. We acquire knowledge by “a step-by-step” process. Each step being an experiment or test. The outcome of such an experiment or test is unknown to us beforehand, but after being known it restricts our attention to a smaller set of possibilities. A sequence of experiments, test or actions comprises a strategy of knowledge acquisition. This model is in many respects similar to Hintikka’s “oracle” (see \[Hin86\].) In Hintikka’s model the “inquirer” asks a series of questions to an external information source, called “oracle”. The oracle answers yes or no and the inquirer increases her knowledge by this piece of additional evidence. This framework can be expressed by adding actions to the language. Preliminary work of ours used quantales for modelling such processes. A similar work without knowledge considerations appears in \[AV90\].
3. Since we can express concepts like affirmative or refutative assertions, which are closed under infinite disjunctions and conjunctions respectively, it is very natural to add infinitary connectives or fixed points operators (the latter as a finite means to express the infinitary connectives.) This would serve the purpose of specifying such properties of programs as “emits an infinite sequence of ones” (see \[Abr91\] for a relevant discussion.) An interesting direction of linking topological spaces with programs can be found in \[Par83\].
4. Our work in the algebras of $`\mathrm{𝐌𝐏}^{}`$ looks very promising. GMAs (see 5) have very interesting properties. A subalgebra of a GMA corresponds to a complete space and this duality can be further investigated with the algebraic machinery of modal logic (see \[Lem66a\], \[Lem66b\], \[Blo80\]) or category theoretic methods.
5. Axiom 10 forces monotonicity in our systems. If we drop this axiom, an application of effort no longer implies a further increase in our knowledge. Any change of our state of knowledge is possible. A non-monotonic version of the systems presented in this thesis can be given along the lines of \[Par87a\].
6. It would be interesting to consider a framework of multiple agents. Adding a modality $`𝖪_i`$ for each agent $`i`$ and assigning a different set of subsets or topology to to each agent we can study their interaction or communication by set-theoretic or topological means.
7. From our work became clear that both systems considered here are linked with intuitionistic logic. We have embed intuitionistic logic to $`\mathrm{𝐌𝐏}`$ or $`\mathrm{𝐌𝐏}^{}`$ and it would be interesting to see how much of the expressiveness of these logics can be carried in an intuitionistic framework.
8. Finally, in another direction Rohit Parikh considers an enrichment of the language to express more (and purely) topological properties such as separation properties and compactness.
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# Semiclassical theory of shot noise in disordered SN contacts
## Abstract
We present a semiclassical theory of shot noise in diffusive superconductor \- normal metal contacts. At subgap voltages, we reproduce the doubling of shot noise with respect to conventional normal-metal contacts, which is interpreted in terms of an energy balance of electrons. Above the gap, the voltage dependence of the noise crosses over to the standard one with a voltage-independent excess noise. The semiclassical description of noise leads to correlations between currents at different electrodes of multiterminal SN contacts which are always of fermionic type, i.e. negative. Using a quantum extension of the Boltzmann - Langevin method, we reproduce the peculiarity of noise at the Josephson frequency and obtain an analytical frequency dependence of noise at above-gap voltages.
PACS numbers: 72.70.+m, 74.40+k, 74.50+r
In recent years, the noise properties of hybrid contacts involving superconducting (S) and normal (N) metals attracted considerable attention. Theoretical work in this direction was pioneered by Khlus, who found that Andreev reflection, in which electrons incoming from normal metals are reflected from the SN interface as holes, play a key role in the shot noise of SN structures. For clean SN contacts, he predicted that the shot noise vanishes at subgap voltages $`eV<\mathrm{\Delta }`$. de Jong and Beenakker addressed dirty SN contacts and found that the shot noise at subgap voltages is doubled with respect to its value in a normal contact with the same resistance. This doubling of shot noise has already been experimentally confirmed. The distribution function of current fluctuations was calculated and it was found that at subgap voltages, it describes independent transfers of discrete charge $`2e`$ through the contact. Lesovik et al. found that the frequency dependence of the shot noise exhibits a peculiarity at the Josephson frequency $`\omega =2eV`$ instead of $`\omega =eV`$ in normal contacts. More recently, the same authors obtained additional singularities in the frequency dependence of the noise at above-gap voltages. These results could lead to the impression that the motion of electrons and holes is correlated not only at the SN interface, but that the current is carried also through the normal part of the contact in portions of $`2e`$ instead of $`e`$. Recent findings that the correlations of currents at two normal contacts attached to a superconductor may be positive, as in the case of Bose statistics seemed to reinforce such a view. On the other hand, from these works it is not clear whether the above results require phase coherence of electrons nor if there are restrictions for the contact length.
Recently, Gramespacher and one of the authors discussed the fluctuations in multiterminal SN contacts using electronlike and holelike distribution functions, which were expressed in terms of quantum-mechanical injectivities of the terminals. It was shown there, in particular, that positive correlations vanish in the limit of a large channel number. In this paper, we present a semiclassical theory of shot noise in disordered SN contacts. We show that the doubling of shot noise may be obtained within a simple Boltzmann - Langevin approach provided that the appropriate boundary conditions for the average distribution function are used. We relate the increase in the shot noise to the peculiar energy transport through an SN interface and show that this effect is stable with respect to phase-breaking. From the semiclassical nature of electron transport in the N region, it immediately follows that the correlations between currents in different electrodes of a multiterminal diffusive SN contact are always negative. We are also able to explain the peculiarities of the frequency dependence of noise in terms of the shape of the electron distribution function.
Consider a narrow normal-metal microbridge connecting massive normal and superconducting electrodes. The
elastic mean free path $`l`$ of electrons in the microbridge is assumed to be short and the length of the contact $`L`$ is assumed to be much larger than $`(\mathrm{}D/\mathrm{\Delta })^{1/2}`$, where $`\mathrm{\Delta }`$ is the superconducting gap and $`D`$ is the diffusion coefficient in the normal metal. The applied voltage or temperature are considered to be much larger than the Thouless energy $`\epsilon _T=\mathrm{}D/L^2`$. Under these conditions, it is possible to neglect the penetration of the condensate into the microbridge and consider it just as a normal metal with a nonequilibrium distribution function of electrons.
Let the $`x`$ axis be directed along the contact, $`x=0`$ corresponding to the superconducting electrode, and $`x=L`$, to the normal-metal one. The electron energy $`\epsilon `$ in the normal metal is measured from the Fermi level of the superconductor. Introduce an energy-dependent coefficient $`W`$ of normal transmission of an electron from the contact through the NS interface without being Andreev reflected, so that $`W(\epsilon )=0`$ if $`|\epsilon |<\mathrm{\Delta }`$ and $`W(\epsilon )1`$, for $`|\epsilon |\mathrm{\Delta }`$.
Since the electrons in the superconducting electrode are described by an equilibrium Fermi distribution function $`f_0(\epsilon )`$, the distribution function of electrons moving into the contact from the SN interface may be written in the form
$`f(\epsilon ,v_x)=Wf_0(\epsilon )+(1W)[1f(\epsilon ,v_x)],`$ (1)
for $`\epsilon >0`$ and $`v_x>0`$, and the distribution function of holes moving into the contact from the interface may be written in the form
$`1f(\epsilon ,v_x)=W[1f_0(\epsilon )]+`$ $`(1W)f(\epsilon ,v_x),`$ (2)
for $`\epsilon <0`$ and $`v_x>0`$. Since in the bulk of the contact the electron motion is diffusive, we must now obtain a suitable boundary condition at the NS interface for the diffusion equation.
Introduce the symmetric and antisymmetric parts of $`f`$ with respect to $`v_x`$:
$`f_s(\epsilon )=(1/2)\left[f(\epsilon ,v_x)+f(\epsilon ,v_x)\right],`$ (3)
$`f_a(\epsilon )=(1/2)\left[f(\epsilon ,v_x)f(\epsilon ,v_x)\right].`$ (4)
Inside the contact, the distribution function obeys the diffusion equation and may be represented as a sum of an isotropic part $`f(\epsilon ,x)`$ and a small anisotropic part
$$f_1=(v_x/v_F)(l/L)[f(\epsilon ,L)f(\epsilon ,0)].$$
(5)
Following Kupriyanov and Lukichev, we equate now $`f(\epsilon ,0)`$ to $`f_s(\epsilon )`$ and $`f_1`$ to $`(v_x/v_F)f_a(\epsilon )`$. Introducing the distribution function of electrons in the normal electrode $`f_n(\epsilon )`$, one may express $`f_a`$ in terms of $`f_s`$ and $`f_n`$:
$$f_a(\epsilon )=(l/L)[f_s(\epsilon )f_n(\epsilon )].$$
(6)
For the sake of definiteness, assume now $`\epsilon >0`$. By substituting Eqs. (4) into Eqs. (1) for $`f(\epsilon ,v_x)`$ and (2) for $`1f(\epsilon ,v_x)`$ and excluding $`f_a`$ from them by means of (6), one obtains a closed set of equations for $`f_s(\epsilon )`$ and $`f_s(\epsilon )`$:
$$\left(1+l/L\right)f_s(\epsilon )+(1W)\left(1l/L\right)f_s(\epsilon )=1Wf_0(\epsilon )$$
$$+(l/L)f_n(\epsilon )(l/L)(1W)f_n(\epsilon )$$
(7)
and the equation obtained from it by interchanging $`\epsilon `$ and $`\epsilon `$. The solution of the set (7) assumes different forms depending on the relationship between $`W`$ and $`l/L`$. In the subgap region $`|\epsilon |<\mathrm{\Delta }`$, $`W=0`$ and we obtain
$$f_s(\epsilon )=\frac{1}{2}\left[1+f_n(\epsilon )f_n(\epsilon )\right]$$
(8)
for both signs of $`\epsilon `$. Suppose that the normal electrode has a potential $`V`$. Then $`f_n(\epsilon )=f_0(\epsilon eV)`$, and inside the gap,
$$f_s(\epsilon )=\frac{1}{2}\left[f_0(\epsilon eV)+f_0(\epsilon +eV)\right],|\epsilon |<\mathrm{\Delta }.$$
(9)
Outside the subgap region, $`W0`$, and $`l/L`$ may be neglected with respect to it virtually at all energies. In this case, one easily obtains
$$f_s(\epsilon )=f_0(\epsilon ),|\epsilon |>\mathrm{\Delta }.$$
(10)
Inside the contact, the distribution function is obtained by solving the diffusion equation with the boundary conditions (9) and (10). As a result, one obtains:
$$f(\epsilon ,x)=\{\begin{array}{cc}\frac{1}{2}\left(1+\frac{x}{L}\right)f_0(\epsilon eV)+\frac{1}{2}\left(1\frac{x}{L}\right)f_0(\epsilon +eV),\hfill & |\epsilon |<\mathrm{\Delta }\hfill \\ \text{}\left(1\frac{x}{L}\right)f_0(\epsilon )+\frac{x}{L}f_0(\epsilon eV),\hfill & |\epsilon |>\mathrm{\Delta }.\hfill \end{array}$$
(11)
At subgap voltages, the shape of the distribution function is exactly the same as if the contact were extended in negative direction from 0 to $`L`$ and the voltage $`V`$ were applied at this point. The distribution function for an above-gap voltage is shown in Fig. 1. Note that at finite voltages, $`f(\epsilon )`$ is discontinuous at $`|\epsilon |=\mathrm{\Delta }`$ even at nonzero temperature.
Consider now the noise of the contact. Since all the voltage drop takes place inside the normal-metal microbridge and there is no voltage drop across the SN interface, it is reasonable to assume that all the noise is produced by random impurity scattering in the microbridge only. In this case, we can use the well known Langevin equations for the current fluctuations,
$$\delta 𝐣=D\frac{}{𝐫}\delta \rho \sigma \frac{}{𝐫}\delta \varphi +\delta 𝐣^{ext},$$
(12)
where $`D`$ is the diffusion coefficient, $`\sigma `$ is the electric conductivity, $`\delta \rho (𝐫)`$ is the charge-density fluctuation, $`\delta \varphi (𝐫)`$ is the local fluctuation of the electric potential, and the correlator of extraneous currents $`\delta 𝐣^{ext}`$ is given by
$$\delta j_\alpha ^{ext}(𝐫_1)\delta j_\beta ^{ext}(𝐫_2)_\omega =4\sigma \delta _{\alpha \beta }\delta (𝐫_1𝐫_2)T_N(𝐫_1),$$
$$T_N(𝐫)=𝑑\epsilon f(\epsilon ,𝐫)[1f(\epsilon ,𝐫)].$$
(13)
At sufficiently low frequencies, the fluctuation of the total current $`\delta 𝐣`$ may be considered as constant along the contact, and Eqn. (12) may be integrated over its length. As a result, the gradient terms in (12) drop out because of the boundary conditions $`\delta \rho =0`$ and $`\delta \varphi =0`$, and one arrives at the standard expression for the noise in disordered metal contacts:
$$S_I=\frac{4}{R}\frac{dx}{L}T_N(x),$$
(14)
where $`R`$ is the normal resistance of the contact. Substituting now the distribution function (11) into (14) and performing the integration, one obtains that
$$S_I=4\frac{T}{R}\{\frac{2}{3}+\frac{1}{3}\frac{eV}{T}\mathrm{coth}\left(\frac{eV}{T}\right)+\frac{1}{6}[\mathrm{tanh}\left(\frac{\mathrm{\Delta }+eV}{2T}\right)+\mathrm{tanh}\left(\frac{\mathrm{\Delta }eV}{2T}\right)2\mathrm{tanh}\left(\frac{\mathrm{\Delta }}{2T}\right)]$$
$$+\frac{1}{6}[\mathrm{coth}\left(\frac{eV}{2T}\right)2\mathrm{coth}\left(\frac{eV}{T}\right)]\mathrm{ln}\left[\frac{\mathrm{exp}(\mathrm{\Delta }/T)+\mathrm{exp}(eV/T)}{\mathrm{exp}(\mathrm{\Delta }/T)+\mathrm{exp}(eV/T)}\right]\}.$$
(15)
At zero voltage, this expression reduces to the Nyquist formula $`S_I=4T/R`$. At zero temperature yet finite voltage, Eqn. (15) takes a very simple form:
$$S_I=\frac{4}{3}\frac{e|V|}{R}+\frac{2}{3}\theta (e|V|\mathrm{\Delta })\frac{\mathrm{\Delta }e|V|}{R}.$$
(16)
The shot noise in diffusive SN contacts is doubled at voltages $`e|V|<\mathrm{\Delta }`$ as compared to the same contacts between normal metals. The reasons for the shot-noise doubling can be understood in terms of the noise temperature of electrons $`T_N`$ and the energy transport. The electrons in the contact acquire additional energy because of their acceleration by the field. If both electrodes are made of normal metal, the excess energy flows into both electrodes thus effectively cooling the contact and decreasing $`T_N`$. However if one of the electrodes is superconducting, the SN boundary forbids heat transfer in the subgap region. The absence of heat transfer through an SN boundary has been known for many years, and this is precisely what motivated the pioneering paper by Andreev. Because there is only one energy drain now instead of two, naturally $`T_N`$ is higher than in the case of a normal contact.
The above semiclassical derivation shows that the doubling of shot noise is insensitive to phase-breaking. In the case of strong electron-electron scattering, the effective-temperature approximation may be used at $`eV\mathrm{\Delta }`$ and precisely the same heat-balance equation for the effective temperature $`T_e`$ as in normal contacts can be written except that the heat-absorption boundary condition $`T_e=T`$ at the SN interface should be replaced by the zero heat-flux condition $`T_e/x=0`$. It is easily seen that this should double the shot noise making it $`2\times (\sqrt{3}/2)eI=\sqrt{3}eI`$ and thus increasing it above the noninteracting-electron value precisely as for a normal contact. Much like in the case of normal contacts, the noise should be suppressed by energy relaxation caused e. g. by phonon emission.
Using the semiclassical approach, Sukhorukov and Loss recently investigated the correlations between currents in different contacts of multiterminal diffusive normal-metal systems. Here we extend their discussion to the case of multiterminal SN systems. It is convenient to introduce characteristic potentials $`\varphi _n(𝐫)`$, $`n=1,\mathrm{},N`$ associated whith each contact $`n`$ which equal 1 at contact $`n`$ and are zero at all other electrodes obeying the equation
$$(\sigma \varphi _n)=0$$
(17)
in the bulk of the contact. In terms of these potentials, the normal-state conductance matrix of the contact may be written in the form
$$G_{mn}=𝑑𝐒_m\sigma \varphi _m\varphi _n,$$
(18)
where $`S_m`$ is the interface of the system with electrode $`m`$, and the cross-correlated spectral density for contacts $`n`$ and $`m`$ may be written in the form
$$S_{mn}=4d^3𝐫\varphi _n\varphi _m\sigma T_N(𝐫).$$
(19)
Performing twice an integration by parts and making use of (17) and a similar equation for $`f`$ in $`T_N`$, equation (19) is easily brought to a form
$$S_{mn}=2G_{mn}T_N|_{S_m}2G_{nm}T_N|_{S_N}$$
$$4d^3r\varphi _n\varphi _m\sigma 𝑑\epsilon (f)^2.$$
(20)
As $`T_N`$, $`G_{mn}`$, and $`\varphi _m`$ are always positive, this proves that the current correlations are always negative, i. e. of fermionic type. As this proof holds for an arbitrary shape of the distribution function at the interfaces with electrodes, the correlations are negative for an arbitrary number of superconducting electrodes provided that the Josephson effect between them is suppressed.
Consider now the noise at frequencies $`\omega eV`$. This is the range of quantum noise, which clearly cannot be described in terms of the Boltzmann equation. However it was recently shown that for disordered metals, the Langevin scheme may be extended to high frequencies by introducing a frequency-dependent noise temperature
$$T_N(𝐫,\omega )=\frac{1}{2}d\epsilon \{f(𝐫,\epsilon )[1f(𝐫,\epsilon +\omega )]$$
$$+f(𝐫,\epsilon +\omega )[1f(𝐫,\epsilon )]\}$$
(21)
into Eq. (13). If the frequency is smaller than $`1/RC`$, where $`C`$ is the capacity of the contact to all possible external gates, the current fluctuation may still be considered as constant along the contact length, and we arrive again at Eq. (14) with $`T_N(x,\omega )`$ substituted for $`T_N(x)`$. Note that the latter condition is well compatible with $`\omega D/L^2`$ because $`1/RC`$ is larger by an extra factor containing the ratio of the contact diameter to the Thomas - Fermy screening length.
Suppose first that $`eV<\mathrm{\Delta }`$. Integration in (21) and (14) then gives
$$S_I(\omega )=\frac{1}{3R}[4U(\omega )+U(2eV+\omega )+U(2eV\omega )],$$
(22)
with $`U(\omega )=\omega \mathrm{coth}(\omega /2T)`$, which reduces to
$$S_I(\omega )=(1/3R)[4\omega +\mathrm{max}(4eV,2\omega )]$$
(23)
at zero temperature. The frequency dependence of $`S_I`$ exhibits a kink at $`\omega =2eV`$, which would correspond to the Josephson frequency of the contact if the second electrode were also superconducting. The reason for the doubling of the kink frequency with respect to normal contacts is that the range of energies where the drop of the distrbution function takes place is two times larger now.
In the case where $`T=0`$ and $`|eV|>\mathrm{\Delta }`$, integration gives
$$S_I(\omega )=\frac{|\omega |}{R}+\frac{1}{6R}(|\omega +2\mathrm{\Delta }|+|\omega 2\mathrm{\Delta }|+|\omega eV+\mathrm{\Delta }|$$
$$+|\omega eV\mathrm{\Delta }|+|\omega +eV+\mathrm{\Delta }|+|\omega +eV\mathrm{\Delta }|).$$
(24)
in accordance with recent numerical results, this expression exhibits peculiarities at $`\omega =2\mathrm{\Delta }`$ and $`\omega =eV\pm \mathrm{\Delta }`$, which are related with discontinuities in the energy dependence of the distribution function. The slope of $`S(\omega )`$ increases from $`1/R`$ at $`\omega 0`$ to $`2/R`$ at $`\omega \mathrm{}`$.
In this work we have developped a semi-classical approach for the noise in metallic diffusive NS-structures. This approach explains a number of key results in a physically transparent manner and demonstrates that these results are immune to dephasing.
This work was supported by the Swiss National Science Foundation.
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# AN EFFICIENT MULTIPROCESSOR MANAGEMENT SYSTEM FOR EVENT–PARALLEL COMPUTING
Proceedings Preprint, DPF ’96, Minneapolis, Minnesota (10–15 August 1996)
Performance of software using TCP/IP sockets to distribute events to UNIX workstations is described. This simple software was written at the University of Mississippi to control UMiss farm reconstruction of 8 billion raw events, part of Fermilab E791’s data. E791 reconstructed HEP’s largest data set to study charm physics.
Fermilab E791 wrote a big dataset (50 Terabytes, 20 billion events, 24 000 8mm Exabyte tapes) in 1991 and early 1992.<sup>a</sup><sup>a</sup>footnotemark: a<sup>b</sup><sup>b</sup>footnotemark: b<sup>c</sup><sup>c</sup>footnotemark: c Reconstruction challenged available computing, requiring over 10<sup>4</sup> mips-years. The task was larger than at colliders (Table 1). Reconstruction was nevertheless completed using four farm sites. <sup>d</sup><sup>d</sup>footnotemark: d Here we describe the multiprocessor management software <sup>e</sup><sup>e</sup>footnotemark: e developed and run at the University of Mississippi farm (Figs. 1 and 2 show hardware).
HEP events are usually independent. Interprocess I/O isn’t needed. An efficient parallel system just has to input and output events fast enough so clients are never idle. Management software had to do a lot of hard work in early HEP systems.<sup>f</sup><sup>f</sup>footnotemark: f<sup>g</sup><sup>g</sup>footnotemark: g Clients had minimal operating systems. All data had to be formatted in a server and downloaded into clients word by word. Moving from single to multiple CPUs was hard; the division between server and
client code was intricate. With the advent of commercial workstation clients <sup>h</sup><sup>h</sup>footnotemark: h with real operating systems, most work inherent in moving to multiprocessors vanished. Using Network File System software, server
disks can be cross-mounted so that files are accessible by multiple clients. In this model, even inexpensive diskless clients directly read an executable code file, a run number file, calibration files, a raw input record file of events, and write report files and reconstructed event files. The server writes input events from tape to disk files. At the end of a job, the server copies client output event files to tape and combines client reports, as clients work on the next job. Because 85% of E791 events were filtered away after reconstruction, disk output was fast enough for us. Event input by disk also worked, but too slowly. So, our multiprocessor manager bypasses disk for input using instead Transmission Control Protocol/ Internet Protocol.<sup>i</sup><sup>i</sup>footnotemark: i With TCP/IP, processes make a
connection between themselves and pass data back and forth using read\_from\_connection and write\_through\_connection subroutine calls. A test of TCP/IP gave 900 kbyte/s, ending client idleness. Fig. 3 illustrates how the network I/O calls are used. As the server prepares to start a client, it uses make\_socket to “have a phone put in”, so that it will be able to connect to the client. When the client starts, it too uses make\_socket to “have a phone put in”. The server “lists its number” by binding its socket to a port (bind\_socket), and “stays near the phone” listening for an attempt to connect (listen\_socket). The client “calls up” the server (connect\_socket) and the server “picks up the phone” establishing the connection (accept\_socket). When it needs input data, the client “places its order” by writing a message to the server (write\_socket). The server is continually monitoring all of the client connections for requests (select\_socket). When a request comes in, the
server “writes down the order” (read\_socket), and
does its best to satisfy the client’s request. The client and server shuttle messages back and forth (each writing to the other and reading from the other) until the input data is exhausted. Next, the server notifies the clients that there is no more data.
The clients then finish their tasks, close their connections (close\_socket), and exit; the server finishes its tasks and exits. The Fortran-callable C routines that manipulate sockets and connections really are that simple to use. Only read\_socket is more than a C to Fortran interface, and even its trivial. Most of the real work has already been done in UNIX, TCP/IP, and Berkeley Sockets.
Although most of our farm processors are in racks, some are on people’s desks. We have found it satisfactory to allow users to abruptly kill the client process whenever they find its activities on their workstations to be troublesome. A reconstruction code crash also kills a client. In either case, the server is quickly aware of the dead client and adjusts event distribution. A disadvantage of this approach is that a few input events are trapped and lost. Having 20 billion events, we take a rather cavalier attitude. In E791, processing raw events is rather like hauling corn to market in a truck. If a few grains of corn fall out of the truck, no one cares. The alternative approach – treating events as babies in a hospital nursery, where one normally expects a somewhat stricter accounting – only makes sense later with small selections of interesting events.
Before writing our own multiprocessor software we considered extracting the few features needed from large packages under development at Fermilab <sup>j</sup><sup>j</sup>footnotemark: j<sup>k</sup><sup>k</sup>footnotemark: k and Argonne.<sup>l</sup><sup>l</sup>footnotemark: l However, offsite support was unavailable. So we focused on writing software which could do a limited number of things very well; e.g. run many clients per server efficiently and tolerate client crashes and operating system upgrades. Six man–weeks were spent coding. Farm operation required 5 hours over a day. After seeing our approach to moving farm data, the D0 experiment decided to follow a similar strategy.<sup>m</sup><sup>m</sup>footnotemark: m
Funding in June 1993 allowed an expansion of the UMiss farm from 1100 to 2900 mips. By July 1993, the increased computing was acquired and processing data. E791 reconstruction was completed in Sept. 1994. A total of 8 billion events on 10 000 raw data tapes were processed in Mississippi. Before running final reconstruction, dozens of full farm tests of algorithms for actual charm yield were run, each test for a few days. The charm yield tripled. X Window operator control displays written in Tcl/Tk aided bookkeeping. Tape reading was multiply buffered, so that events were almost always available immediately when a client asked for them. During smooth running, timing CPUs showed that at least 97% of client processing cycles were used. Overall efficiency, considering cycles lost for any reason, exceeded 90% over a 2$`\frac{1}{2}`$ year period.
Efficient management of multiple processors has led to the reconstruction of 200 000 charm particles, the world’s largest sample. Results <sup>n</sup><sup>n</sup>footnotemark: n include DPF ’96 papers by N. Copty, L. Cremaldi, K. Gounder, M. Purohit, K.C. Peng, A. Tripathi, R. Zaliznyak, and C. Zhang. We especially thank Lucien Cremaldi and Breese Quinn for their contributions to building and running the UMiss farm. This work was supported in part by U.S. DOE DE-FG05-91ER40622. <sup>1</sup><sup>1</sup>footnotetext: S. Amato et al., Nucl. Instr. Meth. A324 (1993) 535; E791 DA. <sup>2</sup><sup>2</sup>footnotetext: D.J. Summers et al., XXVIIth Rencontre de Moriond, Les Arcs (15-22 March 1992) 417. <sup>3</sup><sup>3</sup>footnotetext: B.R. Kumar, Vertex Detectors, Plenum Press, Erice (21-26 September 1986) 167. <sup>4</sup><sup>4</sup>footnotetext: CBPF–Rio de Janeiro, Fermilab, Kansas State, Mississippi (CREMF). <sup>5</sup><sup>5</sup>footnotetext: Steve Bracker et al., IEEE Trans. Nucl. Sci. NS-43 (October 1996). <sup>6</sup><sup>6</sup>footnotetext: Paul F. Kunz et al., IEEE Trans. Nucl. Sci. NS-27 (1980) 582; IBM 168 emulator. <sup>7</sup><sup>7</sup>footnotetext: J. Biel et al., Computer Physics Communications 45 (1987) 331; ACP. <sup>8</sup><sup>8</sup>footnotetext: C. Stoughton and D.J. Summers, Computers in Physics 6 (1992) 371. <sup>9</sup><sup>9</sup>footnotetext: Sidnie Feit, TCP/IP: Architecture, Protocols, and Implementation, McGraw-Hill, 1993. <sup>10</sup><sup>10</sup>footnotetext: F. Rinaldo and S. Wolbers, Computers in Physics 7 (1993) 184. <sup>11</sup><sup>11</sup>footnotetext: Aleardo Manacero, CPS Performance under Different Network Loads, Fermi Pub-94-33. <sup>12</sup><sup>12</sup>footnotetext: R.J. Harrison, International Journal of Quantum Chemistry 40 (1991) 847; TCGMSG. <sup>13</sup><sup>13</sup>footnotetext: Kirill Denisenko et al., D0 Farm Production System, CHEP ’94, 152; Fermilab $`p\overline{p}.`$ <sup>14</sup><sup>14</sup>footnotetext: E.M. Aitala et al. (E791 Collaboration), Observation of D–$`\pi `$ Production Correlations in 500 GeV $`\pi ^{}`$–N Interactions (submitted to PRL); Search for $`D^0`$$`\overline{D}^0`$ Mixing in Semileptonic Decays, Phys. Rev. Lett. 77 (1996) 2384; Mass Splitting and Production of $`\mathrm{\Sigma }_\text{c}^0`$ and $`\mathrm{\Sigma }_\text{c}^{++}`$ Measured in 500 GeV/c $`\pi ^{}`$–N Interactions, Phys. Lett. B379 (1996) 292; Asymmetries Between the Production of $`D^+`$ and $`D^{}`$ Mesons from 500 GeV/c $`\pi ^{}`$–Nucleus Interactions as a Function of $`x_F`$ and $`p_t^2`$, Phys. Lett. B371 (1996) 157; Search for the Flavor-Changing Neutral Current Decays $`D^+\pi ^+\mu ^+\mu ^{}`$ and $`D^+\pi ^+e^+e^{},`$ Phys. Rev. Lett. 76 (1996) 364.
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# Production of 𝑎₀-mesons in the reactions 𝜋𝑁→𝑎₀𝑁 and 𝑝𝑝→𝑑𝑎₀⁺ at GeV energies Supported by DFG, RFFI and GSI Darmstadt
## 1 Introduction
The scalar mesons play a very important role in the physics of hadrons since they carry the quantum numbers of the vacuum. Nevertheless, the structure of the lightest scalar mesons $`a_0(980)`$ and $`f_0(980)`$ is not yet understood and an important topic of hadronic physics (see e.g. Clo ; Gen ; Jan ; Ani ; Hadron99a ; Hadron99b ; Close2000 and references therein). It has been discussed that they could be either “Unitarized $`q\overline{q}`$ states”, “Four-quark cryptoexotic states”, $`K\overline{K}`$ molecules or vacuum scalars (Gribov’s minions) (see e.g. Ref. Hadron99a ). Nowadays, theory gives some preference to the unitarized quark model proposed by Törnqvist Tornqvist (cf. Hadron99a ; Hadron99b ). However, other options cannot be ruled out so far. Since there is a strong mixing between the uncharged $`a_0(980)`$ and the $`f_0(980)`$ due to a coupling to $`K\overline{K}`$ intermediate states Jan ; Kerbikov , it is important to study independently the uncharged and charged components of the $`a_0(980)`$ because the latter ones do not mix with the $`f_0(980)`$ and preserve their original quark content. It is generally expected, furthermore, that the different $`a_0(980)`$ production cross sections in $`\pi N`$ and $`NN`$ reactions will provide valuable information on its internal structure.
Until now the charged components of the $`a_0(980)`$ have been studied dominantly in the $`\eta \pi ^+`$ or $`\eta \pi ^{}`$ decay channels PDG . Recent experimental data from the E852 Collaboration at BNL give for the charged $`a_0^+`$ meson a mass of $`0.9983\pm 0.0040`$ GeV/c<sup>2</sup> and a width of $`0.072\pm 0.0010`$ GeV/c<sup>2</sup> E852 . Note, that the mass of the $`a_0`$ reported by the E852 Collaboration is significantly larger than the average value of $`0.9834\pm 0.0009`$ GeV/c<sup>2</sup> presented in the last issue of the PDG PDG .
The branching ratios to the two main $`a_0`$ decay channels ($`\eta \pi `$ and $`K\overline{K}`$) are still unclear: the $`\eta \pi `$ mode is quoted by the PDG PDG as ‘dominant’ and the $`K\overline{K}`$ mode as ‘seen’. We point out, that the data from only two experiments Deb ; Abele , where the decay of the $`a_0(980)`$ to $`K\overline{K}`$ was observed, have been used for the PDG analysis PDG . The authors of Ref. Abele report a ratio of branching ratios
$`Br(\overline{p}pa_0\pi ;a_0K\overline{K})/Br(\overline{p}pa_0\pi ;a_0\pi \eta )`$
$`=0.23\pm 0.05.`$ (1)
However, the second branching ratio taken from Ref. Amsler94 might have a systematic uncertainty stemming from a strong interference of the $`a_0`$ signal with the broad resonance $`a_0(1450)`$, which has a width of about 265 MeV. As a consequence the $`a_0(980)`$ maximum in the reaction $`\overline{p}p\eta \pi ^0\pi ^0`$ might be distorted. Moreover, the invariant-mass resolution in Refs. Abele ; Amsler94 is only $`27`$ MeV/c<sup>2</sup>.
In another recent study WA102 the WA102 collaboration reported the branching ratio
$$\mathrm{\Gamma }(a_0K\overline{K})/\mathrm{\Gamma }(a_0\pi \eta )=0.166\pm 0.01\pm 0.02,$$
(2)
which was determined from the measured branching ratio for the $`f_1(1285)`$-meson,
$$\mathrm{\Gamma }(f_1K\overline{K}\pi )/\mathrm{\Gamma }(f_1\pi \pi \eta )=0.166\pm 0.01\pm 0.08,$$
(3)
produced centrally in the reaction $`ppp_f(X_0)p_s`$ at 450 GeV/c. However, the authors assumed that the $`f_1(1285)`$ decays effectively by 100% to $`a_0(980)\pi `$ while the PDG quotes only a branching $`Br(f_1(1285)a_0(980)\pi )=0.34\pm 0.08`$.
Therefore, it is necessary to measure the branching fractions of the two main $`a_0`$ decay channels ($`\eta \pi `$ and $`K\overline{K}`$) under different dynamical conditions with a higher mass resolution ($`\mathrm{\Delta }m<10`$ MeV/c<sup>2</sup>) and lower background in an independent experiment. A related experiment to detect the $`a_0^+`$ in both main decay modes in the reaction $`ppda_0^+`$ will be performed at COSY (Jülich) COSY55 . An important dynamical feature of the latter reaction is that the production of the $`a_0^+(980)`$ near threshold cannot be related to an intermediate production of the $`f_1(1285)`$ (see below).
In this paper we investigate the $`a_0`$\- production cross section in the reactions $`\pi Na_0N`$ and $`ppda_0^+`$ near threshold and at medium energies. In Sect. 2 we present a short overview on the uncertainties of the $`a_0`$\- decay parameters according to present knowledge. To analyze different contributions to the cross section of the reaction $`\pi Na_0N`$ we employ an effective Lagragian approach as well as the Regge-pole model in Sect. 3. The results of this analysis then are used in Sect. 4 to calculate the differential and total cross sections of the reaction $`ppda_0^+`$ within the framework of the two-step model (TSM), in which two nucleons produce an $`a_0`$-meson via $`\pi `$-meson exchange and fuse to a deuteron. The TSM has been used before in Refs. Grishina1 ; Grishina2 for the analysis of $`\eta `$, $`\eta ^{}`$, $`\omega `$ and $`\varphi `$ production in the reaction $`pndM`$ near threshold. An important difference of our analysis here is that the $`S`$-wave channel in the reaction $`ppda_0^+`$ is forbidden due to angular momentum conservation and the Pauli principle and that this reaction is dominated near threshold by the $`P`$-wave contribution. A summary of our work is presented in Sect. 5.
## 2 Models and data on the $`K\overline{K}`$ and $`\pi \eta `$ decay channels of the $`a_0(980)`$
Within the framework of a coupled-channel formalism an appropriate parametrization of the shape of the $`a_0(980)`$ in each ($`\eta \pi `$ or $`K\overline{K}`$) channel can be taken in the form proposed by Flatté Flatte ,
$`|A_i|^2=\mathrm{Const}{\displaystyle \frac{|\mathrm{\Gamma }_i(M)|M_r^2}{(M^2M_r^2)^2+M_r^2|\mathrm{\Gamma }_{tot}^2(M)|}}`$ (4)
where $`M_r`$ is the K-matrix pole, $`\mathrm{\Gamma }_{tot}(M)=\mathrm{\Gamma }_1(M)+\mathrm{\Gamma }_2(M)=g_1\rho _1+g_2\rho _2`$, while $`g_1`$ and $`g_2`$ are coupling constants to the two final states and $`\rho _i`$ is given by the momenta of the final particles $`q_i`$ as $`\rho _i=2q_i/M`$. Note that molecular or ”threshold cusp” cases would imply a dominance of the $`|K\overline{K}`$ component in Fock space and therefore correspond to a relatively large ratio $`R=(g_2/g_1)1`$. In Table 1 we present the most recent results for the $`a_0(980)`$ parameters $`R,M_r`$ and $`g_1`$, which show a sizeable variation especially in the coupling $`g_1`$ and ratio $`R`$, respectively.
In Ref. Bugg it has been shown that, when fitting the $`\eta \pi `$ mass distribution without any additional constraints, the parameters $`M_r`$, $`R`$ and $`g_1`$ cannot be fixed very well. These parameters are strongly correlated and if one of them is moved in steps, the value of $`\chi ^2`$ changes rather slowly, but $`M_r`$, $`R`$ and $`g_1`$ move together. Thus additional constraints are used in most fits. In Ref. E852 a Breit-Wigner (BW) fit of the $`a_0(980)`$ shape in the $`\eta \pi `$ channel has been performed where the mass and width of the $`a_0^+`$ were determined to be 0.9964$`\pm `$0.0016 and 0.062$`\pm `$0.006 GeV/c<sup>2</sup>, respectively. The two extractions of the $`a_0`$ mass and width (BW and Flatté) were found to be statistically consistent. Since in a Breit-Wigner parametrization only two parameters enter, it is not sensitive at all to the ratio $`R`$. This implies that for a reliable determination of $`R`$ the measurements of both channels are necessary. Recall that two zero’s of the function $`D(M)=M^2M_r^2+iM_r(g_1\rho _1(M)+g_2\rho _2(M))`$ define two T-matrix poles on sheets II and III where only the position of the pole in sheet II defines the mass ($`m_0`$) and width ($`\mathrm{\Gamma }_0`$) of the $`a_0(980)`$. Note that the pole mass $`m_0`$ is usually different from the resonance mass $`M_r`$ in Eq. (4). According to the PDG PDG the average value of the $`a_0(980)`$ mass is $`0.9834\pm 0.0009`$ GeV/c<sup>2</sup> for the $`\eta \pi `$ final state (without the new result of the E852 CollaborationE852 ($`0.9983\pm 0.004`$ GeV/c<sup>2</sup>)) and $`0.9808\pm 0.0027`$ GeV/c<sup>2</sup> for the $`K\overline{K}`$ final state Abele . The width of the $`a_0(980)`$ is quoted as $`0.092\pm 0.008`$ GeV in the $`K\overline{K}`$ final state Abele and $`0.072\pm 0.01`$ GeV in the $`\eta \pi `$ final state E852 .
The values of the ratio $`R`$ presented in Table 1 are not in favor of a pure molecular or pure ”threshold cusp” interpretation of the $`a_0(980)`$. This statement is also in line with the results of Ref. Jan , where it was shown that the pure ”threshold cusp” model gives an $`a_0`$ width of about 200 MeV, which is much larger than the experimental value. Nevertheless, there is still a comparatively large uncertainty in $`g_1`$ and $`g_2`$ : the values of $`g_1`$ may vary from 0.12 to 0.32 GeV and $`R=g_2/g_1`$ from 0.9 to 2.05. A better knowledge of $`g_1`$ and $`g_2`$ will help to understand the $`a_0(980)`$ internal structure or its decomposition in Fock space, respectively.
## 3 The reaction $`\pi Na_0N`$
### 3.1 An effective Lagrangian Approach
The most simple mechanisms for $`a_0`$ production in the reaction $`\pi Na_0N`$ near threshold are described by the pole diagrams shown in Fig. 1 a–d. It is known experimentally that the $`a_0`$ couples strongly to the channels $`\pi \eta `$ and $`\pi f_1(1285)`$ because $`\pi \eta `$ is the dominant decay channel of the $`a_0`$ while $`\pi a_0`$ is one of the most important decay channels of the $`f_1(1285)`$ PDG . The amplitudes, which correspond to the $`t`$-channel exchange of $`\eta (550)`$\- and $`f_1(1285)`$\- mesons (a,b), can be written as
$`M_\eta ^t(\pi ^{}pa_0^{}p)=g_{\eta \pi a_0}g_{\eta NN}\overline{u}(p_2^{})\gamma _5u(p_2)`$ (5)
$`\times `$ $`{\displaystyle \frac{1}{tm_\eta ^2}}F_{\eta \pi a_0}(t)F_{\eta NN}(t)`$
$`M_{f_1}^t(\pi ^{}pa_0^{}p)=g_{f_1\pi a_0}g_{f_1NN}`$ (6)
$`\times `$ $`(p_1+p_1^{})_\mu \left(g_{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{m_{f_1}^2}}\right)\overline{u}(p_2^{})\gamma _\nu \gamma _5u(p_2)`$
$`\times `$ $`{\displaystyle \frac{1}{tm_{f_1}^2}}F_{f_1\pi a_0}(t)F_{f_1NN}(t).`$
Here $`p_1`$ and $`p_1^{}`$ are the four momenta of $`\pi ^{},a_0^{}`$, whereas $`p_2`$ and $`p_2^{}`$ are the four momenta of the initial and final protons, respectively; furthermore, $`q=p_2^{}p_2`$, $`t=(p_2^{}p_2)^2`$. The functions $`F_j`$ present form factors at the different vertices $`j`$ ($`j=f_1NN,\eta NN`$), which are taken of the monopole form
$`F_j(t)={\displaystyle \frac{\mathrm{\Lambda }_j^2m_j^2}{\mathrm{\Lambda }_j^2t}},`$ (7)
where $`\mathrm{\Lambda }_j`$ is a cut-off parameter. In the case of $`\eta `$ exchange we use $`g_{\eta NN}=3`$, $`\mathrm{\Lambda }_{\eta NN}`$=1.5 GeV from Ref. Holinde and $`g_{\eta \pi a_0}`$=2.46 GeV which results from the width $`\mathrm{\Gamma }(a_0\eta \pi `$) = 80 MeV. The contribution of the $`f_1`$ exchange is calculated for two parameter sets; set $`A`$: $`g_{f_1NN}=11.2`$, $`\mathrm{\Lambda }_{f_1NN}=1.5`$ GeV from Ref. Bonnf1 , set $`B`$: $`g_{f_1NN}=14.6`$, $`\mathrm{\Lambda }_{f_1NN}=2.0`$ GeV from Ref. Kirchbach and $`g_{f_1a_0\pi }`$=2.5 for both cases. The latter value for $`g_{f_1a_0\pi }`$ corresponds to $`\mathrm{\Gamma }(f_1a_0\pi )=24`$ MeV and $`Br(f_1a_0\pi )=34\%`$.
In Fig. 2 (upper part) we show the differential cross sections $`d\sigma /dt`$ for the reaction $`\pi ^{}pa_0^{}p`$ at 2.4 GeV/c corresponding to $`\eta `$ (dash-dotted) and $`f_1`$ exchanges with set $`A`$ (solid line) and set $`B`$ (dashed line). A soft cut-off parameter (set $`A`$) close to the mass of the $`f_1`$ implies that all the contributions related to $`f_1`$ exchange become negligibly small. On the other hand, for the parameter values given by set $`B`$, the $`f_1`$ exchange contribution is much larger than that from $`\eta `$ exchange. Note, that this large uncertainty in the cut-off presently cannot be controlled by data and we will discuss the relevance of the $`f_1`$ exchange contribution for all reactions separately throughout this study. For set $`B`$ the total cross section for the reaction $`\pi ^{}pa_0^{}p`$ can be about 0.5 mb at 2.4 GeV/c (cf. Fig. 3 (upper part)) while the forward differential cross section can be about 1 mb/GeV<sup>2</sup>.
The $`\eta `$ and $`f_1`$ exchange, however, do not contribute to the amplitude of the charge exchange reaction $`\pi ^{}pa_0^0n`$. In this case we have to consider the contributions of the $`s`$\- and $`u`$-channel diagrams (Fig. 1 c and d):
$`M_N^s(\pi ^{}pa_0^0n)=g_{a_0NN}{\displaystyle \frac{f_{\pi NN}}{m_\pi }}{\displaystyle \frac{1}{sm_N^2}}F_N(s)`$ (8)
$`\times `$ $`p_{1\mu }\overline{u}(p_2^{})\left[(p_1+p_2)_\alpha \gamma _\alpha +m_N\right]\gamma _\mu \gamma _5u(p_2);`$
$`M_N^u(\pi ^{}pa_0^0n)=g_{a_0NN}{\displaystyle \frac{f_{\pi NN}}{m_\pi }}{\displaystyle \frac{1}{um_N^2}}F_N(u)`$ (9)
$`\times `$ $`p_{1\mu }\overline{u}(p_2^{})\gamma _\mu \gamma _5\left[(p_2p_1^{})_\alpha \gamma _\alpha +m_N\right]u(p_2),`$
where $`s=(p_1+p_2)^2,u=(p_2p_1^{})^2`$ and $`m_N`$ is the nucleon mass.
The $`\pi NN`$ coupling constant is taken as $`f_{\pi NN}^2/4\pi =0.08`$ Holinde and the form factor for each virtual nucleon is taken in the form Feuster
$`F_N(u)={\displaystyle \frac{\mathrm{\Lambda }_N^4}{\mathrm{\Lambda }_N^4+(um_N^2)^2}}`$ (10)
with a cut-off parameter $`\mathrm{\Lambda }_N=1.2÷1.3`$ GeV.
The dotted and dash-double-dotted lines in the lower part of Fig. 2 show the differential cross section for the charge exchange reaction $`\pi ^{}pa_0^0n`$ at 2.4 GeV/c corresponding to $`s`$\- and $`u`$\- channel diagrams, respectively. Due to isospin only the $`s`$\- channel contributes to the $`\pi ^{}pa_0^{}p`$ reaction (dotted line in the upper part of Fig. 2). In these calculations the cut-off parameter $`\mathrm{\Lambda }_N`$ = 1.24 GeV and $`g_{a_0NN}^2/4\pi `$=1.075 is taken from Ref. Kirchbach . The solid line in the lower part of Fig. 2 describes the coherent sum of the $`s`$\- and $`u`$\- channel contributions. Except for the very forward region the $`s`$\- channel contribution (dotted line) is rather small compared to the $`u`$\- channel for the charge exchange reaction $`\pi ^{}pa_0^0n`$, which may give a backward differential cross section of about 1 mb/GeV<sup>2</sup> . The corresponding total cross section can be about 0.3 mb at this energy (cf. Fig. 3, middle part).
Unfortunately, there are no experimental data for the total cross section of $`a_0`$ production in $`\pi N`$ collisions near the threshold. Some crude estimates can only be done by comparing the $`a_0`$ production with $`\rho `$ and $`\omega `$ production. For example, the WA57 collaboration has measured inclusive photoproduction of $`a_0^\pm (980)`$ mesons at photon energies of 25 – 55 GeV WA57 . It was found that the cross section of this process is rather large and about $``$ 1/6 of the cross sections for the corresponding non-diffractive production of leading $`\rho ^0,\omega ,\rho ^+`$ and $`\rho ^{}`$ mesons. Furthermore, in the LBL experiment Abolins the measured cross sections $`d\sigma /d\mathrm{\Omega }`$ for the reaction $`ppda_0^+(980)`$ at 3.8 – 6.3 GeV/c are $`(1/4÷1/6)`$ of the cross section for $`\rho ^+`$ production (Table 2).
In view of these arguments we also compare the cross sections for the reactions $`\pi ^{}pa_0^0n`$ and $`\pi ^{}p\rho ^0(\omega )n`$ at 2.4 GeV/c. According to the parametrization of Ref. Sibirtsev we have $`\sigma (\pi ^{}p\rho ^0n)2\sigma (\pi ^{}p\omega n)1.8\mathrm{mb}`$; our estimate then gives $`\sigma (\pi ^{}pa_0^0n)0.15÷0.3`$ mb, which is in a reasonable agreement with the $`u`$\- channel mechanism as well as $`f_1`$ exchange contribution with parameters from set $`B`$ (cf. Fig. 3).
There is a single experimental point for the forward differential cross section of the reaction $`\pi ^{}pa_0^0n`$ at 2.4 GeV/c (Ref. Cheshire , lower part of Fig. 2),
$$\frac{d\sigma }{dt}(\pi ^{}pa_0^0n)|_{t0}=0.49\mathrm{mb}/\mathrm{GeV}^2.$$
Since in the forward region ($`t`$ 0) the $`s`$\- and $`u`$\- channel diagrams only give a smaller cross section, the charge exchange reaction $`\pi ^{}pa_0^0n`$ is most probably dominated at small $`t`$ by the isovector $`b_1(1^+)`$\- and $`\rho _2(2^{})`$\- meson exchanges (see e.g. Ref. Achasov ). Though the couplings of these mesons to $`\pi a_0`$ and $`NN`$ are not known, we can estimate $`\frac{d\sigma }{dt}(\pi ^{}pa_0^0n)`$ in the forward region using the Regge-pole model as developed by Achasov and Shestakov Achasov . Note, that the Regge-pole model is expected to provide a reasonable estimate for the cross section at medium energies of about a few GeV and higher (see e.g. Refs. Kaidalov1 ; Kondrat and references therein).
### 3.2 The Regge-pole model
The $`s`$\- channel helicity amplitudes for the reaction $`\pi ^{}pa_0^0n`$ can be written as
$`M_{\lambda _2^{}\lambda _2}(\pi ^{}pa_0^0n)=\overline{u}_{\lambda _2^{}}(p_2^{})[A(s,t)`$ (11)
$`+`$ $`(p_1+p_1^{})_\alpha \gamma _\alpha {\displaystyle \frac{B(s,t)}{2}}]\gamma _5u_{\lambda _2}(p_2),`$
where the invariant amplitudes $`A(s,t)`$ and $`B(s,t)`$ do not contain kinematical singularities and (at fixed $`t`$ and large $`s`$) are related to the helicity amplitudes as
$`M_{++}sB,M_+M_{++}\sqrt{t_{\mathrm{min}}t}A.`$ (12)
The differential cross section then can be expressed through the helicity amplitudes in the standard way as
$`{\displaystyle \frac{d\sigma }{dt}}(\pi ^{}pa_0^0n)={\displaystyle \frac{1}{64\pi s}}{\displaystyle \frac{1}{(p_1^{\mathrm{cm}})^2}}(|M_{++}|^2+|M_+|^2).`$ (13)
Usually it is assumed that the reaction $`\pi ^{}pa_0^0n`$ at high energies is dominated by the $`b_1`$ Regge-pole exchange. However, as shown by Achasov and Shestakov Achasov this assumption is not compatible with the angular dependence of $`d\sigma /dt(\pi ^{}pa_0^0n)`$ observed at Serpukhov at 40 GeV/c Serpukhov ; Serpukhov1 and Brookhaven at 18 GeV/c Brookhaven . The reason is that the $`b_1`$ Regge trajectory contributes only to the amplitude $`A(s,t)`$ giving a dip in differential cross section at forward angles, while the data show a clear forward peak in $`d\sigma /dt(\pi ^{}pa_0^0n)`$ at both energies. To interpret this phenomenon Achasov and Shestakov introduced a $`\rho _2`$ Regge-pole exchange conspiring with its daughter trajectory. Since the $`\rho _2`$ Regge trajectory contributes to both invariant amplitudes, $`A(s,t)`$ and $`B(s,t)`$, its contribution does not vanish at $`\mathrm{\Theta }=0`$ thus giving a forward peak due to the term $`|M_{++}|^2`$ in $`d\sigma /dt`$. At the same time the contribution of the $`\rho _2`$ daughter trajectory to the amplitude $`A(s,t)`$ is necessary to cancel the kinematical pole at $`t=0`$ introduced by the $`\rho _2`$ main trajectory (conspiracy effect). In this model the $`s`$\- channel helicity amplitudes can be expressed through the $`b_1`$ and the conspiring $`\rho _2`$ Regge trajectories exchange as
$`M_{++}M_{++}^{\rho _2}(s,t)=\gamma _{\rho _2}(t)\mathrm{exp}[i{\displaystyle \frac{\pi }{2}}\alpha _{\rho _2}(t)]\left({\displaystyle \frac{s}{s_0}}\right)^{\alpha _{\rho _2}(t)},`$ (14)
$`M_+M_+^{b_1}(s,t)`$ $`=`$ $`\sqrt{(t_{\mathrm{min}}t)/s_0}\gamma _{b_1}(t)`$ (15)
$`\times `$ $`i\mathrm{exp}[i{\displaystyle \frac{\pi }{2}}\alpha _{b_1}(t)]\left({\displaystyle \frac{s}{s_0}}\right)^{\alpha _{b_1}(t)},`$
where $`\gamma _{\rho _2}(t)=\gamma _{\rho _2}(0)\mathrm{exp}(b_{\rho _2}t)`$, $`\gamma _{b_1}(t)=\gamma _{b_1}(0)\mathrm{exp}(b_{b_1}t)`$, $`t_{\mathrm{min}}m_N^2(m_{a_0}^2m_\pi ^2)/s^2`$, $`s_01`$ GeV<sup>2</sup> while the meson Regge trajectories have the linear form $`\alpha _j(t)=\alpha _j(0)+\alpha _j^{}(0)t`$.
Achasov and Shestakov describe the Brookhaven data on the $`t`$ distribution at 18 GeV/c for $`t_{\mathrm{min}}t0.6`$ GeV<sup>2</sup> Brookhaven by the expression
$`{\displaystyle \frac{dN}{dt}}=C_1\left[e^{\mathrm{\Lambda }_1t}+(t_{\mathrm{min}}t){\displaystyle \frac{C_2}{C_1}}e^{\mathrm{\Lambda }_2t}\right],`$ (16)
where the first and second terms describe the $`\rho _2`$ and $`b_1`$ exchanges, respectively. They found two fits: a) $`\mathrm{\Lambda }_1=4.7`$ GeV$`{}_{}{}^{2},C_2/C_1=0`$; b) $`\mathrm{\Lambda }_1=7.6`$ GeV$`{}_{}{}^{2},C_2/C_12.6`$ GeV$`{}_{}{}^{2},\mathrm{\Lambda }_2=5.8`$ GeV<sup>-2</sup>. This implies that the $`b_1`$ contribution is equal to zero for fit a) and yields only 1/3 of the integrated cross section for fit b) at 18 GeV/c. Moreover, using the available data on the reaction $`\pi ^{}pa_2^0(1320)n`$ at 18 GeV/c and comparing them with the data on the $`\pi ^{}pa_0^0n`$ reaction they estimated the total and forward differential cross sections $`\sigma (\pi ^{}pa_0^0n\pi ^0\eta n)200`$ nb and $`[d\sigma /dt(\pi ^{}pa_0^0n\pi ^0\eta n)]_{t=0}940`$ nb/GeV$`^2.`$ Taking $`Br(a_0^0\pi ^0\eta )0.8`$ we find $`\sigma (\pi ^{}pa_0^0n)0.25`$ $`\mu `$b and $`[d\sigma /dt(\pi ^{}pa_0^0n)]_{t=0}1.2`$ $`\mu `$b/GeV<sup>2</sup>. In this way all the parameters of the Regge model can be fixed and we will employ it for the energy dependence of the $`\pi ^{}pa_0^0n`$ cross section to obtain an estimate at lower energies, too.
The mass of the $`\rho _2(2^{})`$ is expected to be about 1.7 GeV (see Kokoski and references therein) and the slope of the meson Regge trajectory in the case of light ($`u,d`$) quarks is 0.9 GeV<sup>-2</sup> Kaidalov . Therefore, the intercept of the $`\rho _2`$ Regge trajectory is $`\alpha _{\rho _2}(0)=20.9m_{\rho _2}^20.6`$. Similarly – in the case of the $`b_1`$ trajectory – we have $`\alpha _{b_1}(0)0.37`$. At forward angles we can neglect the contribution of the $`b_1`$ exchange (see discussion above) and write the energy dependence of the differential cross section in the form
$`{\displaystyle \frac{d\sigma _{Regge}}{dt}}(\pi ^{}pa_0^0n)|_{t=0}`$ $``$ $`{\displaystyle \frac{d\sigma _{\rho _2}}{dt}}|_{t=0}`$ (17)
$``$ $`{\displaystyle \frac{1}{(p_1^{cm})^2}}\left({\displaystyle \frac{s}{s_0}}\right)^{2.2}.`$
This provides the following estimate for the forward differential cross section at 2.4 GeV/c,
$`{\displaystyle \frac{d\sigma _{Regge}}{dt}}(\pi ^{}pa_0^0n)|_{t=0}0.6\mathrm{mb}/\mathrm{GeV}^2,`$ (18)
which is in agreement with the experimental data point Cheshire (lower part of Fig. 2). Since the $`b_1`$ and $`\rho _2`$ Regge trajectories have isospin 1, their contribution to the cross section for the reaction $`\pi ^{}pa_0^{}p`$ is twice smaller,
$`{\displaystyle \frac{d\sigma _{Regge}}{dt}}(\pi ^{}pa_0^{}p)={\displaystyle \frac{1}{2}}{\displaystyle \frac{d\sigma _{Regge}}{dt}}(\pi ^{}pa_0^0n).`$ (19)
In Fig. 2 the short-dotted lines show the resulting differential cross sections for $`d\sigma _{Regge}(\pi ^{}pa_0^{}p)/dt`$ (upper part) and $`d\sigma _{Regge}(\pi ^{}pa_0^0n)/dt`$ (lower part) at 2.4 GeV/c corresponding to $`\rho _2`$ Regge exchange (fit a) ), whereas the dash-dotted lines indicate the contribution for $`\rho _2`$ and $`b_1`$ Regge trajectories (fit b) ). For $`t0`$ both Regge parametrizations agree, however, at large $`|t|`$ the solution including the $`b_1`$ exchange gives a smaller cross section. The cross section $`d\sigma _{Regge}(\pi ^{}pa_0^{}p)/dt`$ in the forward region exceeds the contributions of $`\eta `$, $`f_1`$ (set $`A`$) and $`s`$\- channel exchanges, however, is a few times smaller than the $`f_1`$ exchange contribution for set $`B`$. On the other hand, the cross section $`d\sigma _{Regge}(\pi ^{}pa_0^0n)/dt`$ is much larger than the $`s`$\- and $`u`$\- channel contributions in the forward region, but much smaller than the $`u`$\- channel contribution in the backward region.
The integrated cross sections for $`\pi ^{}pa_0^{}p`$ (upper part) and $`\pi ^{}pa_0^0n`$ (middle and lower part) for the Regge model are shown in Fig. 3 as a function of the pion lab. momentum by short-dotted lines for $`\rho _2`$ exchange and by short dash-dotted lines for $`\rho _2,b_1`$ trajectories. In the few GeV region the cross sections are comparable with the $`u`$-channel and $`f_1`$ -exchange contribution (set $`B`$). At higher energies it decreases as $`s^{3.2}`$ in contrast to the non-Reggeized $`u`$-channel and $`f_1`$\- exchange contributions which anyhow should only be employed close to the threshold region.
The main conclusions of this Section are as follows: In the region of a few GeV the dominant mechanisms of $`a_0`$ production in the reaction $`\pi Na_0N`$ are $`u`$-channel nucleon and $`t`$-channel $`f_1`$ -meson exchanges which give cross sections for $`a_0`$ production about $`0.3÷0.4`$ mb (cf. upper part of Fig. 3). Similar cross sections ($`0.4÷1`$ mb) are predicted by the Regge model with conspiring $`\rho _2`$ (or $`\rho _2`$ and $`b_1`$) exchanges, normalized to the Brookhaven data at 18 GeV/c (lower part of Fig. 3). The contributions of $`s`$-channel nucleon and $`t`$-channel $`\eta `$ -meson exchanges are small (cf. upper and middle parts of Fig. 3).
## 4 The reaction $`ppda_0^+`$
The missing mass spectrum in the reaction $`ppdX`$ for deuterons produced at $`0^{}`$ in the laboratory and incident momenta of 3.8, 4.5 and 6.3 GeV/c has been measured at LBL (Berkeley) Abolins . It is interesting, that apart from the missing mass peaks corresponding to $`\pi `$ and $`\rho `$ production, there is a distinctive structure in the missing mass spectrum at 0.95 GeV<sup>2</sup>, which was identified as $`a_0`$ production.
In order to estimate the cross section for the reaction $`ppda_0^+`$ at lower momenta (available at COSY) we use the two-step model (TSM) (cf. Refs. Grishina1 ; Grishina2 ). The contributions of hadronic intermediate states to the $`P`$-wave amplitude of the reaction $`ppda_0^+`$ (within the framework of the TSM) are described by the diagrams $`ad`$ in Fig. 4. We consider three different contributions from the amplitude $`\pi Na_0N`$: i) the $`f_1(1285)`$\- meson exchanges (Fig. 4 a); ii) the $`\eta `$\- meson exchange (Fig. 4 b); iii) $`s`$\- and $`u`$\- channel nucleon exchanges (Fig. 4 c and d). As follows from the analysis in Sect. 3 the contributions of the $`\eta `$\- exchange and $`s`$\- channel nucleon can be neglected. We thus restrict to the $`f_1`$\- exchange (set $`B`$) and the $`u`$\- channel nucleon current.
The cut-off $`\mathrm{\Lambda }_N`$ for nucleon exchange (Eq. (10)) is considered as a free parameter now within the interval 1.2 – 1.3 GeV. In order to preserve the correct structure of the amplitude under permutations of the initial nucleons (which are antisymmetric in the isovector state) the amplitude is written as the difference of $`t`$\- and $`u`$\- channel contributions in the form
$$T_{ppdM}^\pi (s,t,u)=A_{ppdM}(s,t)A_{ppdM}(s,u),$$
(20)
where $`M`$ stands for the $`a_0^+`$\- meson. Furthermore, $`s=(p_1+p_2)^2`$, $`t=(p_3p_1)^2`$, $`u=(p_3p_2)^2`$ where $`p_1`$, $`p_2`$, $`p_3`$, and $`p_4`$ are the 4-momenta of the initial protons, meson $`M`$ and the deuteron, respectively. The structure of the amplitude (20) guarantees that the S-wave part vanishes since it is forbidden by angular momentum conservation and the Pauli principle.
Since we are interested in the $`ppda_0^+`$ cross section near threshold, where the momentum of the deuteron is comparatively small, we use a non-relativistic description of this particle by neglecting the 4th component of it’s polarization vector. Correspondingly, the relative motion of nucleons in the deuteron is also treated non-relativistically. Then one can write the first ($`t`$-channel) term on the r.h.s. of Eq. (20) as (Grishina1 )
$`A_{ppda_0^+}(s,t)={\displaystyle \frac{f_{\pi NN}}{m_\pi }}g_{f_1NN}g_{f_1a_0\pi }`$
$`\times `$ $`\sqrt{(p_1^0+m_N)(p_2^0+m_N)}`$
$`\times `$ $`M^{jl}(\stackrel{}{p}_1,\stackrel{}{p}_3)\phi _{\lambda _2}^T(\stackrel{}{p}_2)(i\sigma _2)\sigma ^j\stackrel{}{\sigma }\stackrel{}{ϵ}^{(d)}\sigma ^l\phi _{\lambda _1}(\stackrel{}{p}_1),`$
where $`\stackrel{}{ϵ}^{(d)}`$ is the polarization vector of the deuteron; $`p_1^0=p_2^0=\sqrt{\stackrel{}{p}_1^{\mathrm{\hspace{0.17em}2}}+m_N^2}`$, while $`\phi _\lambda `$ are the (2-component) spinors of the nucleons in the initial state. The tensor function $`M^{jl}(\stackrel{}{p}_1,\stackrel{}{p}_3)`$ is defined by the integral
$`M^{jl}(\stackrel{}{p}_1,\stackrel{}{p}_3)=\sqrt{2m_N}{\displaystyle \frac{d^3p_2^{}}{(2\pi )^{3/2}}}`$
$`\times `$ $`\sqrt{(p_1^0+m_N)(p_2^0+m_N)}\left\{{\displaystyle \frac{p_1^j}{p_1^0+m_N}}+{\displaystyle \frac{p_2^j}{p_2^0+m_N}}\right\}`$
$`\times `$ $`I\mathrm{\Phi }_{\pi ^0Na_0^0N}^l(\stackrel{}{p}_2^{},\stackrel{}{p}_1,\stackrel{}{p}_3){\displaystyle \frac{F_{\pi NN}(q_\pi ^2)}{q_\pi ^2m_\pi ^2}}\mathrm{\Psi }_d(\stackrel{}{p}_2^{}+\stackrel{}{p_3}/2),`$
where the contribution of $`f_1`$\- exchange is given by
$`\mathrm{\Phi }_{\pi ^0Na_0^0N(f_1)}^l(\stackrel{}{p}_2^{},\stackrel{}{p}_1,\stackrel{}{p}_3)=g_{f_1NN}g_{f_1a_0\pi }{\displaystyle \frac{F_{f_1NN}(q_{f_1}^2)}{q_{f_1}^2m_{f_1}^2}}`$ (23)
$`\times `$ $`\{2p_3^l{\displaystyle \frac{2(p_3+p_2^{})^l}{p_1^0+m_N}}(m_N[1+{\displaystyle \frac{m_{a_0}^2q_2^2}{m_{f_1}^2}}]p_3^0)`$
$``$ $`{\displaystyle \frac{2p_1^l}{p_1^0+m_N}}(m_N[1+{\displaystyle \frac{m_{a_0}^2q_2^2}{m_{f_1}^2}}]+p_3^0)\}.`$
The $`u`$\- channel contribution reads
$`\mathrm{\Phi }_{\pi ^0Na_0^0N(u)}^l(\stackrel{}{p}_2^{},\stackrel{}{p}_1,\stackrel{}{p}_3)=g_{a_0NN}{\displaystyle \frac{f_{\pi NN}}{m_\pi }}2m_N`$
$`\times `$ $`\{p_3^l+{\displaystyle \frac{(p_3+p_2^{})^l}{p_1^{\mathrm{\hspace{0.17em}0}}+m_N}}({\displaystyle \frac{m_N}{2}}[3+{\displaystyle \frac{q_N^2}{m_N^2}}]p_3^0)`$
$`+`$ $`{\displaystyle \frac{p_1^l}{p_1^0+m_N}}({\displaystyle \frac{m_N}{2}}[3+{\displaystyle \frac{q_N^2}{m_N^2}}]+p_3^0)\}{\displaystyle \frac{F_N(q_N^2)}{q_N^2m_N^2}}.`$
Here $`\mathrm{\Psi }_d(\stackrel{}{p}_2^{}+\stackrel{}{p_3}/2)`$ is the deuteron wave function for which we use the Paris model Lacombe . In (4) $`I`$ is the isospin factor which depends on the mechanism of the reaction $`pp(pn)a_0^+`$. For $`f_1`$ and $`u`$\- channel exchange we have $`I_{(f_1)}=1`$ and $`I_{(u)}=3\sqrt{2}`$, respectively. Further kinematical quantities, which also dependent on the momenta $`\stackrel{}{p}_1`$, $`\stackrel{}{p}_3`$ and $`\stackrel{}{p}_2^{}`$, are defined as
$`q_\pi ^2=\delta _0(\stackrel{}{p}_2^2+\beta _\pi (\stackrel{}{p}_1))2\stackrel{}{p}_1\stackrel{}{p}_2^{},`$
$`q_{f_1}^2=\delta _0\left(\stackrel{}{p}_2^2+\beta _{f_1}(\stackrel{}{p}_1,\stackrel{}{p}_3)\right)+{\displaystyle \frac{p_3^0}{m_N}}\stackrel{}{p}_2^2`$
$`2\stackrel{}{p}_1\stackrel{}{p}_2^{}2\stackrel{}{p}_3\stackrel{}{p}_2^{}2\stackrel{}{p}_3\stackrel{}{p}_1,`$
$`q_N^2=m_N^2+m_{a_0}^22p_1^0p_3^0+2\stackrel{}{p}_1\stackrel{}{p}_3,`$
$`\beta _\pi (\stackrel{}{p}_1)=(\stackrel{}{p}_1^{\mathrm{\hspace{0.17em}2}}T_1^2)/\delta _0,`$ (25)
$`\beta _{f_1}(\stackrel{}{p}_1,\stackrel{}{p}_3)=\beta _\pi (\stackrel{}{p}_1)m_{a_0}^2/\delta _0+p_3^0m_N`$
$`\delta _0=p_1^0/m_N,T_1=\sqrt{\stackrel{}{p}_1^{\mathrm{\hspace{0.17em}2}}+m_N^2}m_N,`$
$`p_2^{\mathrm{\hspace{0.17em}0}}=\sqrt{\stackrel{}{p}_2^2+m_N^2},p_3^0=\sqrt{\stackrel{}{p}_3^{\mathrm{\hspace{0.17em}2}}+m_{a_0}^2},`$
$`p_1^0=\sqrt{(\stackrel{}{p}_2^{}+\stackrel{}{p}_3)^{\mathrm{\hspace{0.17em}2}}+m_N^2}.`$
with $`m_{a_0}`$ denoting the mass of the $`a_0`$ meson. The form factors $`F_{f_1NN}`$ and $`F_{\pi NN}`$ are taken in the form (7) within $`\mathrm{\Lambda }_{\pi NN}=1.3`$ GeV for the $`\pi NN`$ vertex according to Ref. Holinde and parameter set $`B`$ for the $`f_1NN`$ vertex. The $`u`$\- channel term $`A_{ppda_0^+}(s,u)`$ in Eq. (20) can be obtained from (4) by the substitution $`p_1p_2`$, $`\phi _{\lambda _1}\phi _{\lambda _2}`$.
The differential cross section $`ppda_0^+`$ then can be written as
$`{\displaystyle \frac{d\sigma _{ppda_0^+}}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{64\pi s}}{\displaystyle \frac{1}{(p_1^{\mathrm{cm}})^2}}`$
$`\times `$ $`\overline{|A_{ppda_0^+}(s,t)A_{ppda_0^+}(s,u)|^2}.`$
The calculated forward differential cross section (as a function of the proton-beam momentum) is presented in Fig. 5. The dash-dotted and solid lines describe the results of the TSM for different values of the nucleon cut-off parameter: $`\mathrm{\Lambda }_N=1.2`$ and 1.3 GeV, respectively. A rather good description of the existing data Abolins is achieved for $`\mathrm{\Lambda }_N=1.3`$ GeV (solid line). We recall that in Sect. 3 we have used $`\mathrm{\Lambda }_N=`$ 1.24 GeV from Ref. Kirchbach which gives a cross section in between the dash-dotted and solid line. Our predictions for this cross section practically do not depend on the couplings of the $`f_1`$ since the $`f_1`$ exchange contribution turns out to be very small even for parameter set $`B`$ (dashed line in Fig. 5). The arrow in Fig. 5 indicates the maximum proton momentum presently available at COSY. At this energy a differential cross section of $`0.1÷0.2`$ $`\mu `$b/sr should be expected according to our calculations.
In Fig. 6 the calculated angular differential cross section for the reaction $`ppda_0^+`$ is shown as a function of the center-of-mass angle $`\mathrm{\Theta }`$ which can be parametrized as
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}=A+B\mathrm{cos}^2\mathrm{\Theta }+C\mathrm{cos}^4\mathrm{\Theta }.`$ (27)
The results of our calculations in the framework of the TSM for $`\mathrm{\Lambda }_N=1.3`$ GeV are:
$`A=21.3`$ nb/sr, $`B=15.3`$ nb/sr, $`C=2.1`$ nb/sr
at $`T_{\mathrm{lab}}=2.52`$ GeV ($`\sigma _{tot}=330`$ nb);
$`A=68`$ nb/sr, $`B=76`$ nb/sr, $`C=22`$ nb/sr
at $`T_{\mathrm{lab}}=2.6`$ GeV ($`\sigma _{tot}=1120`$ nb);
$`A=78`$ nb/sr, $`B=97`$ nb/sr, $`C=31`$ nb/sr
at $`T_{\mathrm{lab}}=2.62`$ GeV ($`\sigma _{tot}=1310`$ nb).
We note that an understanding of the $`a_0(980)`$ production mechanism may also give interesting information on its internal structure. For example, the WA57 collaboration has interpreted the relatively strong production of the $`a_0^\pm (980)`$ in photon induced reactions at energies of 25 – 55 GeV as evidence for a $`q\overline{q}`$ state rather than a $`qq\overline{q}\overline{q}`$ state WA57 . This argument can also be used here. If measurements at COSY will confirm a comparatively large value of the $`a_0^+(980)`$\- production cross section as presented in this work, this will provide further evidence that the $`a_0^+(980)`$ has an essential admixture of a $`q\overline{q}`$ component.
## 5 Conclusions
In this work we have estimated $`a_0`$ production cross sections in the reaction $`\pi Na_0N`$ near threshold and at medium energies by considering the $`a_0(980)`$-resonance as a usual $`q\overline{q}`$-meson. Using an effective Lagragian approach we have analyzed different contributions to the differential and total cross sections, i.e. $`t`$\- channel $`\eta `$\- and $`f_1`$\- meson exchanges as well as $`s`$\- and $`u`$-channel nucleon exchanges, and have found that the $`f_1`$\- and $`u`$\- channel contributions are dominant in the $`\pi ^{}pa_0^{}p`$ and $`\pi ^{}pa_0^0n`$ reactions, respectively. We have analyzed also predictions of the Regge model with conspiring $`\rho _2`$ exchange normalized to the data at 18 GeV/c. We found that this model gives (in the few GeV region) a cross section comparable to the $`f_1`$\- and $`u`$\- channel mechanisms.
The latter results have been used to calculate the differential and total cross section of the reaction $`ppda_0^+`$ within the framework of the two-step model, where the amplitude of the $`NNda_0`$ reaction can be expressed through the amplitude of the $`\pi Na_0N`$ reaction and a structure integral containing the deuteron wave function in the non-relativistic limit. It is found that the cross section of the $`ppda_0^+`$ reaction is dominated almost entirely by the $`u`$\- channel mechanism reaching a value of about 1 $`\mu `$b at $`T_{lab}`$ = 2.6 GeV. An experimental confirmation of this comparatively large production cross section would imply that the $`a_0^+(980)`$ has an essential admixture of a $`q\overline{q}`$ component.
## Acknowledgements
We are grateful to A. Sibirtsev for helpful discussions.
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# Untitled Document
Arakelov-type inequalities for Hodge bundles
Chris Peters Department of Mathematics University of Grenoble I Saint-Martin d’Hères, France
June 20 2000
Prépublication de l’Institut Fourier n<sup>o</sup> 511 (2000) http://www-fourier.ujf-grenoble.fr/prepublications.html Keywords: Variations of Hodge structure, Higgs bundles, Higgs field, period map, Hodge bundles.AMS Classification: 14D07, 32G20
§ 0. Introduction
The inequalities from the title refer back to Arakelov’s article \[Arakelov\]. The main result of that paper is:
Theorem. Fix a complete curve $`C`$ of genus $`>1`$ and a finite set $`S`$ of points on $`C`$. There are at most finitely many non-isotrivial families of curves of given genus over $`C`$ that are smooth over $`CS`$.
The proof consists of two parts. First one proves that there are only finitely many such families (this is a boundedness statement) by bounding the degree $`d`$ of the relative canonical bundle in terms of the genus $`p`$ of $`C`$, the genus $`g`$ of the fiber and the cardinality of the set $`S`$ :
$$0d(2p2+\mathrm{\#}S)\frac{g}{2}.$$
The second part consists of establishing rigidity for a non-isotrivial family. It follows upon identifying the deformation space of the family with the $`H^1`$ of the inverse of the relative canonical bundle, which is shown to be ample. Kodaira vanishing then completes the proof.
This approach can be carried out for other situations as well. In fact \[Faltings\] deals with the case of abelian varieties and shows that boundedness always holds and that for rigidity one has to impose further conditions besides non-isotriviality. Subsequently the rigidity statement has been generalized in \[Peters90\] and using his result, the case of K3-surfaces, resp. Abelian varieties could be treated completely by Saito and Zucker in \[Saito-Zucker\], resp. by Saito in \[Saito\].
The boundedness statement is an inequality for the degree $`d`$ of the direct image of the relative canonical bundle, i.e., the (canonical extension) of the Hodge bundle and one can ask for bounds for the degrees of the other Hodge bundles. In fact, the main result of this note gives such a bound for the Hodge components of complex variations of Hodge structures in terms of ranks of iterates of the Higgs field (the linear map between Hodge components induced by the Gauss-Manin connection). By way of an example, we have :
Proposition. Let $`V=_{p=0}^wV^{p,wp}`$ be a real weight $`w`$ variation of Hodge structures over a punctured curve $`CS`$ with unipotent local monodromy-operators. Let $`\sigma ^p:V^{p,wp}V^{p1,wp+1}`$ the $`k`$-th component of the Higgs field and $`\sigma ^k=\sigma _{wk+1}\mathrm{}\sigma _w:V^{w,0}V^{wk,k}`$ the $`k`$-th iterate. Then
$$0degV^{w,0}(2p2+\mathrm{\#}S)\left(\underset{r=1}{\overset{w}{}}\frac{r}{2}\left(rank\sigma ^{r1}rank\sigma ^r\right)\right).$$
Some time ago I sketched a proof of a similar, but weaker inequality in \[Peters86\], but the details of this proof never appeared. Deligne found an amplification of my argument when the base is a compact curve leading to optimal bounds for complex variations of Hodge structures (letter to the author 18/2/1986). The principal goal of this note is to give a complete proof of the refined inequalities (for a complex variation with quasi-unipotent local monodromy-operators over a quasi-projective smooth curve) based on this letter in the light of later developments which I sketch below.
The reason for writing up this note stems from a recent revival of interest in this circle of ideas: on the one hand Parshin posed me some questions related to this. On the other hand, Jost and Zuo sent me a preprint \[JostZuo\] containing similar (but weaker) bounds obtained by essentially the same method.
Continuing with the historical development, Deligne’s letter and subsequent correspondence between Deligne and Beilinson together with basic ideas and results of Hitchin paved the way for the theory of Higgs bundles, the proper framework for such questions. See \[Simpson92\], \[Simpson94\] and \[Simpson95\] for a further discussion of these matters.
In relation with Simpson’s work, I should remark that the boundedness result from \[Simpson94\] (Corollary 3.4 with $`\mu =0`$ and $`P=0`$) immediately implies that the Chern numbers of Hodge bundles underlying a complex variation of given type over a compact projective variety can only assume finitely many values. It follows that there are a priori bounds on these Chern numbers. If the base is a (not necessarily compact) curve, Simpson shows in \[Simpson90\] that the same methods give boundedness for the canonical extensions of the Hodge bundles. Explicit bounds were not given however. Conversely, knowing that bounds on the degrees of the Hodge bundles exist can be used to simplify the rather technical proof for boundedness, as outlined in §4.
As to the further contents of this note, in §1 I rephrase known bounds on the curvature of Hodge bundles (with respect to the Hodge metric) in a way that shows how to adapt these in the case of non-compact algebraic base manifolds. The main argument is given in §2 while in §3 I explain what one has to change in the non-compact case.
In closing, I want to mention Eyssidieux work which treats a generalization of the true Arakelov inequality (weight one Hodge structures on a compact curve) to higher dimensional base manifolds. The inequalities concern variations of Hodge structure over a compact Kähler manifold $`M`$ such that the period map is generically finite onto its image. See \[Eyss\].
§ 1. Curvature bounds for Hodge bundles
Let me recall that a complex variation of Hodge bundles of weight $`w`$ on a complex manifold $`M`$ consists of a complex local system $`V`$ on $`M`$ with a direct sum decomposition into complex subbundles
$$V=\underset{p+q=w}{}V^{p,q}$$
with the property that the canonical flat connection $``$ satisfies the transversality condition
$$:V^{p,q}𝒜^{1,0}(V^{p1,q+1})𝒜^1(V^{p,q})𝒜^{0,1}(V^{p+1.q1}).$$
$`()`$
The local system $`V`$ defines a holomorphic vector bundle denoted by the same symbol. The bundles $`V^{p,q}`$ are not necessarily holomorphic, but the transversality condition implies that the filtration defined by
$$F^pV:=\underset{rp}{}V^{r,wr}$$
is holomorphic. We have a $`C^{\mathrm{}}`$-isomorphism between the Hodge bundle
$$V_{\mathrm{Hdg}}:=\underset{p}{}F^{p+1}V/F^pV$$
and the bundle $`V`$ which we shall use to transport the connection to the former.
Of course, the usual (real) variations are examples of complex variations. These satisfy the additional reality constraint $`V^{p,q}=\overline{V}^{q,p}`$. Conversely, a complex variation $`V`$ together with its complex conjugate, defines a real variation of Hodge structures on $`V\overline{V}`$ in the obvious manner. In passing, we observe however that we can have complex variations of any given pure type. If for instance $`V^{p,q}`$ has the property that it is preserved by $``$ i.e., if it is a flat subbundle, it is itself a complex variation.
One says that $`V`$ is polarized by a bilinear form $`b`$ if $`b`$ is preserved by $``$ and the two Riemann bilinear equations are verified:
$$\begin{array}{cc}\hfill b(u,v)& =0uH^{p,q},vH^{r,s},(u,v)(s,r)\hfill \\ \hfill h_C(u,v)& =(\mathrm{i})^wb(Cu,\overline{v})\text{is a positive definite hermitian metric.}\hfill \end{array}$$
Here $`C`$ is the Weil-operator which equals multiplication with $`\mathrm{i}^{pq}`$ on $`H^{p,q}`$. Instead of $`b`$ one can also consider the hermitian form
$$h(u,v)=\underset{p}{}(1)^ph_C(u,v)=\mathrm{i}^wb(u,\overline{w}),$$
preserved by $``$. As before, we shall transport $`b`$, $`h`$ and $`h_C`$ to $`V_{\mathrm{Hdg}}`$ using the $`C^{\mathrm{}}`$-isomorphism $`V\begin{array}{c}\hfill \end{array}V_{\mathrm{Hdg}}`$.
If we decompose
$$=\sigma ^{}+D+\sigma ^+$$
according to the transversality condition (\*), the operator $`D`$ is a connection on $`V_{\mathrm{Hdg}}`$ and the operators $`\sigma ^+`$ and $`\sigma ^{}`$ are $`𝒜_M^0`$-linear. Moreover, the operator $`\sigma ^+`$ is the $`h_C`$-conjugate of $`\sigma ^{}`$ so that one may write
$$\begin{array}{cc}\hfill \sigma :=& \sigma ^{}\hfill \\ \hfill \sigma ^{}=& \sigma ^+\hfill \end{array}$$
The $`\overline{}`$-operator for the holomorphic structure on $`V_{\mathrm{Hdg}}`$ is given by the $`(0,1)`$-part $`D^{0,1}`$ of $`D`$ and $`D`$ preserves the metric $`h_C`$. So $`D`$ is the Chern connection, i.e., the unique metric connection on the Hodge bundle whose $`(0,1)`$-part is the $`\overline{}`$-operator.
Decomposing the equation $`=0`$ into types yields various equalities. The first
$$0=\overline{}(\sigma ):=\overline{}\sigma +\sigma \overline{}$$
says that $`\sigma `$ is a holomorphic endomorphism. The second
$$\sigma \sigma =0$$
implies that the pair $`(V_{\mathrm{Hdg}},\sigma )`$ is a so-called Higgs bundle. By definition this is a pair $`(E,\theta )`$ consisting of a holomorphic bundle $`E`$ with a holomorphic map $`\theta :E\mathrm{\Omega }^1(E)`$ satisfying $`\theta \theta =0`$ in $`\mathrm{\Omega }^2End(E)`$. The endomorphism $`\theta `$ is also called the Higgs field. Recall also that a hermitian metric $`k`$ on a Higgs bundle is called harmonic if its Chern connection $`D_k`$ combines with $`\theta `$ and its $`k`$-conjugate $`\theta ^{}`$ to give a flat connection $`\theta +D_k+\theta ^{}`$. In our case, the Hodge metric $`h_C`$ is indeed harmonic.
Summarizing the preceding discussion, from a polarized complex variation of Hodge bundles over $`M`$ we have constructed a Higgs bundle equipped with a harmonic metric. If $`M`$ is compact, such bundles can be shown to be semi-stable. Let me recall that a harmonic Higgs bundle $`(V,h)`$ is stable resp. semi-stable if for any proper Higgs subsheaf $`WV`$ i.e., a coherent submodule preserved by the Higgs field, one has an inequality of slopes
$$\mu (W)<\mu (V),\text{resp.}\mu (W)\mu (V).$$
The slope for vector bundle $`E`$ on a Kähler manifold $`(M,\omega )`$ is
$$\mu (E)=deg(E)/rank(E),$$
with the degree of $`E`$ is defined using the Kähler metric:
$$deg(E)=c_1(E)[\omega ]^{m1},m=dimM.$$
Note that $`deg(E)=deg(^rE)`$, $`r=rankE`$. For a torsion free coherent sheaf $`E`$, one has to replace $`E`$ by the line bundle which is the double dual of $`^rE`$.
Semi-stability of harmonic Higgs bundles can be proved as a consequence of the curvature formula, which is a straightforward consequence of the transversality relation:
1.1. Lemma. The curvature of the Chern connection $`D`$ on $`V_{\mathrm{Hodge}}`$ is given by:
$$F_D=[\sigma ,\sigma ^{}].$$
Below (see Lemma 1.4) I give a proof of a refined version of the semi-stability property. Let me observe however that semi-stability is part of the complete characterization of harmonic Higgs bundles as found by Simpson (see \[Simpson92\]):
Theorem. A Higgs bundle over a compact Kähler manifold $`M`$ with a harmonic metric is the direct sum of stable Higgs bundles with the same slope. Any local system which is the direct sum of irreducible local systems admits the structure of a Higgs bundle with a harmonic metric. The category of Higgs bundles admitting a harmonic metric is equivalent to the category of semi-simple local systems, i.e., those that are direct sums of irreducible local systems.
Observe that a Higgs bundle has zero Chern classes since it carries a flat connection and hence semi-stability means that any Higgs subbundle has non-positive first Chern class. I need a refinement of this in terms of the first Chern forms
$$\gamma _1(E,h)=\frac{\mathrm{i}}{2\pi }\text{(Trace of the curvature of the Chern connection)}.$$
So I need to estimate the curvature (with respect to the Hodge metric) of a graded Higgs subbundle $`W=_pW^{p,wp}`$, i.e.$`W^{p,wp}V^{p,wp}`$ and $`\sigma W^{p,wp}\mathrm{\Omega }_M^1(W^{p1,w+p+1})`$. The curvature estimate we are after reads:
1.2. Lemma. For a subsystem $`WV`$ we have
$$\mathrm{i}Tr[\sigma |W_{\mathrm{Hdg}},(\sigma |W_{\mathrm{Hdg}})^{}]0$$
with equality everywhere if and only if the orthogonal complement $`W_{\mathrm{Hdg}}^{}`$ with respect to $`h_C`$ is preserved by $`\sigma `$ so that we have a direct sum decomposition of Higgs bundles
$$W_{\mathrm{Hdg}}W_{\mathrm{Hdg}}^{},$$
or, equivalently, of complex systems of Hodge bundles. Proof: If we split $`V=WW^{}`$ into a $`C^{\mathrm{}}`$-orthogonal sum with respect to the Hodge metric and write $`\sigma `$ into block-form, the fact that $`\sigma `$ preserves $`W`$ means that this block-form takes the shape
$$\sigma =\left(\begin{array}{cc}S& T\\ 0& S^{}\end{array}\right).$$
Then $`Tr[\sigma |W,(\sigma |W)^{}]=Tr(SS^{}SS^{}+TT^{})=TrTT^{}`$. If we multiply this with $`\mathrm{i}`$ this is a positive definite $`(1,1)`$-form. The trace is zero if and only if $`T=0`$ which means that $`\sigma `$ preserves also $`W^{}`$.
The last clause is a consequence of the following discussion. Q.E.D.
Let me compare the flat structure $`(V,)`$ and the holomorphic structure $`(V_{\mathrm{Hdg}},\sigma )`$ using the principle of pluri-subharmonicity : on a compact complex manifold there are only constant pluri-subharmonic functions. The result is:
1.3. Lemma. Suppose $`M`$ is a smooth projective variety. A holomorphic section of $`V_{\mathrm{Hdg}}`$ satisfies $`\sigma (s)=0`$ if and only if $`(s)=0`$. In other words the flat sections are precisely the holomorphic sections of the Hodge bundle killed by the Higgs field. Proof: The well known Bochner-type formula (see \[Schmid, §7\])
$$\begin{array}{cc}\hfill \overline{}h_C(s,s)=& h_C(D_hs,D_hs)h_C(F_hs,s)h_C([\sigma (s),\sigma ^{}(s)],s)\hfill \\ \hfill =& h_C(\sigma (s),\sigma (s))+h_C(\sigma ^{}(s),\sigma ^{}(s))\hfill \end{array}$$
shows that when $`\sigma (s)=0`$, the function $`h_C(s,s)`$ is pluri-subharmonic and so constant. Hence $`D_hs=0=\sigma ^{}(s)`$, implying $`(s)=0`$. Conversely, a flat holomorphic section satisfies $`\sigma (s)=0`$ by type considerations. Q.E.D.
1.4. Corollary. A complex subbundle $`WV`$ is a subsystem, i.e.$`WW\mathrm{\Omega }_X^1`$ if and only if $`W_{\mathrm{Hdg}}V_{\mathrm{Hdg}}`$ is a Higgs subbundle.
Noticing the minus sign in the curvature of the Hodge metric and recalling that the curvature decreases on subbundles with equality if and only if the quotient bundle is holomorphic, Lemma 1.2 implies:
1.5. Corollary. The first Chern form of a graded Higgs subbundle $`W`$ of $`V`$ (with respect to the Hodge metric) is negative semi-definite and it is zero everywhere if and only if $`W^{}`$ is a (holomorphic) graded Higgs subbundle as well. In this case the complex variation splits as complex polarized variations of Hodge structures $`V=WW^{}`$.
1.6. Remark. These point-wise estimates can be integrated over any compact Kähler variety $`M`$ showing that the slope of a Higgs subbundle is non-positive thereby proving semi-stability.
One also needs to know what happens for a morphism $`f:V_1V_2`$ of graded Higgs bundles. In general neither the kernel nor the image are bundles, although they are preserved by the Higgs fields. Since the kernel of $`f`$ is torsion free, its degree is well-defined. For the image this is not necessarily the case. We have to replace it by its saturation. Recall that a subsheaf $`FV`$ of a locally free (or torsion free) sheaf $`V`$ is saturated if $`V/F`$ is torsion free. Any subsheaf $`FV`$ is a subsheaf of a unique saturated subsheaf $`F^{\mathrm{sat}}`$, its saturation which is defined as the inverse image under $`VV/F`$ of the torsion of the target. For Higgs bundles on curves the natural morphisms
$$\sigma _p:V^{p,wp}V^{p1,w+1p}\mathrm{\Omega }_M^1$$
then can be used to define image bundles
$$Im(\sigma _p)=\left(\sigma _pV^{p,wp}\right)^{\mathrm{sat}}(\mathrm{\Omega }_M^1)^{}.$$
§ 2. A bound for variations over a compact curve
In this section I investigate polarized complex systems of Hodge bundles over a compact curve. The bundles $`V^{p,q}`$ can be given holomorphic structures via the $`C^{\mathrm{}}`$-isomorphism $`V^{p,q}\begin{array}{c}\hfill \end{array}F^p/F^{p+1}V_{\mathrm{Hdg}}`$ and in the sequel we tacitly make this identification.
2.1. Theorem. Let $`V=_pV^{p,wp}`$ be a weight $`w`$ polarized complex system of Hodge over a compact curve $`C`$ of genus $`p`$. Put
$$\chi (C)=deg(\mathrm{\Omega }_C^1)=22p.$$
Let $`\sigma =\sigma _p`$ be the associated Higgs field and let $`\sigma ^k:V^{p,wp}V^{pk,wp+k}`$ the composition $`\sigma _{pk+1}\mathrm{}\sigma _p`$. Then we have the inequality
$$degV^{p,wp}\chi (C)\underset{r1}{}\frac{r}{2}\left(rank\sigma ^{r1}rank\sigma ^r\right)$$
Equality holds if and only if for some $`k0`$ the maps $`\sigma _p,\mathrm{},\sigma _{pk}`$ are all isomorphisms and $`\sigma _{pk1}`$ is zero. In this case $`V^{p,wp}\mathrm{}V^{pk1,wp+k+1}`$ is a complex subvariation.
Proof (Deligne): Consider the following graded Higgs subbundles $`W_rV`$, $`r=1,2,\mathrm{}`$ which only are non-zero in degrees $`p,\mathrm{},pr+1`$ and which — using the notation of the image under $`\sigma `$ as defined in the previous section — are given by
$$\begin{array}{cc}\hfill W_r^{p,wp}:=Ker\sigma ^rV^{p,wp},& W_r^{p1,wp+1}:=Im(\sigma Ker\sigma ^r)V^{p1,wp+1},\mathrm{}\hfill \\ \hfill \mathrm{}& W_r^{pr+1,wp+r1}:=Im(\sigma ^{r1}Ker\sigma ^r)V^{pr+1,wp+r1}.\hfill \end{array}$$
Apply the stability property to these bundles. For simplicity we put
$$\begin{array}{cc}\hfill d_r& =deg(Ker\sigma ^r/Ker\sigma ^{r1})\hfill \\ \hfill l_r& =dim(Ker\sigma ^r/Ker\sigma ^{r1})=rank\sigma ^{r1}rank\sigma ^r\hfill \end{array}$$
Now compute the degree of each of these $`r`$ bundles $`W_r^{pk,wp+k}`$ using the isomorphism
$$Ker\sigma ^r/Ker\sigma ^k\begin{array}{c}\hfill \end{array}\sigma ^k(Ker\sigma ^r)^{\mathrm{sat}}V^{pk,wp+k}\left(\mathrm{\Omega }_C^1\right)^k.$$
The bundle on the left has rank $`l_r+\mathrm{}+l_{k+1}`$ and degree $`d_r+\mathrm{}+d_{k+1}`$. It follows that $`degW_r^{pk,wp+k}=deg\left(Im(\sigma ^kKer\sigma ^r)\right)=deg\left((Ker\sigma ^r/Ker\sigma ^k)\left(\mathrm{\Omega }_C^1\right)^k\right)=d_r+\mathrm{}+d_{k+1}+(l_r+\mathrm{}+l_{k+1})k\chi `$ and adding these for $`k=0,\mathrm{},r1`$ one finds the total degree
$$degW_r=\left(\underset{p=1}{\overset{r}{}}pd_p+l_p\frac{p(p1)}{2}\chi (C)\right)0.$$
Now take a weighted sum of these inequalities in order to make appear $`degV^p=d_p`$. Indeed
$$0\underset{k=1}{\overset{\mathrm{}}{}}\left(\frac{1}{k}\frac{1}{k+1}\right)\left(\underset{p=1}{\overset{k}{}}pd_p+l_p\frac{p(p1)}{2}\chi (C)\right)=\underset{p1}{}d_p+\underset{p1}{}l_p\frac{p1}{2}\chi (C)$$
You have equality if and only if all of the equalities you started out with are equalities which means that the $`W_r`$ are direct factors as graded Higgs subbundle of $`V`$. This translates into the $`\sigma _r`$ being either zero or an isomorphism. To show this, consider the bundle $`W_1`$, i.e.$`Ker\sigma _pV^{p,wp}`$. If this is direct factor of $`V`$ as a graded Higgs bundle, the orthogonal complement has the structure of a graded Higgs bundle. Since the latter, if not zero, must have type different from $`(p,wp)`$ it follows that $`\sigma _p`$ is either $`0`$ or has maximal rank. If it is the zero map, the kernel $`V^{p,wp}`$ is itself a subvariation. If it has maximal rank, its kernel is zero and hence has degree $`0`$ and we repeat the argument with $`\sigma _{p1}`$. If this map is zero $`V^{p,wp}V^{p1,wp+1}`$ is a subvariation and if $`\sigma _{p1}`$ is an isomorphism we continue the procedure. In this way we produce a chain $`\sigma _p,\mathrm{},\sigma _{pk}`$ of isomorphisms such that the next map $`\sigma _{pk1}`$ is zero as asserted. Q.E.D.
Next, we remark that these bounds applied to the dual variation of Hodge structure gives bounds in the other direction. We shall make these explicit for real variations, which are self-dual in the obvious sense. In fact, the reality constraint $`H^{p,q}=\overline{H}^{q,p}`$ implies $`degH^{p,q}=degH^{q,p}`$. In the real case, again one has equality only for a consecutive string $`\sigma _p,\mathrm{},\sigma _{pk}`$ of isomorphisms such that the next one is zero. Note however that the complex subvariation splitting the real variation itself need not be real. The reality constraint itself imposes however some additional constraints.
2.2. Corollary. Let $`V`$ be a variation of polarized real weight $`w`$ Hodge structures over a compact curve $`C`$ with Higgs field $`\sigma =\sigma _p`$. Put $`\overline{\sigma }_p=\sigma _{wp}`$ and $`\sigma ^k=\sigma _{pk1}\mathrm{}\sigma _p:V^{p,wp}V^{pk,wp+k}`$ and $`\overline{\sigma }^k=\overline{\sigma }_{pk1}\mathrm{}\overline{\sigma }_p:V^{wp,p}V^{wp+k,pk}`$. Then we have the bounds
$$\begin{array}{cc}\hfill \chi (C)\underset{r1}{}\frac{r}{2}\left(rank\overline{\sigma }^{r1}rank\overline{\sigma }^r\right)& degV^{p,wp}\hfill \\ & \chi (C)\underset{r1}{}\frac{r}{2}\left(rank\sigma ^{r1}rank\sigma ^r\right).\hfill \end{array}$$
In particular, if $`V^{p,wp}0`$ precisely in the interval $`p=0,\mathrm{},w`$, the Hodge bundle $`V^{w,0}`$ satisfies the bounds
$$0degV^{w,0}\chi (C)\underset{r=1}{\overset{w}{}}\frac{r}{2}\left(rank\sigma ^{r1}rank\sigma ^r\right)$$
and equality on the left holds if and only if $`V^{w,0}`$ is a flat subbundle, while equality on the right holds if and only if all the maps $`\sigma _k`$ are isomorphisms.
Let me translate this Corollary in terms of period maps. The period domains involved are the Griffiths domain $`D`$ resp. $`D^{}`$ parametrizing real polarized Hodge structures of weight $`w`$ of fixed given type, resp. the partial Hodge flags $`V^{w,0}V`$ of the same type. Giving a polarized variation of Hodge structures of this type over the curve $`C`$ is the same as giving its period map
$$p:CD/\mathrm{\Gamma },$$
where $`\mathrm{\Gamma }`$ is the monodromy group. Let me refer to \[Griffiths\] for the necessary background. Likewise, the partial Hodge flag is given by the partial period map
$$p^{}:CD^{}/\mathrm{\Gamma }^{}$$
defined as the composition of $`p`$ and the forgetful map $`q:D/\mathrm{\Gamma }D^{}/\mathrm{\Gamma }^{}`$. The subbundle $`V^{w,0}=F^w`$ being flat means exactly that $`p^{}`$ is constant. At the other end of the spectrum, using that the maps $`\sigma _k`$ measure the derivative of the period map, the latter is everywhere an immersion if and only if all $`\sigma _k`$ are isomorphisms. Summarizing:
2.3. Corollary.
i) One has $`degW^{w,0}=0`$ if and only if the image of the full period map lands in one of the fibers of $`q`$.
ii) The equality $`degV^{w,0}=\chi (C)_{r=1}^w\frac{r}{2}\left(rank\sigma ^{r1}rank\sigma ^r\right)`$ holds if and only if the perod map is everywhere an immersion.
2.4. Remark. Put
$$\begin{array}{cc}\hfill v^{p,q}:=& dimV^{p,q}\hfill \\ \hfill v_0^{p,q}:=& dimKer(\sigma _p:V^{p,q}V^{p1,q+1}).\hfill \end{array}$$
Using the estimate
$$rank\sigma _{wr}rank\sigma ^{r+1}$$
and a telescoping argument, the above inequality implies the (a priori weaker) bound
$$degE^{w,0}\chi (C)\underset{p}{}(v^{p,wp}v_o^{p,wp})$$
and this is due to \[JostZuo\] (they combine this with the obvious symmetry among the Hodge numbers and, similarly, the numbers $`v_0^{p,q}`$).
§ 3. The case of a smooth non-compact curve
We let $`C`$ be a smooth compactification and we let $`S=CC_0`$. Assume now that $`V_0`$ is a complex variation of Hodge structure with the property that the local monodromy operators $`\gamma _s`$ around the points $`sS`$ are quasi-unipotent. This is true if the complex variation is defined over the integers, e.g. if the variation comes from geometry. This is the monodromy theorem, whose proof can be found in \[Schmid\].
The assumption imply that $`V_0`$ admits a vector bundle extension $`V`$ over $`C`$ and all the Hodge bundles $`F^pV_0`$ extend to holomorphic bundles $`F^pVV`$ on $`C`$. The Higgs fields are now maps
$$\sigma :VV\mathrm{\Omega }_C^1(\mathrm{log}S).$$
Asymptotic analysis of the Hodge metric $`h_p`$ on every Hodge bundle $`F^p`$ around the punctures shows that its Chern form is integrable and that there are non-negative rational numbers $`\alpha _s^p`$, the residues of the operators $`\gamma _p`$ such that
$$deg(F^p)=_C\gamma _1(F^p,h_p)+\underset{sS}{}\alpha _s^p.$$
For details of this I refer to \[Peters84\]. The residues are always zero in the unipotent case and conversely, if the Hodge bundles are flat the vanishing of all the residues implies that the local monodromy operators are all unipotent. In general, the residue is given by
$$\alpha _s^p=\underset{\alpha [0,1)}{}\alpha dim(F^pV_\alpha )$$
$`()`$
where $`V_\alpha `$ is the subspace of $`V`$ on which the local monodromy operator $`T_s`$ acts with eigenvalue $`exp(2\pi i\alpha )`$.
The estimates of the previous section go through for the first summand and also for its successive quotients, the Hodge bundles. So we find
3.1. Theorem. Let $`V`$ be a complex variation over a curve $`C_0=CS`$ with quasi-unipotent local monodromy operators. Put $`genus(C)=g`$. Let $`\sigma =\sigma _p`$ be the associated Higgs field and let $`\sigma ^k:V^{p,wp}V^{pk,wp+k}`$ the composition $`\sigma _{pk+1}\mathrm{}\sigma _p`$. Then the degree of the Hodge bundles are bounded by
$$degV^{p,wp}(2g2+\mathrm{\#}S)\underset{r}{}\frac{r}{2}\left(rank\sigma ^{r1}rank\sigma ^r\right)+\underset{sS}{}(\alpha _s^p\alpha _s^{p+1}),$$
with $`\alpha _s`$ given by (\*).
If the local monodromy-operators are unipotent, the last summand vanishes and one has a description for those variations where the bound is attained:
3.2. Addition. In the situation of the preceding theorem, assume that the local monodromy operators are unipotent. Equality is attained if and only if for some $`k0`$ the maps $`\sigma _p,\mathrm{},\sigma _{pk}`$ are all isomorphisms and $`\sigma _{pk1}`$ is zero. In this case $`V_0^{p,wp}\mathrm{}V_0^{pk1,wp+k+1}`$ is a complex subvariation of $`V_0`$. The proof is the same as before; it depends on the validity of Lemma 1.3 in the case of non-compact curves. The principle of pluri-subharmonicity needed here can be found in \[Schmid\], proof of Theorem (7.22). It really depends on his $`SL(2)`$-orbit theorem and so it is considerably less trivial.
3.3. Remark. For real variations there are bounds as in Corollary 2.2:
$$0\underset{sS}{}\alpha _s^pdegV^{w,0}(2p2+\mathrm{\#}S)\underset{r=1}{\overset{w}{}}\frac{r}{2}\left(rank\sigma ^{r1}rank\sigma ^r\right)+\underset{sS}{}\alpha _s^p$$
and if for instance $`degV^{w,0}=0`$, all the local monodromy operators are unipotent and $`V^{w,0}`$ is a flat subbundle of $`V`$. See also \[Peters84\].
§ 4. Boundedness revisited
Here we give the proof that a bound on the degrees of all of the Hodge bundles imply boundedness. To be precise, we need both upper and lowerbounds so that only finitely many degrees are possible. We repeat here that lower bounds are obtained from the upper bounds of the dual variation.
First we need to see that only those variations that occur on a fixed local systems need to be considered. For this we invoke Deligne’s boundedness result from \[Deligne\]:
Theorem. Given a smooth algebraic variety $`S`$ there are only finitely many classes of representations of $`\pi _1(S)`$ on a given rational vector space $`V`$ such that the resulting flat vector bundle underlies a polarizable complex variation of Hodge structure of a given weight.
I would like to point out that the proof of this theorem is not too difficult, except for one point: it uses the (by now standard fact) that a polarizable complex variation of Hodge structure is semi-simple; in any case it is much easier than Simpson’s methods used to show boundedness.
We complete the argument by showing:
4.1. Lemma. Let there be given a flat vector bundle on a curve $`C`$ and a finite set of integers $`\{d_p\}`$, $`pI`$. Complex variations $`V=_{pI}V^{p,wp}`$ with $`degV^{p,wp}=d_p`$ are parametrized by an open set of an algebraic variety. Proof: Consider the bundle $``$ over $`C`$ whose fiber over $`c`$ consists of all flags in $`V_c`$, the fiber of the given complex variation $`V`$ at $`c`$, which are of the same type as the Hodge flags. Then the total space $``$ is a projective variety. There is a subvariety $`𝒢`$ of $``$ formed by flags satisfying the first bilinear relation. Over $``$ we have the tautological flag of vectorbundles $`F^p`$. A complex variation underlying $`V`$ is the same thing as a section $`s`$ of $`𝒢`$ which in each fiber lands in the set of flags satisfying the second period relation. The degree of the $`p`$-th Hodge bundle is the degree of $`s^{}F^p`$, together making up the flag degree $`\{d_p\}`$, $`pI`$. Since these are all fixed, the desired collection of complex variations forms an open set in the variety $`S`$ of sections of the partial flag bundle $`𝒢`$ over $`C`$ of fixed flag degree. This variety $`S`$ can be viewed as the Hilbert scheme of curves in $`𝒢`$ of fixed multidegree with respect to the embeddings defined by the flags. Q.E.D.
References
\[Arakelov\] Arakelov, A.: Families of algebraic curves with fixed degeneracies, Izv. Ak. Nauk. S.S.S.R, ser. Math. 35 (1971) 1277–1302.
\[Deligne\] Deligne, P.: Un théorème de finitude pour la monodromie, in Discrete groups in geometry and analysis, Papers in honor of G. D. Mostow’s 60th Birthday, Prog. Math. 67 (1987) 1–19.
\[Eyss\] Eyssidieux, Ph.: La caractéristique d’Euler du complexe de Gauss-Manin, Journ. f. reine und angew. Math. 490 (1997) 155–212.
\[Faltings\] Faltings, G.: Arakelov’s theorem for abelian varieties, Invent. Math. 73 (1983) 337–348.
\[Griffiths\] Griffiths, P. A.: Topics in transcendental algebraic geometry, Princeton Univ. Press (1984) .
\[JostZuo\] Jost, J., K. Zuo: Arakelov type inequalities for Hodge bundles over algebraic varieties, Part I: Hodge bundles over algebraic curves with unipotent monodromies around singularities, Preprint (1999)
\[Peters84\] Peters, C.A.M. : A criterion for flatness of Hodge bundles over curves and geometric applications, Math. Ann. 268 (1984) 1–19.
\[Peters86\] Peters, C.A.M.: On Arakelov’s finiteness theorem for higher dimensional varieties, Rend. Sem. Mat. Univ. Politec. Torino (1986) 43–50.
\[Peters90\] Peters, C.A.M. : Rigidity for variations of Hodge structure and Arakelov-type finiteness theorems, Comp. Math 75 (1990) 113–126.
\[Saito\] Saito, M.-H.: Classification of non-rigid families of abelian varieties, Tohoku Math. J 45 (1993) 159–189.
\[SaitoZucker\] Saito, M.-H., S. Zucker: Classification of non-rigid families of K3-surfaces and a finiteness theorem of Arakelov type, Math. Ann. 289 (1991) 1–31.
\[Schmid\] Schmid, W.: Variation of Hodge structure: the singularities of the period mapping, Invent. Math. 22 (1973) 211–319.
\[Simpson90\] Simpson, C.: Harmonic bundles on non-compact curves, J. Am. Math. Soc. 3 (1990) 713–770.
\[Simpson92\] Simpson, C.: Higgs bundles and local systems, Publ. Math. I.H.E.S. 75 (1992) 5–95.
\[Simpson94\] Simpson, C.: Moduli of representations of the fundamental group of a smooth variety, Publ. Math. I.H.E.S. 79 (1994) 47–129.
\[Simpson95\] Simpson, C.: Moduli of representations of the fundamental group of a smooth projective variety. II, Publ. Math. I.H.E.S. 80 (1995) 5-79.
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# On the Eigenproblems of 𝒫𝒯-Symmetric Oscillators
## 1. Introduction
We are considering the eigenproblem
(1)
$$u^{\prime \prime }(x)+[P(x^2)(ix)^{2n+1}]u(x)=\lambda u(x)\text{for }\mathrm{}<x<\mathrm{}$$
with $`u(\pm \mathrm{})=0`$, where $`P(x)`$ is a polynomial of degree at most $`n1`$ with all nonnegative real coefficients (possibly $`P0`$).
This is an example of a class of problems, the so-called $`𝒫𝒯`$-symmetric non-Hermitian Hamiltonian problems, which have arisen in recent years in a number of physical contexts . D. Bessis conjectured in 1995 that:
###### Conjecture .
Eigenvalues of $`H=\frac{d^2}{dx^2}(ix)^3`$ are all real and positive.
Many numerical and asymptotic results support this conjecture. And later for $`n>1`$ it was conjectured that the equation (1) also has positive real eigenvalues, under different boundary conditions . However, there is no rigorous proof of this to date.
This paper is organized as follows: In Section 2, we prove that eigenvalues of the equation (1) lie in the sector $`|\mathrm{arg}\lambda |\frac{\pi }{2n+3}`$. This goes part way to proving that the eigenvalues are real and positive. We generalize this result to $`H=\frac{d^2}{dx^2}+[P(x^2)+ixQ(x^2)]`$ for some real polynomials $`P`$ and $`Q`$. In particular, for the potentials $`(ix)^3`$ and $`x^2+igx^3`$ with any real $`g`$, we have that $`|\mathrm{arg}\lambda |\frac{\pi }{5}`$. Then next in Section 3, for the case $`H=\frac{d^2}{dx^2}(ix)^3`$, we fairly precisely locate the zeros of the eigenfunctions and their first derivatives in the complex plane. Conversely we find a large zero-free region. In Section 4, still with $`H=\frac{d^2}{dx^2}(ix)^3`$, we find a large class of polynomials that are orthogonal to $`|u|^2`$ on each horizontal line. And finally in the last section, we discuss related open problems.
For the rest of Introduction, we provide some more background information on (1). First, a $`𝒫𝒯`$-symmetric Hamiltonian is a Hamiltonian which is invariant under the product of the parity operation $`𝒫(:xx)`$ and the time reversal operation $`𝒯(:ii)`$. Certainly (1) is $`𝒫𝒯`$-symmetric while, for example, $`\frac{d^2}{dx^2}+x(ix)^3`$ is not $`𝒫𝒯`$-symmetric. If $`H=\frac{d^2}{dx^2}+V(x)`$ is $`𝒫𝒯`$-symmetric, then $`\overline{V(x)}=V(x)`$ and so $`\mathrm{Re}V(x)`$ is an even function and $`\mathrm{Im}V(x)`$ is an odd function. Hence if $`V(x)`$ is a polynomial, then $`V(x)=P(x^2)+ixQ(x^2)`$ for some real polynomials $`P`$ and $`Q`$.
Next by the work of Caliceti et al. , it is known that the $`𝒫𝒯`$-symmetric Hamiltonian $`H=\frac{d^2}{dx^2}+x^2g(ix)^3`$ has discrete spectrum, for $`g`$ real, and these eigenvalues are positive real if $`g`$ is small enough. However, there are some $`𝒫𝒯`$-symmetric Hamiltonians that have no eigenvalues \[18, §1\], or non-real eigenvalues \[12, footnote on page 26\].
Lastly, for any $`\lambda `$ there are two linearly independent solutions of (1), if the boundary conditions are not imposed. In generic cases, the solutions blow up at both $`+\mathrm{}`$ and $`\mathrm{}`$, while in exceptional cases, the solutions decay to zero as $`x`$ approaches $`+\mathrm{}`$ or $`\mathrm{}`$. Only in very exceptional cases (when $`\lambda `$ is an eigenvalue!) does one find a solution that decays to zero at both $`+\mathrm{}`$ and $`\mathrm{}`$ (see Lemma 1 for details).
## 2. The eigenvalues lie in a sector
In this section, we prove that the eigenvalues $`\lambda `$ of (1) lie in the sector $`|\mathrm{arg}\lambda |\frac{\pi }{2n+3}`$ and we extend this result for more general cases. To do this we will use results of Hille \[16, §7.4\].
For any $`\lambda `$ the equation (1) without the boundary conditions allows two linearly independent solutions. If $`u(x)`$ solves the ODE (1), then since $`P(z^2)(iz)^{2n+1}`$ is an entire function (analytic in the whole complex plane), there exists an entire function $`u(z)`$ which agrees with $`u(x)`$ on the real line and satisfies $`u^{\prime \prime }(z)+[P(z^2)(iz)^{2n+1}]u(z)=\lambda u(z)`$. We begin by describing the asymptotic behavior of $`u`$ near infinity. Recall that $`\mathrm{deg}Pn`$.
###### Definition .
Let
(4) $`\theta _j=2\pi {\displaystyle \frac{j}{2n+3}}{\displaystyle \frac{\mathrm{arg}(i^{2n+1})}{2n+3}}=\{\begin{array}{cc}\hfill \frac{2\pi j\frac{\pi }{2}}{2n+3}& \text{if }n\text{ is even,}\hfill \\ \hfill \frac{2\pi j+\frac{\pi }{2}}{2n+3}& \text{if }n\text{ is odd.}\hfill \end{array}`$
We define Stokes regions
$$S_j=\{z:\theta _j<\mathrm{arg}z<\theta _{j+1}\},$$
for $`j=0,1,2,\mathrm{},2n+2`$. And for notational convenience, we define $`S_{j+2n+3}=S_j`$ for all $`j`$. Also we denote
$$S_{j,ϵ}=\{z:\theta _j+ϵ<\mathrm{arg}z<\theta _{j+1}ϵ\},$$
for $`0<ϵ<\frac{\pi }{2n+3}.`$
Notice $`\theta _j`$ is neither $`0`$ nor $`\pi `$. Thus the negative and the positive real axes lie within two of the Stokes regions (see Figure 1). We call these the left- and the right-hand Stokes regions, respectively. Also we call the rays $`\{\mathrm{arg}z=\theta _j\}`$ “critical rays”.
###### Lemma 1.
Every solution of $`u^{\prime \prime }(z)+[P(z^2)(iz)^{2n+1}]u(z)=\lambda u(z)`$ is asymptotic to
(5)
$$(const.)z^{\frac{2n+1}{4}}\mathrm{exp}[\pm \frac{2}{2n+3}(iz)^{\frac{2n+3}{2}}(1+o(1))]$$
as $`z\mathrm{}`$ in $`S_{j,ϵ}`$, for each $`0<ϵ<\frac{\pi }{2n+3}.`$ The error $`o(1)`$ is uniform in $`\mathrm{arg}z`$ in the sense that $`lim_r\mathrm{}sup\{|o(1)|:zS_{j,ϵ},|z|=r\}=0`$.
Also $`u`$ has infinitely many zeros in $``$ but only finitely many in $`_jS_{j,ϵ},`$ for each $`0<ϵ<\frac{\pi }{2n+3}`$.
The asymptotic expressions imply in particular that in each Stokes region, $`u(z)`$ either decays to 0 or blows up, as $`z`$ approaches infinity in $`S_{j,ϵ}`$.
###### Proof.
See Hille’s book \[16, §7.4\] for a proof of a more general result. An outline of the proof is as follows: Hille first transforms the equation into another complex $`Z`$-plane by using the Liouville transform. And then he compares $`u`$ with the solutions of the sine equation $`w^{\prime \prime }(Z)+w(Z)=0`$ and finally transforms back to the original complex $`z`$-plane. So the above asymptotic expressions are the asymptotic expressions for solutions of the sine equation (in the $`Z`$-variable) expressed in terms of the original $`z`$-variable. The Stokes regions are determined by the Liouville transformation.
Also we can deduce the last assertion of the theorem from \[16, §7.4\]. This is proved in \[13, Theorem 5\] for more general equations. ∎
###### Remark 1.
Under the Liouville transformation, a neighborhood of infinity in each Stokes region in the complex $`z`$-plane maps to a neighborhood of infinity in either the upper or lower half $`Z`$-plane. So if $`u`$ decays in a Stokes region $`S_j`$ for some $`j`$, then $`u`$ must blow up in the Stokes regions $`S_{j+1}`$ and $`S_{j1}`$. Otherwise, there would be a solution of the sine equation in the $`Z`$-plane which decays to zero in all directions. This is a contradiction. However, $`u`$ might blow up in many consecutive Stokes regions (even in all Stokes regions) (see \[16, §7.4\]).
###### Definition .
Let $`\lambda `$ and let $`u(z)0`$ be an analytic function on $``$ that satisfies (1). We say $`u`$ is an eigenfunction and $`\lambda `$ is an eigenvalue, for (1), if $`u(z)`$ decays to zero along rays to infinity in the left- and right-hand Stokes regions (that is, if $`u`$ has decaying asymptotics in (5), in these two regions).
###### Remark 2.
Given a Stokes region $`S_j`$, there always exists a solution of $`u^{\prime \prime }(z)+[P(z^2)(iz)^{2n+1}]u(z)=\lambda u(z)`$ that blows up in $`S_j`$ \[16, §7.4\]. So if there were two linearly independent eigenfunctions with the same eigenvalue, then all the solutions of $`u^{\prime \prime }(z)+[P(z^2)(iz)^{2n+1}]u(z)=\lambda u(z)`$ would satisfy $`u(\pm \mathrm{}+0i)=0`$ and there would be no solutions that blow up in the left- and right-hand Stokes regions. Thus there are no repeated eigenvalues, and all eigenvalues are simple.
###### Remark 3.
Note that if $`u(z)`$ is an eigenfunction with eigenvalue $`\lambda `$, then $`\overline{u}(\overline{z})`$ is an eigenfunction with eigenvalue $`\overline{\lambda }`$ (an upper bar denotes the complex conjugate). So if an eigenvalue is real then $`u(z)=c\overline{u}(\overline{z})`$ by Remark 2, and clearly $`|c|=1`$. Writing $`c=e^{2i\varphi }`$ and replacing $`u`$ by $`e^{i\varphi }u`$, we get that eigenfunctions with real eigenvalues are symmetric with respect to the imaginary axis.
The main result of this paper is:
###### Theorem 2.
If $`\lambda `$ is an eigenvalue of (1), then $`\lambda 0`$ and $`|\mathrm{arg}\lambda |\frac{\pi }{2n+3}`$.
That the eigenvalues have positive real part was known already (according to Mezincescu ); our proof below includes a very simple argument for this fact. In the proof and elsewhere, we will use the following:
Since $`u(z)`$ decays exponentially along rays to infinity in the left- and right-hand Stokes regions, so does $`u^{}`$ by the Cauchy integral formula. Therefore $`p(r)|u(re^{i\theta })|^2`$ and $`p(r)|u^{}(re^{i\theta })|^2`$ are integrable along the line $`rre^{i\theta }`$ in $``$ for any polynomial $`p(r)`$, provided $`|\theta |<\frac{\pi }{2(2n+3)}`$ (so that the ends of the line stay in the left- and right-hand Stokes regions).
###### Proof of Theorem 2.
Let $`u`$ be an eigenfunction with eigenvalue $`\lambda `$, so that
$$u^{\prime \prime }(z)+[P(z^2)+(iz)^{2n+1}]u(z)=\lambda u(z),$$
where $`P(z)=\mathrm{\Sigma }_{k=0}^na_kz^k`$ for some $`a_k0`$, $`k=0,1,2,\mathrm{},n`$.
Write
$$\lambda =\alpha +i\beta ,\alpha ,\beta .$$
Fix $`\theta `$ with $`|\theta |<\frac{\pi }{2(2n+3)}`$. Let $`v(r)=u(re^{i\theta })`$. Then $`v^{}(r)=u^{}(re^{i\theta })e^{i\theta }`$ and $`v^{\prime \prime }(r)=u^{\prime \prime }(re^{i\theta })e^{2i\theta }.`$ Thus our ODE becomes
$$v^{\prime \prime }(r)+\left\{[\alpha +i\beta P(r^2e^{2i\theta })]e^{2i\theta }+i^{2n+1}r^{2n+1}e^{i(2n+3)\theta }\right\}v(r)=0.$$
Then we multiply this by $`e^{i(2n+3)\theta }\overline{v}(r)`$, integrate and use integration by parts to get
$`e^{i(2n+3)\theta }{\displaystyle _{\mathrm{}}^{\mathrm{}}}|v^{}|^2𝑑r`$
$`=`$ $`(\alpha +i\beta )e^{i(2n+1)\theta }{\displaystyle _{\mathrm{}}^{\mathrm{}}}|v|^2𝑑r{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{i(2n+1)\theta }P(r^2e^{2i\theta })|v|^2𝑑r+i^{2n+1}{\displaystyle _{\mathrm{}}^{\mathrm{}}}r^{2n+1}|v|^2𝑑r,`$
for all $`|\theta |<\frac{\pi }{2(2n+3)}`$, where we note that the line $`re^{i\theta }`$ stays in the left- and right-hand Stokes regions where $`u`$ (and hence $`u^{}`$) decays exponentially to zero as $`z`$ approaches infinity.
Taking the real part of (2) gives (since $`|\theta |<\frac{\pi }{2(2n+3)}`$)
$`0`$ $`<`$ $`\mathrm{cos}(2n+3)\theta {\displaystyle _{\mathrm{}}^{\mathrm{}}}|v^{}|^2𝑑r`$
$`=`$ $`\left\{\alpha \mathrm{cos}(2n+1)\theta +\beta \mathrm{sin}(2n+1)\theta \right\}{\displaystyle _{\mathrm{}}^{\mathrm{}}}|v|^2𝑑r`$
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{Re}[e^{i(2n+1)\theta }P(r^2e^{2i\theta })]|v|^2𝑑r.`$
But $`\mathrm{Re}[e^{i(2n+1)\theta }P(r^2e^{2i\theta })]=\mathrm{\Sigma }_{k=0}^na_kr^{2k}\mathrm{cos}(2n2k+1)\theta 0`$ if $`a_k0`$ and $`|\theta |<\frac{\pi }{2(2n+1)}`$ (certainly true if $`|\theta |<\frac{\pi }{2(2n+3)}`$). So from (2) we conclude that
$$\alpha \mathrm{cos}(2n+1)\theta +\beta \mathrm{sin}(2n+1)\theta >0,$$
for all $`|\theta |<\frac{\pi }{2(2n+3)}`$. That is,
$$\alpha >|\beta |\mathrm{tan}(2n+1)\theta ,$$
for all $`0\theta <\frac{\pi }{2(2n+3)}.`$
Taking $`\theta =0`$ gives $`\alpha >0`$, in particular $`\lambda 0`$ and taking $`\theta \frac{\pi }{2(2n+3)}`$ gives
$$\alpha |\beta |\mathrm{tan}\frac{(2n+1)\pi }{2(2n+3)}.$$
Then finally using $`\mathrm{tan}\varphi =\mathrm{cot}(\frac{\pi }{2}\varphi )`$, we have
$$\mathrm{tan}\frac{\pi }{2n+3}\frac{|\beta |}{\alpha }.$$
That is, $`|\mathrm{arg}\lambda |\frac{\pi }{2n+3}.`$
###### Remark 4.
We can extend Theorem 2 by allowing $`P`$ to have some negative coefficients as long as $`P`$ satisfies $`\mathrm{Re}[e^{i(2n+1)\theta }P(r^2e^{2i\theta })]0`$ for $`|\theta |<\frac{\pi }{2(2n+3)}`$. For example, with $`n=3`$ and $`c`$, let $`P(z)=z^3+cz^2+z`$; then $`\mathrm{Re}[e^{7i\theta }P(r^2e^{2i\theta })]=r^2[r^4\mathrm{cos}\theta +cr^2\mathrm{cos}(3\theta )+\mathrm{cos}(5\theta )]`$. So if $`c^2\mathrm{cos}^2(3\theta )4\mathrm{cos}\theta \mathrm{cos}(5\theta )0`$ for $`|\theta |<\frac{\pi }{18}`$, i.e. $`|c|\sqrt{\frac{16}{3}\mathrm{cos}(\frac{\pi }{18})\mathrm{cos}(\frac{5\pi }{18})}`$ $`1.837`$, then $`\mathrm{Re}[e^{7i\theta }P(r^2e^{2i\theta })]0`$. So the theorem holds for this $`P`$ provided $`c\sqrt{\frac{16}{3}\mathrm{cos}(\frac{\pi }{18})\mathrm{cos}(\frac{5\pi }{18})}`$.
Also by simple change of variables, we get the same result for $`H=\frac{d^2}{dz^2}+[P(z^2)g(iz)^{2n+1}]`$ for any non-zero real $`g`$.
Moreover, by translations in $``$, we have the same result for $`H=\frac{d^2}{dz^2}+P((z\xi )^2)gi^{2n+1}(z\xi )^{2n+1}`$ for any $`\xi `$. For example, if $`u`$ solves $`u^{\prime \prime }(z)iz^3u(z)=\lambda u(z)`$, then $`v(z)=u(z+ai)`$ solves $`v^{\prime \prime }(z)+[(3az^2a^3)iz(z^23a^2)]v(z)=\lambda v(z)`$ for any real number $`a`$. Observe $`v`$ still satisfies the boundary conditions: $`v(\pm \mathrm{}+0i)=u(\pm \mathrm{}+ai)=0`$.
###### Remark 5.
The readers should notice that our boundary conditions are different, for $`n2`$, from those Bender and Boettcher take. In , the zero boundary conditions of the problems $`u^{\prime \prime }(iz)^Nu=\lambda u`$ for $`N4`$ are taken not on Stokes regions containing the real axis but instead on Stokes regions which are near the negative imaginary axis for large $`N`$.
The next theorem extends Theorem 2.
###### Theorem 3.
Let $`\lambda `$ and $`n1`$. Suppose that $`u`$ solves the ODE
(8)
$$u^{\prime \prime }[P(z^2)+izQ(z^2)]u=\lambda u,u(\pm \mathrm{}+0i)=0,$$
for some real polynomials $`P(z)=\mathrm{\Sigma }_{k=0}^na_kz^k`$ and $`Q(z)=\mathrm{\Sigma }_{k=0}^nb_kz^k`$ with all nonnegative $`a_k`$ and with $`b_n\{0\}`$. If for all $`k<n`$ the coefficients $`a_k,b_k`$ satisfy
(9)
$$\frac{\mathrm{sin}^2(2n2k)\theta }{\mathrm{cos}(2n2k+1)\theta \mathrm{cos}(2n2k1)\theta }b_k^2\{\begin{array}{cc}\hfill 4a_ka_{k+1}& \text{if }n=1\text{ and }k=0\hfill \\ \hfill 2a_ka_{k+1}& \text{if }n>1\text{ and }k=0,n1\hfill \\ \hfill a_ka_{k+1}& \text{if }n>1\text{ and }1kn2\hfill \end{array}$$
at $`\theta =\frac{\pi }{2(2n+3)}`$, then $`|\mathrm{arg}\lambda |\frac{\pi }{2n+3}`$.
For $`n=3`$, the coefficients of $`b_k^2`$ in (9) are approximately $`3.41,0.74`$ and $`0.14`$ for $`k=0,1,2,`$ respectively.
Theorem 3 contains Theorem 2, just by taking $`b_k=0`$ for $`k=0,1,\mathrm{},n1`$ (in which case (9) is trivially satisfied).
###### Proof.
The main idea of the proof is the same as that of the proof of Theorem 2. Even if the equation (8) is little different from the equation for (5), Stokes regions for (8) are the same as for (5) (if $`b_n`$ has the same sign as $`(1)^{n+1}`$) or else are rotated by $`180^{}`$ (if $`b_n`$ has the opposite sign). See Section $`7.4`$ in for details. But in either case, the lines $`rre^{i\theta }`$ with $`|\theta |<\frac{\pi }{2(2n+3)}`$ lie within the left- and right-hand Stokes regions, where we impose the zero boundary conditions. And this gives the integrabilities in the proof.
Let $`v(r)=u(re^{i\theta })`$. Then like we derived (2) in the proof of Theorem 2 we have
$`\left\{\alpha \mathrm{cos}(2n+1)\theta +\beta \mathrm{sin}(2n+1)\theta \right\}{\displaystyle _{\mathrm{}}^{\mathrm{}}}|v|^2𝑑r`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\{\mathrm{cos}[(2n+3)\theta ]|v^{}|^2+\mathrm{\Sigma }_{k=0}^na_k\mathrm{cos}[(2n2k+1)\theta ]r^{2k}|v|^2`$
$`+\mathrm{\Sigma }_{k=0}^{n1}b_k\mathrm{sin}[(2n2k)\theta ]r^{2k+1}|v|^2\}dr,\text{for }|\theta |<{\displaystyle \frac{\pi }{2(2n+3)}}.`$
Since $`a_k0`$ for all $`k`$, we get $`\alpha >0`$ by letting $`\theta =0`$ in (2). (This is true for any $`Q`$ with real coefficients and $`\mathrm{deg}Q=n`$.)
If we find conditions on $`a_k`$ and $`b_k`$ such that
(11)
$$\mathrm{\Sigma }_{k=0}^na_k\mathrm{cos}[(2n2k+1)\theta ]r^{2k}|v|^2+\mathrm{\Sigma }_{k=0}^{n1}b_k\mathrm{sin}[(2n2k)\theta ]r^{2k+1}|v|^20,$$
for every $`r`$ and $`|\theta |<\frac{\pi }{2(2n+3)}`$, then we have from (2) with $`\theta \pm \frac{\pi }{2(2n+3)}`$, that
$$\alpha \mathrm{cos}\frac{(2n+1)\pi }{2(2n+3)}\pm \beta \mathrm{sin}\frac{(2n+1)\pi }{2(2n+3)}0.$$
So then $`|\mathrm{arg}\lambda |\frac{\pi }{2n+3}`$ as desired, like in the proof of Theorem 2.
When $`n3`$, we can rewrite the expression in (11) as $`|v|^2`$ times
$`\left\{a_0\mathrm{cos}(2n+1)\theta +rb_0\mathrm{sin}2n\theta +r^2{\displaystyle \frac{a_1}{2}}\mathrm{cos}(2n1)\theta \right\}`$
(12) $`+\mathrm{\Sigma }_{k=1}^{n2}\left\{{\displaystyle \frac{a_k}{2}}\mathrm{cos}(2n2k+1)\theta +rb_k\mathrm{sin}(2n2k)\theta +r^2{\displaystyle \frac{a_{k+1}}{2}}\mathrm{cos}(2n2k1)\theta \right\}r^{2k}`$
$`+\left\{{\displaystyle \frac{a_{n1}}{2}}\mathrm{cos}3\theta +rb_{n1}\mathrm{sin}2\theta +r^2a_n\mathrm{cos}\theta \right\}r^{2n2}.`$
Now (12) is nonnegative if each quadratic in $`r`$ has non-positive discriminant:
$`b_0^2\mathrm{sin}^22n\theta 2a_0a_1\mathrm{cos}(2n+1)\theta \mathrm{cos}(2n1)\theta `$ $``$ $`0,`$
$`b_k^2\mathrm{sin}^2(2n2k)\theta a_ka_{k+1}\mathrm{cos}(2n2k+1)\theta \mathrm{cos}(2n2k1)\theta `$ $``$ $`0\text{for}1kn2,`$
$`b_{n1}^2\mathrm{sin}^22\theta 2a_{n1}a_n\mathrm{cos}3\theta \mathrm{cos}\theta `$ $``$ $`0,`$
which is (9). The coefficients of $`b_k^2`$ in (9) are all increasing functions of $`0\theta <\frac{\pi }{2(2n+3)}`$, and so it suffices that (9) hold at $`\theta =\frac{\pi }{2(2n+3)}`$.
Now when $`n=1,2`$, it is easy to see similarly that the theorem holds. This completes the proof. ∎
###### Remark 6.
In (11), the sign of $`_{\mathrm{}}^{\mathrm{}}r^{2k+1}|v|^2𝑑r`$ is difficult to determine because $`r`$ can be negative as well as positive.
The conditions in Theorem 3 are sufficient but not necessary, as is clear from the proof.
We used $`a_k=\frac{a_k}{2}+\frac{a_k}{2}`$ to get (12) from (2). If we use $`a_k=\delta _ka_k+(1\delta _k)a_k`$ for some $`0<\delta _k<1`$, we will get new sufficient conditions for the theorem.
## 3. The zero–free region for $`u`$ and $`u^{}`$
The results in the previous section are based on the eigenfunction $`u`$ decaying to zero as $`z`$ approaches infinity in the left- and the right-hand Stokes regions. So consideration of the finite zeros of $`u`$ may be useful for further results on our eigenproblem.
For the next two sections, we will suppose $`H=\frac{d^2}{dz^2}(iz)^3`$. See Figure 2 for the asymptotic behavior of the eigenfunction $`u`$.
In this section, we provide a zero-free region for the eigenfunction $`u`$ of
(13)
$$u^{\prime \prime }iz^3u=\lambda u\text{with }u(\pm \mathrm{}+0i)=0$$
and for its derivative $`u^{}`$. And we give some answers on how zeros of the eigenfunction should be arranged in $``$.
It is obvious that $`u`$ and $`u^{}`$ do not share a common zero. Otherwise, by (13), all the derivatives of $`u`$ and $`u`$ itself would vanish at the zero, and so $`u0`$.
The following lemma is needed for our argument. Recall $`\lambda =\alpha +i\beta `$.
###### Lemma 4.
Let $`z:[c,d]`$ be a smooth curve with $`z^{}(t)0`$ for $`t[c,d]`$. If $`u`$ solves (13), then writing $`z(t)=x(t)+iy(t)`$,
$`\mathrm{Re}(u^{}\overline{u})|_{z(c)}^{z(d)}`$
$`=`$ $`{\displaystyle _c^d}x^{}|u_x(z(t))|^2𝑑t+{\displaystyle _c^d}\left[x^{}\mathrm{Re}(iz^3(t)\lambda )y^{}\mathrm{Im}(iz^3(t)\lambda )\right]|u(z(t))|^2𝑑t,`$
and
$`\mathrm{Im}(u^{}\overline{u})|_{z(c)}^{z(d)}`$
$`=`$ $`{\displaystyle _c^d}y^{}|u_x(z(t))|^2𝑑t+{\displaystyle _c^d}\left[y^{}\mathrm{Re}(iz^3(t)\lambda )+x^{}\mathrm{Im}(iz^3(t)\lambda )\right]|u(z(t))|^2𝑑t.`$
Hille calls this lemma the Green’s transform \[16, §11.3\], and he uses it to get information on zero-free regions of solutions of linear second order equations (mainly with coefficient functions that are real on the real line).
###### Proof.
Let $`f(t)=u(z(t))`$ for $`t[c,d]`$. Then $`f^{}(t)=z^{}(t)u^{}(z(t))`$ and
$$(\frac{f^{}(t)}{z^{}(t)})^{}=z^{}(t)u^{\prime \prime }(z(t))=z^{}(t)[iz^3(t)\lambda ]f(t).$$
Hence by integration by parts,
$$\left(\frac{f^{}(t)}{z^{}(t)}\right)\overline{f}(t)|_c^d=_c^d\frac{|f^{}|^2}{z^{}}𝑑t+_c^dz^{}[iz^3\lambda ]|f|^2𝑑t.$$
Now by the formula $`f^{}(t)=z^{}(t)u^{}(z(t))`$ and splitting real and imaginary parts of the above, we get the lemma. ∎
Now we examine the consequences of the lemma. First, if $`\mathrm{Re}(u^{}\overline{u})`$ were not one-to-one on the imaginary axis, that would imply that the eigenvalue would be real by (4) with $`z(t)=it`$.
###### Remark 7.
Second, another immediate consequence of Lemma 4 is that on any vertical line segments on which $`\mathrm{Im}(iz^3\lambda )`$ doesn’t change its sign, $`\mathrm{Re}(u_x\overline{u})`$ as a function of $`y`$ is one-to-one. On horizontal line segments on which $`\mathrm{Im}(iz^3\lambda )`$ doesn’t change its sign, $`\mathrm{Im}(u_x\overline{u})`$ as a function of $`x`$ is one-to-one (Mezincescu \[18, §3\] observed this last fact on the real axis, where $`y0`$). These observations are special cases of Hille’s Theorem 11.3.3 in .
Third, let us define open regions $`A_j`$ and $`B_j`$, $`j=1,2,3,4`$ as in Figure 3 and 4. The following two theorems provide a large zero-free region for an eigenfunction $`u`$ of (13) and its derivative $`u^{}`$, assuming $`\lambda `$ is non-real. Perhaps these theorems might help show that $`\lambda `$ must actually be real. The underlying ideas of the proofs are taken from Hille’s book \[16, §11.3\].
###### Theorem 5.
If $`\beta :=\mathrm{Im}\lambda >0`$ then $`\mathrm{Im}(u^{}\overline{u})<0`$ on
$$B_1\{zB_4:\mathrm{Re}z\sqrt[3]{\beta /2}\}\{zA_1:\mathrm{Im}z\sqrt[3]{\beta /2}\},\text{see Figure }\text{5}\text{.}$$
Mezincescu \[18, §3\] has previously observed that the eigenfunction has no zeros on the real axis, which obviously lies in the shaded region of Figure 5. Moreover, we see that all the zeros of $`u`$ and $`u^{}`$ in $`\mathrm{Im}z0`$ must be in $`A_1`$.
Note that the lowest point in the closure $`cl(B_2)`$ of $`B_2`$ is $`\sqrt[3]{\beta /2}+i\sqrt[3]{\beta /2}`$.
###### Proof of Theorem 5.
For any $`y`$, by (4) with $`z(t)=t+iy`$ and by $`u(\pm \mathrm{}+iy)=0=u^{}(\pm \mathrm{}+iy)`$, we have that
$$\mathrm{Im}[u^{}(x+iy)\overline{u}(x+iy)]=_{\mathrm{}}^x\mathrm{Im}(iz(t)^3\lambda )|u|^2𝑑t,$$
and this is negative for $`x+iyB_4`$ with $`|y|\sqrt[3]{\beta /2}`$, because then $`z(t)=t+iyB_4`$ for all $`\mathrm{}<t<x`$ and so $`\mathrm{Im}(iz(t)^3\lambda )<0`$.
This argument also shows that $`\mathrm{Im}(u^{}\overline{u})<0`$ in $`\{zB_4:\mathrm{Re}z\sqrt[3]{\beta /2}\}`$; see Figure 5.
Similarly in $`B_1`$, for all $`y`$ we have that $`\mathrm{Im}(u^{}\overline{u})=_x^{\mathrm{}}\mathrm{Im}(iz(t)^3\lambda )|u(t+iy)|^2𝑑t<0`$.
For $`zA_1\text{ with }\mathrm{Im}z\sqrt[3]{\beta /2}`$ (so that $`zA_4`$), we use (4) along vertical line segments starting from points on the line $`\mathrm{Im}z=\sqrt[3]{\beta /2}`$ to conclude $`\mathrm{Im}(u^{}\overline{u})<0`$ in this region. ∎
Note that $`\mathrm{Im}(u^{}\overline{u})=\frac{1}{2}\frac{}{y}|u(x+iy)|^2`$. So in the region in Theorem 5, $`|u(x+iy)|`$ is an increasing function of $`y.`$
###### Theorem 6.
Assume $`\beta >0`$. Then
(i) $`\mathrm{Re}(u^{}\overline{u})>0`$ on the union of the regions $`A_2`$, the region below $`A_2`$ and the region $`B_4`$ between $`A_2`$ and $`B_2`$ with the real part less than or equal to that of the zero $`\omega _3`$ of $`iz^3\lambda `$ in the third quadrant. See Figure 7.
(ii) $`\mathrm{Re}(u^{}\overline{u})<0`$ on the union of the regions $`A_3`$, the region below $`A_3`$, and the region in $`B_1`$ with the real part greater than or equal to that of the zero $`\omega _4`$ of $`iz^3\lambda `$ in the fourth quadrant. See Figure 7.
Obviously $`iz^3\lambda `$ has three zeros. When $`\beta >0`$, one of the zeros is in the second quadrant, one $`\omega _3`$ in the third and one $`\omega _4`$ in the fourth quadrant. Certainly these are the three points at which the boundaries of the $`A_i`$ and $`B_i`$ intersect.
Theorem 2 with $`n=1`$ shows that $`\frac{|\beta |}{\alpha }\mathrm{tan}\frac{\pi }{5}0.73.`$ It is easy to see that the rightmost point of $`cl(A_2)`$ is $`\sqrt[3]{\frac{\alpha }{2}}(1i)`$, at which $`\mathrm{Im}(iz^3\lambda )=x^33xy^2\beta >0`$ by $`0<\beta <\alpha `$. Thus the rightmost point of $`cl(A_2)`$ lies inside $`B_3`$ as shown in Figure 7. Similarly, the leftmost point of $`cl(A_3)`$ is $`\sqrt[3]{\frac{\alpha }{2}}(1i)`$, at which $`\mathrm{Im}(iz^3\lambda )=x^33xy^2\beta <0`$. Thus the leftmost point of $`cl(A_3)`$ lies outside $`B_1`$ as shown in Figure 7.
###### Proof of Theorem 6.
In the regions $`A_2`$ and $`A_3`$, we use (4) with horizontal lines to infinity to get the statements in parts $`(i)`$ and $`(ii)`$ of this theorem. In the region $``$ between $`A_2`$ and $`B_2`$ with the real part less than or equal to that of the zero of $`iz^3\lambda `$ in the third quadrant (see Figure 7), we use (4) with vertical lines $`z(t)=x+it`$ to show that $`\mathrm{Re}(u^{}\overline{u})>0`$. That is, we use $`\mathrm{Re}(u^{}\overline{u})|_{x+ic}^{x+id}=_c^d\mathrm{Im}(iz(t)^3\lambda )|u(x+it)|^2𝑑t`$. If $`x+id`$, we can find $`x+iccl(A_2B_3^c)`$, so that $`\mathrm{Re}[u^{}(x+ic)\overline{u}(x+ic)]>0`$, and $`\mathrm{Im}(iz(t)^3\lambda )>0`$. Hence, the above integral is an increasing function of $`d`$. Hence we have the desired result in this region $``$.
The region below $`A_2`$ is contained in $`B_3`$ since the rightmost point of $`cl(A_2)`$ lies in $`B_3`$ (see Figure 7). So a similar argument shows that $`\mathrm{Re}(u^{}\overline{u})>0`$ in the region below $`A_2`$. Also, the region below $`A_3`$ is contained in $`B_4`$ (see Figure 7) and so modified arguments show that the other statements of this theorem in part $`(ii)`$ are true. ∎
###### Corollary 7.
When $`\mathrm{Im}\lambda =\beta >0`$, the zero-free region of $`u`$ and $`u^{}`$ contains the union of the three shaded regions in Figure 5, 7 and 7.
Note that in case $`\mathrm{Im}\lambda =\beta <0`$ we can get similar theorems corresponding to the above two, since $`\overline{u}(\overline{z})`$ is an eigenfunction with eigenvalue $`\overline{\lambda }`$. The regions involved are simply the reflections of the above with respect to the imaginary axis.
In case $`\beta =0`$, so that $`\lambda `$ is real and $`\lambda =\alpha >0`$, the regions $`B_1,B_2,B_3`$ degenerate to the sectors $`\{\frac{\pi }{6}<\mathrm{arg}z<\frac{\pi }{6}\},\{\frac{\pi }{2}<\mathrm{arg}z<\frac{5\pi }{6}\},\{\frac{5\pi }{6}<\mathrm{arg}z<\frac{\pi }{2}\}`$ respectively, and we get the following theorem on zero-free regions.
###### Theorem 8.
Suppose $`\lambda `$ is real. Then $`\mathrm{Im}(u^{}\overline{u})<0`$ on $`\{\frac{\pi }{6}\mathrm{arg}z\frac{7\pi }{6}\}A_1`$ (which is a degenerate case of Figure 5), while $`\mathrm{Re}(u^{}\overline{u})`$ behaves as in Figure 7 and 7 with $`B_1,B_2,B_3`$ being sectors as above.
Also $`\mathrm{Re}(u^{}\overline{u})<0`$ in $`cl(A_1)\{\mathrm{Re}z<0\}`$ and $`\mathrm{Re}(u^{}\overline{u})>0`$ in $`cl(A_1)\{\mathrm{Re}z>0\}`$.
###### Corollary 9.
When $`\lambda `$ is real, the zero-free region of $`u`$ and $`u^{}`$ contains the union of all regions in Theorem 8; see Figure 8. That is, $`u`$ and $`u^{}`$ can only have zeros in $`\{iy:y>\sqrt[3]{\lambda }\}`$
$$\{zA_4:\frac{5\pi }{6}<\mathrm{arg}z<\frac{\pi }{6},\mathrm{Im}z>\sqrt[3]{\lambda /2}\}\{z:|\mathrm{Re}z|<\sqrt[3]{\lambda /2},\mathrm{Im}z\sqrt[3]{\lambda /2}\}.$$
###### Remark 8.
Bender et al. find numerically that $`u`$ has some zeros along an “arch” within the unshaded region in Figure 8, when $`\lambda `$ is real.
In proving Theorem 8, we will use the following lemma.
###### Lemma 10.
Suppose $`\zeta A_1`$, $`\mathrm{Re}(u^{}\overline{u})=0`$ at $`\zeta `$. Then $`\mathrm{Re}(u^{}\overline{u})<0`$ at all $`\zeta tA_1,t>0`$ and $`\mathrm{Re}(u^{}\overline{u})>0`$ at all $`\zeta +tA_1,t>0`$.
Note that there is no restriction on the sign of $`\mathrm{Im}\lambda `$, in this lemma.
###### Proof.
On horizontal line segments $`z(t)=t+iy`$ in $`A_1`$, (4) becomes
$$\mathrm{Re}(u^{}\overline{u})|_{c+iy}^{d+iy}=_c^d|u_x(z(t))|^2𝑑t+_c^d\mathrm{Re}(iz^3(t)\lambda )|u(z(t))|^2𝑑t.$$
Since $`\mathrm{Re}(iz^3(t)\lambda )>0`$ in $`A_1`$, $`\mathrm{Re}(u^{}\overline{u})`$ is a strictly increasing function of $`x`$ on each horizontal line segment in $`A_1`$. ∎
###### Proof of Theorem 8.
The proofs of Theorems 5 and 6 give everything except the last statement of the theorem. For that, recall we can take $`u(z)=\overline{u}(\overline{z})`$ by Remark 3; this implies $`\mathrm{Re}(u^{}\overline{u})`$ is an odd function with respect to reflection in the imaginary axis, so $`\mathrm{Re}(u^{}\overline{u})=0`$ on the whole imaginary axis. Now we use Lemma 10 to complete the proof. ∎
By the last statement of Lemma 1, in the sector $`S_{1,\frac{\pi }{20}}`$ that contains the negative imaginary axis, the eigenfunction $`u`$ has only finitely many zeros. Now with the zero-free region in the Theorems 6 and 8, we see that $`u`$ has only finitely many zeros in $`\mathrm{Im}z<0`$. Since $`u`$ has infinitely many zeros, $`u`$ must have infinitely many zeros in $`\mathrm{Im}z0`$. When $`\beta >0`$ (hence when $`\beta <0`$ as well), by Theorem 5, $`u`$ must have infinitely many zeros in $`A_1`$. Also when $`\beta =0`$, by Theorem 8, $`u`$ has infinitely many zeros on the positive imaginary axis.
The next theorem gives some information on how zeros of $`u`$ and $`u^{}`$ in $`A_1`$ should be arranged, when $`\beta >0`$. Note that all the zeros of $`u`$ and $`u^{}`$ in $`\mathrm{Im}z0`$ lie in $`A_1`$ by Theorem 5.
###### Theorem 11.
Suppose $`u(z)`$ is an eigenfunction of (13) with eigenvalue $`\lambda `$ with $`\mathrm{Im}\lambda =\beta >0`$. Then
(i) $`\mathrm{Re}(u^{}\overline{u})0`$ for some point on the imaginary axis if and only if $`uu^{}`$ has infinitely many zeros in $`A_1B_2`$ and at most finitely many zeros in $`A_1B_2^c`$; and
(ii) $`\mathrm{Re}(u^{}\overline{u})<0`$ for every point on the imaginary axis if and only if $`uu^{}`$ has no zeros in $`\{zA_1:\mathrm{Re}z0\}`$ and infinitely many in $`\{zA_1:\mathrm{Re}z>0\}`$.
We will use the following lemma along with Lemma 10.
###### Lemma 12.
Assume $`\mathrm{Im}\lambda =\beta >0`$. Suppose $`\mathrm{Re}\zeta _1\mathrm{Re}\zeta _2`$ and $`\mathrm{Re}(u^{}\overline{u})=0`$ at $`\zeta _1`$, $`\zeta _2`$ $`(\text{where }\zeta _1\zeta _2)`$. Then
(i) $`\zeta _1,\zeta _2cl(A_1B_2)`$ $`\mathrm{Re}\zeta _1<\mathrm{Re}\zeta _2`$ and $`\mathrm{Im}\zeta _1<\mathrm{Im}\zeta _2`$, and
(ii) $`\zeta _1,\zeta _2cl(A_1B_2^c)`$ $`\mathrm{Re}\zeta _1<\mathrm{Re}\zeta _2`$ and $`\mathrm{Im}\zeta _1>\mathrm{Im}\zeta _2`$.
###### Proof of part (i).
We will first prove this for $`\zeta _1,\zeta _2A_1B_2`$. Suppose that $`\mathrm{Re}\zeta _1=\mathrm{Re}\zeta _2`$. Then we could find a vertical line segment $`z(t)`$ in $`A_1B_2`$ whose end points are $`\zeta _1`$ and $`\zeta _2`$. We apply (4) to this line segment to get
$$0=_c^d\mathrm{Im}(iz^3(t)\lambda )|u(z(t))|^2𝑑t.$$
This would imply $`u0`$ on the curve $`z(t)`$ since $`\mathrm{Im}(iz^3(t)\lambda )>0`$ in $`B_2`$. So then since $`u`$ is analytic, $`u0`$ in $``$. This is a contradiction. Hence $`\mathrm{Re}\zeta _1<\mathrm{Re}\zeta _2`$.
Similarly, suppose that $`\mathrm{Im}\zeta _1\mathrm{Im}\zeta _2`$. Then we could find a smooth curve $`z(t)=x(t)+iy(t)`$ in $`A_1B_2`$ such that $`z(c)=\zeta _1`$, $`z(d)=\zeta _2`$, $`x^{}(t)>0`$ and $`y^{}(t)0`$. Note that $`\mathrm{Im}(iz^3(t)\lambda )>0`$ and $`\mathrm{Re}(iz^3(t)\lambda )>0`$ in $`A_1B_2`$. This contradicts (4) like for the case of $`\mathrm{Re}\zeta _1=\mathrm{Re}\zeta _2`$. We now see that the above argument still holds for $`\zeta _1,\zeta _2cl(A_1B_2)`$.
Proof of part (ii). We use (4) again and a similar argument like in the proof of part $`(i)`$. ∎
###### Proof of Theorem 11.
Suppose $`u(z)`$ is an eigenfunction of (13) with eigenvalue $`\lambda `$ with $`\mathrm{Im}\lambda =\beta >0`$. Since $`u`$ has infinitely many zeros in $`A_1`$ (by the paragraph shortly before Theorem 11), certainly $`uu^{}`$ also has infinitely many zeros in $`A_1`$.
Proof of part (i). Suppose that $`\mathrm{Re}(u^{}\overline{u})0`$ for some point $`iy`$ on the imaginary axis. By (4) with $`z(t)=it`$, we have that
(16)
$$\mathrm{Re}(u^{}\overline{u})|_{ic}^{id}=\beta _c^d|u(iy)|^2𝑑y.$$
So then since $`\beta >0`$, $`\mathrm{Re}(u^{}\overline{u})>0`$ at every point $`it`$ for $`t>y`$. Now by Lemma 10, we have that $`\mathrm{Re}(u^{}\overline{u})>0`$ at every $`x+itA_1`$ for $`t>y`$ and $`x0`$. Thus $`uu^{}`$ does not have any zeros in $`\{zA_1B_2^c:\mathrm{Re}z0,\mathrm{Im}z>y\}`$.
The entire function $`uu^{}`$ does not have infinitely many zeros in any bounded region. So if $`uu^{}`$ had infinitely many zeros in $`A_1B_2^c`$, then $`uu^{}`$ would have infinitely many zeros in $`\{zA_1B_2^c:\mathrm{Re}z<0\}`$. But if $`uu^{}`$ has a zero $`z_1`$ in $`\{zA_1B_2^c:\mathrm{Re}z<0\}`$, then by Lemma 12 $`(ii)`$, $`uu^{}`$ has no zeros in $`\{zA_1B_2^c:\mathrm{Re}z_1\mathrm{Re}z<0\}`$. So then $`uu^{}`$ would have infinitely many zeros in a bounded region. This is a contradiction. Thus $`uu^{}`$ has infinitely many zeros in $`A_1B_2`$ and at most finitely many zeros in $`A_1B_2^c`$.
Conversely, suppose that $`uu^{}`$ has infinitely many zeros in $`A_1B_2`$ and at most finitely many zeros in $`A_1B_2^c`$. Choose a zero $`z_0`$ in $`A_1B_2`$. Then by Lemma 10, we see that $`\mathrm{Re}(u^{}\overline{u})>0`$ at $`i\mathrm{Im}z_0`$ since $`\mathrm{Re}z_0<0`$.
Proof of part (ii). Suppose that $`\mathrm{Re}(u^{}\overline{u})<0`$ for every point on the imaginary axis. Then by Lemma 10, $`\mathrm{Re}(u^{}\overline{u})<0`$ for every point in $`\{zA_1:\mathrm{Re}z0\}`$. So then $`uu^{}`$ has no zeros in $`\{zA_1:\mathrm{Re}z0\}`$. Now since we know that $`uu^{}`$ has infinitely many zeros in $`A_1`$, $`uu^{}`$ must have infinitely many zeros in $`\{zA_1:\mathrm{Re}z>0\}`$.
Conversely, suppose $`\mathrm{Re}(u^{}\overline{u})0`$ for some point on the imaginary axis. Then $`uu^{}`$ would have at most finitely many zeros in $`\{zA_1:\mathrm{Re}z>0\}`$ by the argument as in the proof of part $`(i)`$. This completes the proof. ∎
###### Remark 9.
Since the negative imaginary axis is in the middle of a blowing-up Stokes region (see Figure 2), $`u(iy)`$ blows up as $`y`$ tends to $`\mathrm{}`$. On the other hand, the positive imaginary axis is a critical ray. We can show that $`|u(iy)|^2(const.)y^{\frac{3}{2}}`$ for all $`y`$ near positive infinity, by Theorem 7.4.4 in .
So the right-hand side of (16) approaches $`+\mathrm{}`$ as $`c`$ tends to $`\mathrm{}`$ (while $`d`$ is fixed). Thus we see that $`\mathrm{Re}[u^{}(ic)\overline{u}(ic)]<0`$ for all $`c`$ near negative infinity. However, the right-hand side of (16) is convergent as $`d`$ tends to $`+\mathrm{}`$ (while $`c`$ is fixed). So $`\mathrm{Re}(u^{}\overline{u})`$ may or may not become positive near infinity along the positive imaginary axis.
The next lemma gives some information on zeros of $`u`$ and $`u^{}`$ in $`\mathrm{Im}z<0`$, if any exist. There can only be finitely many such zeros, by the paragraph shortly before Theorem 11.
###### Lemma 13.
Assume $`\mathrm{Im}\lambda =\beta 0`$. Suppose $`\mathrm{Im}\zeta _1\mathrm{Im}\zeta _2`$ and $`\mathrm{Im}(u^{}\overline{u})=0`$ at $`\zeta _1`$, $`\zeta _2`$ $`(\text{where }\zeta _1\zeta _2)`$. Then
(i) $`\zeta _1,\zeta _2cl(A_4B_3)`$ $`\mathrm{Im}\zeta _1<\mathrm{Im}\zeta _2`$ and $`\mathrm{Re}\zeta _1<\mathrm{Re}\zeta _2`$, and
(ii) $`\zeta _1,\zeta _2cl(A_4B_4)`$ $`\mathrm{Im}\zeta _1<\mathrm{Im}\zeta _2`$ and $`\mathrm{Re}\zeta _1>\mathrm{Re}\zeta _2`$.
###### Proof.
We omit the proof because it is very similar to the proof of Lemma 12. We use (4) instead of (4), and also make use of Figures 7 and 7. ∎
Roughly speaking, then, the zeros move up and to the right in the third quadrant, and down and to the right in the fourth quadrant. This observation supports that when $`\lambda `$ is real, zeros of $`u`$ in $`\mathrm{Im}z<0`$ lie on an arch-shaped curve as in Figures 5 and 6 in .
## 4. Other properties of eigenfunctions
Here we present a possible way of proving the conjecture that the eigenvalues $`\lambda `$ of $`H=\frac{d^2}{dz^2}(iz)^3`$ are positive real. Given an eigenfunction $`u`$ with eigenvalue $`\lambda `$, Theorem 14 below gives a class $`𝒪`$ of polynomials $`p(x,y)`$ which are orthogonal to $`|u|^2`$ in the sense that $`_{\mathrm{}}^{\mathrm{}}p(x,y)|u(x+iy)|^2𝑑x=0`$ for all $`y`$. One can perhaps prove the conjecture as follows. Suppose $`\mathrm{Im}\lambda 0`$; if $`𝒪`$ is large enough then $`|u|^20`$, giving a contradiction.
Let $`u`$ be an eigenfunction of $`H=\frac{d^2}{dz^2}(iz)^3`$ with eigenvalue $`\lambda =\alpha +i\beta `$.
###### Theorem 14.
Let $`𝒪=\{\text{polynomials }p(,):_{\mathrm{}}^{\mathrm{}}p(x,y)|u(x+iy)|^2𝑑x=0\text{ for all }y\}`$. Then:
(i) $`x^33xy^2\beta 𝒪`$,
(ii) for all $`m0`$,
$`{\displaystyle \frac{4}{m+1}}({\displaystyle \frac{x^{m+5}}{m+5}}3y^2{\displaystyle \frac{x^{m+3}}{m+3}}\beta {\displaystyle \frac{x^{m+2}}{m+2}})(x^33y^2x\beta )`$
$`m(m1)x^{m2}4x^m(3x^2yy^3+\alpha ){\displaystyle \frac{12}{m+1}}yx^{m+2}`$ $``$ $`𝒪,`$
(iii) if $`p𝒪`$ then $`p_y+2(x^33xy^2\beta )_0^xp(t,y)𝑑t𝒪`$, and
(iv) if $`p𝒪`$ then $`p_{xx}+p_{yy}+12x^2yp+4(x^33xy^2\beta )_0^xp_y(t,y)𝑑t𝒪`$.
For example the following polynomials are in $`𝒪`$:
$`p_3(x,y)`$ $`=`$ $`x^33xy^2\beta ,\text{ by }(i),`$
$`p_7(x,y)`$ $`=`$ $`x^79x^5y^25x^4\beta +18x^3y^4+18x^2y^2\beta +4x(\beta ^23y),`$
$`\text{by applying }(iii)\text{ to }p_3\text{ and multiplying by 2},`$
$`p_8(x,y)`$ $`=`$ $`2x^816x^6y^27x^5\beta +30x^4y^4+25x^3y^2\beta +5x^2(\beta ^212y)+10y^310\alpha ,`$
by applying $`(ii)`$ with $`m=0`$ and multiplying by $`\frac{5}{2}`$,
$`p_9(x,y)`$ $`=`$ $`2x^915x^7y^26x^6\beta +27x^5y^4+4x^3(\beta ^227y)+24xy^3+21x^4y^2\beta 24x\alpha ,`$
by applying $`(ii)`$ with $`m=1`$ and multiplying by $`6`$, and
$`p_{10}(x,y)`$ $`=`$ $`20x^{10}144x^8y^255x^7\beta +252x^6y^4+189x^5y^2\beta 35x^4(48y\beta ^2)`$
$`+420x^2(y^3\alpha )210,\text{ by applying }(ii)\text{ with }m=2\text{ and multiplying by }105\text{.}`$
We do not know whether Theorem 14 generates all the polynomials in $`𝒪`$.
###### Proof of Theorem 14.
It is useful to have the following two formulas, which follow from multiplying (13) by $`\overline{u}`$ and separating real and imaginary parts:
(17)
$$\mathrm{Im}[u_x(x+iy)\overline{u}(x+iy)]_x=(x^33xy^2\beta )|u(x+iy)|^2$$
and
(18)
$$\mathrm{Re}[u_x(x+iy)\overline{u}(x+iy)]_x=|u_x|^2+(3x^2y+y^3\alpha )|u|^2.$$
At the end of the proof we will justify the fact that we can differentiate through the integrals that follow.
$`(i)`$ This is clear by integrating (17), using the zero boundary conditions in the left- and right-hand Stokes regions.
$`(ii)`$ Suppose $`m`$ is a nonnegative integer. Then
$`{\displaystyle \frac{d}{dy}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}x^m|u|^2𝑑x={\displaystyle _{\mathrm{}}^{\mathrm{}}}x^m{\displaystyle \frac{}{y}}|u|^2𝑑x`$ $`=`$ $`2{\displaystyle _{\mathrm{}}^{\mathrm{}}}x^m\mathrm{Im}(u_x\overline{u})𝑑x`$
$`=`$ $`{\displaystyle \frac{2}{m+1}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}x^{m+1}\mathrm{Im}(u_x\overline{u})_x𝑑x,`$
where the last step is by integration by parts.
So using (17), we have that
$$\frac{d}{dy}_{\mathrm{}}^{\mathrm{}}x^m|u|^2𝑑x=\frac{2}{m+1}_{\mathrm{}}^{\mathrm{}}x^{m+1}(x^33xy^2\beta )|u|^2𝑑x.$$
Hence,
$`{\displaystyle \frac{d^2}{dy^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}x^m|u(x+iy)|^2𝑑x`$
$`=`$ $`{\displaystyle \frac{2}{m+1}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}[6x^{m+2}y|u|^22(x^{m+4}3x^{m+2}y^2\beta x^{m+1})\mathrm{Im}(u_x\overline{u})]𝑑x.`$
Then again using the integration by parts and (17), we have that this equals
(20)
$$\frac{2}{m+1}_{\mathrm{}}^{\mathrm{}}[6x^{m+2}y+2(\frac{x^{m+5}}{m+5}3y^2\frac{x^{m+3}}{m+3}\beta \frac{x^{m+2}}{m+2})(x^33xy^2\beta )]|u|^2𝑑x.$$
Also, we differentiate (4) without applying integration by parts:
(21) $`{\displaystyle \frac{d^2}{dy^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}x^m|u|^2𝑑x`$ $`=`$ $`2{\displaystyle _{\mathrm{}}^{\mathrm{}}}x^m[\mathrm{Re}(u_x\overline{u})_x2|u_x|^2]𝑑x`$
(22) $`=`$ $`2{\displaystyle _{\mathrm{}}^{\mathrm{}}}x^m[(3x^2y+y^3\alpha )|u|^2|u_x|^2]𝑑x\text{by (}\text{18}\text{)}.`$
Also, applying integration by parts twice to the right-hand side of (21), we have that (21) equals
(23)
$$m(m1)_{\mathrm{}}^{\mathrm{}}x^{m2}|u|^2𝑑x+4_{\mathrm{}}^{\mathrm{}}x^m|u_x|^2𝑑x.$$
By equating (22) and (23), we get
(24)
$$_{\mathrm{}}^{\mathrm{}}x^m|u_x|^2𝑑x=\frac{m(m1)}{2}_{\mathrm{}}^{\mathrm{}}x^{m2}|u|^2𝑑x_{\mathrm{}}^{\mathrm{}}x^m(3x^2y+y^3\alpha )|u|^2𝑑x.$$
Hence equating (20) and (22) and substituting (24) give $`(ii)`$.
$`(iii)`$ Suppose that $`_{\mathrm{}}^{\mathrm{}}p(x,y)|u|^2𝑑x=0\text{ for all }y`$. Then
$`0`$ $`=`$ $`{\displaystyle \frac{d}{dy}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}p(x,y)|u|^2𝑑x`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}[p_y(x,y)|u|^22p(x,y)\mathrm{Im}(u_x\overline{u})]𝑑x`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[p_y(x,y)|u|^2+2\left\{{\displaystyle _0^x}p(t,y)𝑑t\right\}\mathrm{Im}(u_x\overline{u})_x\right]𝑑x,`$
by integration by parts. This with (17) gives $`(iii)`$.
$`(iv)`$ Suppose $`_{\mathrm{}}^{\mathrm{}}p(x,y)|u|^2𝑑x=0\text{ for all }y`$.
Then we differentiate through (4) with respect to $`y`$ again to get
$`0`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[p_{yy}(x,y)|u|^24p_y(x,y)\mathrm{Im}(u_x\overline{u})2p(x,y)\mathrm{Im}(iu_{xx}\overline{u}i|u_x|^2)\right]𝑑x`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[p_{yy}(x,y)|u|^2+4\left\{{\displaystyle _0^x}p_y(t,y)𝑑t\right\}\mathrm{Im}(u_x\overline{u})_x2p(x,y)(\mathrm{Re}(u_x\overline{u})_x2|u_x|^2)\right]𝑑x`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[p_{yy}(x,y)+4\left\{{\displaystyle _0^x}p_y(t,y)𝑑t\right\}(x^33xy^2\beta )\right]|u|^2𝑑x`$
$`2{\displaystyle _{\mathrm{}}^{\mathrm{}}}p(x,y)\left[\mathrm{Re}(u_x\overline{u})_x+2(3x^2y+y^3\alpha )|u|^2\right]𝑑x,\text{by (}\text{17}\text{) and (}\text{18}\text{)}`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[p_{yy}(x,y)+4\left\{{\displaystyle _0^x}p_y(t,y)𝑑t\right\}(x^33xy^2\beta )4p(x,y)(3x^2y+y^3\alpha )\right]|u|^2𝑑x`$
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}p_x(x,y)2\mathrm{R}\mathrm{e}(u_x\overline{u})𝑑x.`$
But $`(y^3\alpha )_{\mathrm{}}^{\mathrm{}}p|u|^2𝑑x=0`$, and so applying integration by parts again to the last term gives $`(iv)`$.
Now to complete the proof we need to show that we can differentiate through the above integrals, which reduces to showing that
$$\frac{d}{dy}_{\mathrm{}}^{\mathrm{}}x^m|u|^2𝑑x=_{\mathrm{}}^{\mathrm{}}x^m\frac{}{y}|u|^2𝑑x,\text{for each }m0.$$
So we estimate the following:
$`\left|{\displaystyle _{\mathrm{}}^{\mathrm{}}}x^m[{\displaystyle \frac{1}{h}}\{|u(x+i(y+h))|^2|u(x+iy)|^2\}{\displaystyle \frac{}{y}}|u(x+iy)|^2]𝑑x\right|`$
$``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}|x|^m|{\displaystyle \frac{1}{h}}\{|u(x+i(y+h))|^2|u(x+iy)|^2\}{\displaystyle \frac{}{y}}|u(x+iy)|^2|dx`$
$``$ $`|h|{\displaystyle _{\mathrm{}}^{\mathrm{}}}|x|^m\left|{\displaystyle \frac{^2(|u|^2)}{y^2}}(x+i\xi (x))\right|dx,\text{by the Mean Value Theorem,}`$
where $`|\xi (x)y||h|.`$
So then it is enough to show that for each $`M>0`$, there exist $`C_1>0`$ and $`C_2>0`$ such that
(26)
$$|u(x+iy)|+|u^{}(x+iy)|+|u^{\prime \prime }(x+iy)|\frac{C_1}{\mathrm{exp}[|x|^{C_2}]}\text{for all}|y|M.$$
Proof of (26) is as follows:
Suppose that $`z=re^{i\theta }`$ with $`|\theta |\frac{\pi }{20}`$. Then $`z`$ is in the decaying Stokes regions (see Figure 2). By the asymptotic expression (5) we get that for some $`C>0`$, $`|u(re^{i\theta })|\mathrm{exp}[|r|^C]`$ for all $`|\theta |\frac{\pi }{20}`$ and large $`|r|`$, say $`|r|R`$. Also choose $`RM+1`$.
Now the Cauchy integral formula says that
$$u^{(k)}(z)=\frac{k!}{2\pi i}_{|\zeta z|=1}\frac{u(\zeta )}{(\zeta z)^{k+1}}𝑑\zeta .$$
So then
$$|u^{(k)}(z)|\frac{k!}{2\pi }_{|\zeta z|=1}|u(\zeta )||d\zeta |k!\mathrm{max}\{|u(\zeta )|:|\zeta z|=1\}k!\mathrm{exp}[(|z|1)^C]$$
where the last inequality holds if $`\{\zeta :|\zeta z|=1\}\{\zeta :\zeta =\rho e^{i\varphi },|\rho |R,|\varphi |\frac{\pi }{20}\}`$.
Choose $`0<C_2<C`$. Then $`\mathrm{exp}[(|z|1)^C]\mathrm{exp}[|z|^{C_2}]`$ if $`|z|R_2`$ for large $`R_2R`$.
The region where $`|z|R_2`$ and $`\{\zeta :|\zeta z|=1\}\{\zeta :\zeta =\rho e^{i\varphi },|\rho |R,|\varphi |\frac{\pi }{20}\}`$ covers all of $`|y|M`$ but a bounded region.
Since the minimum of $`\mathrm{exp}[|z|^{C_2}]`$ in this bounded region is strictly positive, we can find a large $`C_1>0`$ so that the left-hand side of (26) is bounded by $`C_1\mathrm{exp}[|z|^{C_2}]`$. Thus (26) holds since $`|x||z|`$. And this completes the proof. ∎
###### Corollary 15.
Let $`u(z)`$ be an eigenfunction of (13). Then $`_{\mathrm{}}^{\mathrm{}}|u(x+iy)|^2𝑑x`$ is a convex function.
###### Proof.
This is a consequence of (23) with $`m=0`$, or it can be proved using the subharmonicity of $`|u|^2.`$
## 5. Conclusions
Using simple path integrations, we were able to prove that eigenvalues of (1) lie in the sector $`|\mathrm{arg}\lambda |\frac{\pi }{2n+3}`$ and we extended the result for some more general Hamiltonians. Also we provide zero-free regions of eigenfunctions and their first derivatives, for the potential $`(ix)^3`$. Then finally we have the set $`𝒪`$ of polynomials $`p(x,y)`$ which are orthogonal to $`|u|^2`$ in the sense that $`_{\mathrm{}}^{\mathrm{}}p(x,y)|u|^2𝑑x=0`$ for all $`y`$.
In a recent communication with Mezincescu, he pointed out that for the potential $`(ix)^3`$ if $`\mathrm{Im}\lambda =\beta 0`$, combining $`\frac{|\beta |}{\alpha }\mathrm{tan}\frac{\pi }{5}`$ with the equation ($`23`$) in gives $`|\lambda |>(\frac{2}{5}\mathrm{cos}\frac{\pi }{5})10^53\times 10^4.`$ So if any non-real eigenvalues exist, they are very large.
In this paper we consider only polynomial potentials with odd degrees. However, a number of other authors have worked on even degree potentials, particularly quartic and sextic polynomial potentials. Our techniques in proving Theorems 2 and 3 can be used to get information on eigenvalues for even degree potentials if both ends of a line passing through the origin stay in decaying Stokes regions.
Obvious open problems are to narrow the eigenvalue sectors closer to the positive real axis, and finally to prove that the eigenvalues are real. Since some $`𝒫𝒯`$-symmetric non-Hermitian Hamiltonians do not have all real eigenvalues, one might further want to classify $`𝒫𝒯`$-symmetric non-Hermitian Hamiltonians which do have positive real eigenvalues.
### Acknowledgments
The author was partially supported by NSF grant number DMS-9970228. The author appreciates Gary Gundersen’s help at the initial stage and Carl Bender’s comments later on, and thanks Richard S. Laugesen for encouragement, invaluable suggestions and discussions throughout the work.
email contact: kcshin@math.uiuc.edu
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# Absence of Bose-Einstein condensation with a uniform field and contact interaction
## Abstract
The behavior of a Bose-Einstein ideal gas of particles in a three dimensional space in the presence of a uniform field, such as gravity, and of contact interaction, describing the presence of one impurity, is investigated. It is shown that Bose-Einstein condensation can not occur.
The recent experiments on Bose-Einstein Condensation (BEC) in atomic vapor magnetically trapped, have stimulated a new interest in the theoretical study of Bose gases . In this Letter we propose a theoretical model, that could be relevant to simulate what happens in the presence of some particular external fields such as a gravitational uniform field and a point-like (or at “zero range”, or $`\delta `$-like) interaction, which describes the presence of impurities. We start from the classical one-particle Hamiltonian of the problem, that is
$$H=\frac{𝐩^2}{2m}+mgx_3,$$
(1)
where $`m`$ and $`𝐩=(p_1,p_2,p_3)`$ are the mass and the momentum of the particle respectively. The contact interaction, i.e. the presence of the impurity, will be included in the formalism after quantization, which can be done only under the specification of particular boundary conditions on the wave function satisfying the Schrödinger equation . After the introduction of dimensionless cylindrical coordinates,
$$x_1=\frac{r}{\kappa }\mathrm{cos}\varphi ,x_2=\frac{r}{\kappa }\mathrm{sin}\varphi ,x_3=\frac{z}{\kappa }$$
(2)
$$0<\varphi <2\pi ,\mathrm{\hspace{0.33em}0}<r<\mathrm{},\mathrm{}<z<\mathrm{}$$
(3)
with the quantum gravitational length and energy defined by
$$\lambda _g\kappa ^1=\left(\frac{\mathrm{}^2}{2m^2g}\right)^{1/3},E_g=\frac{\mathrm{}^2\kappa ^2}{2m},$$
(4)
the eigenvalues equation $`H\mathrm{\Psi }_E=E\mathrm{\Psi }_E`$ gives us the following general form of the eigenfunctions in terms of Bessel’s, Neumann’s and Airy’s functions :
$`\mathrm{\Psi }_{l,\lambda ,ϵ}(r,z,\varphi )=Z(\eta )e^{il\varphi }R_l(\lambda r)/\sqrt{2\pi },`$ (5)
$`\eta =zϵ+\lambda ^2,`$ (6)
$`R_l(\lambda r)=A_{R,l}J_l(\lambda r)+B_{R,l}N_l(\lambda r),`$ (7)
$`Z(\eta )=A_ZAi(\eta )+B_ZBi(\eta ),`$ (8)
where $`A_{R,l},B_{R,l},A_Z,B_Z`$ are normalization constants to be further determined, whilst the dimensionless quantum numbers are defined to be
$$ϵ=\frac{E}{E_g}𝐑,l𝐙,\lambda =\sqrt{\frac{p_x^2+p_y^2}{\mathrm{}^2\kappa ^2}}0.$$
(9)
In order to evaluate the normalization constants we have to impose the usual orthonormality relations
$`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑zZ(z,ϵ,\lambda )Z(z,ϵ^{},\lambda )=\delta (ϵϵ^{})`$ (10)
$`{\displaystyle _0^{\mathrm{}}}𝑑rrR_l(\lambda r)R_l(\lambda ^{}r)=\delta (\lambda \lambda ^{})`$ (11)
Concerning eq. (11), it is necessary to separate the cases with $`l0`$ and $`l=0`$. This distinction is the main point of the discussion. As a matter of fact, we have to carefully take into account that the particles are in a uniform field and that they feel, from a quantum mechanical point of view, another extremely localized interaction (point-like or $`\delta `$-like interaction) usually called contact interaction describing the impurity. We know that $`\delta `$-like potentials on the plane are not mathematically well defined in quantum mechanics. The only correct quantum mechanical formalism to describe them is in terms of the self-adjoint extensions of the Hamiltonian operator which, in the present case, is given by the differential operator of eq. (1). Actually, that differential operator is only symmetric, until we specify its domain. These self-adjoint extensions of the Hamiltonian are non-trivial in the case $`l=0`$, when both the regular and the irregular wave functions at $`r=0`$ are allowed, in contrast with the case $`l0`$, in which singularities at $`r=0`$ can not be accepted. It is indeed possible to see from the behavior of the Neumann’s function for large argument that $`B_{R,l}=0`$, for the wave-function has to be locally square integrable. Therefore we readily get
$$A_{R,l}(\lambda )=\sqrt{\lambda },l0.$$
(12)
In the case $`l=0`$ the presence of the irregular part of the wave function is allowed since it turns out to be locally square integrable. A straightforward calculation leads to the result that the normalization condition (11) is true if and only if
$$\pi \mathrm{tan}[\pi \nu (\alpha _0,\lambda )]=\mathrm{ln}(\alpha _0/\lambda ^2),\lambda 0,$$
(13)
where $`0\alpha _0`$ is an arbitrary constant parameter. Moreover, taking eq. (13) into account we definitely obtain
$`A_{R,0}(\lambda )=\sqrt{\lambda }\mathrm{sin}(\pi \nu ),B_{R,0}(\lambda )=\sqrt{\lambda }\mathrm{cos}(\pi \nu ).`$ (14)
Now it is very important to understand the role of $`\alpha _0`$. As we have already emphasized, we are dealing with a quantum-mechanical point-like interaction, which can be correctly handled through the formalism of the self-adjoint extensions of the Hamiltonian operator . The latter formalism entails that in order to describe the dynamics a whole one-parameter continuous family of different quantum mechanical Hamiltonians has to be considered. The different self-adjoint Hamiltonian operators are naturally labelled by the parameter $`\alpha _0`$. When $`\alpha _0=0`$ (Friedrichs’ limit), the contact interaction is turned off and, from the mathematical point of view, it means that the irregular part of the wave-function disappears. Actually
$$\mathrm{tan}(\pi \nu )=\mathrm{}\nu =1/2,$$
(15)
i.e., the same result as in the case $`l0`$. Let us finally evaluate the constants $`A_Z`$,$`B_Z`$. Substituting eq. (8) into eq. (10), and noting that from the asymptotic behavior of the Airy’s function $`Bi(z)`$ we have to set $`B_Z=0`$ \- in order to keep square summability on the positive half-line $`z>0`$ \- we finally get
$$(A_Z)^2_{\mathrm{}}^{\mathrm{}}𝑑zAi(z,ϵ,\lambda )Ai(z,ϵ^{},\lambda )=\delta (ϵϵ^{}).$$
(16)
Using the integral representation for the Airy’s function we find
$$A_Z=1/\sqrt{2\pi },B_Z=0,$$
(17)
the wave functions of the eigenstates with vanishing angular momentum and with a continuous degeneracy $`\lambda 0`$ due to the transverse momentum becomes
$$\mathrm{\Psi }_{0,\lambda ,ϵ}(r,z,\varphi )=\frac{\sqrt{\lambda }Ai(\eta )}{2\pi \mathrm{sec}(\pi \nu )}\left[N_0(\lambda r)\mathrm{tan}(\pi \nu )J_0(\lambda r)\right],$$
(18)
whereas the eigenfunctions of the states with non-vanishing angular momentum read
$$\mathrm{\Psi }_{l,\lambda ,ϵ}(r,z,\varphi )=\frac{\sqrt{\lambda }}{2\pi }J_l(\lambda r)Ai(\eta )e^{il\varphi }.$$
(19)
From the previously obtained values for the normalization constants, it follows that the improper eigenfunctions are normalized according to
$$\mathrm{\Psi }_{l,\lambda ,ϵ}|\mathrm{\Psi }_{l^{},\lambda ^{},ϵ^{}}=\delta _{l,l^{}}\delta (\lambda \lambda ^{})\delta (ϵϵ^{}).$$
(20)
Now we are ready to discuss the existence of a characteristic state which faithfully and uniquely specifies any given self-adjoint extension, within the above mentioned one-parameter continuous family of the quantum Hamiltonian operators that are allowed according to the general principles of quantum mechanics. To this purpose, let us start again from the Schrödinger stationary equation, after the replacement $`\lambda i\widehat{\lambda }`$. Then we obtain
$`{\displaystyle \frac{r}{R_l(r)}}{\displaystyle \frac{}{r}}\left(r{\displaystyle \frac{}{r}}\right)R_l(r)\widehat{\lambda }^2r^2l^2`$ $`=`$ $`0,`$ (21)
$`{\displaystyle \frac{1}{Z(z)}}{\displaystyle \frac{^2}{z^2}}Z(z)z+ϵ+\widehat{\lambda }^2`$ $`=`$ $`0.`$ (22)
The first equation is exactly the same that we had previously discussed in the case of states with a continuous degeneracy in the transverse momentum. Our interest relies only in the subspace of vanishing angular momentum because it is only in this subspace that non-trivial self-adjoint extensions of the Hamiltonian actually occur. For $`l0`$ the solution of eq. (21) is not acceptable because it does not belong to $`L^2(𝐑^3)`$. Taking $`l=0`$ and evaluating the wave function of the state that characterizes the extension we get
$$\widehat{\mathrm{\Psi }}_{0,ϵ}(r,z)=\frac{\widehat{\lambda }}{\sqrt{2\pi ^2}}K_0(\widehat{\lambda }r)Ai(z\widehat{ϵ}\widehat{\lambda }^2).$$
(23)
This state has to be orthogonal to any of the improper states with a continuous degeneracy of the transverse momentum $`\lambda 0`$ \- see eq. (18). We expect to find that the orthogonality takes place when the parameter $`\widehat{\lambda }^2`$ of the state characterizing the self-adjoint extensions is exactly equal to the above introduced parameter $`\alpha _0`$. Actually we require
$$\mathrm{\Psi }_{0,\lambda ,ϵ}|\widehat{\mathrm{\Psi }}_{0,ϵ^{}}=0,ϵϵ^{},\lambda 0.$$
(24)
Explicit evaluation leads to the result
$$\pi \mathrm{tan}(\pi \nu )=\mathrm{ln}(\widehat{\lambda }^2/\lambda ^2),$$
(25)
which exactly corresponds to the definition of $`\alpha _0`$ in eq. (13). Consequently, orthogonality occurs iff $`\widehat{\lambda }^2=\alpha _0`$, and from now on we can label the state that uniquely characterizes the extension with the parameter $`\alpha _0`$.
The calculation of the one-particle Heat-Kernel is the first step to get the partition function of the system. It is also interesting to examine the limiting cases in which the background fields are switched off. As a matter of fact, we have to recover the Heat-Kernel of the non-relativistic free particle when both of the external fields are turned off, i.e., $`g0`$ and $`\alpha _00`$ . To get the general form of the Heat-Kernel, we need the spectral decomposition of the quantum Hamiltonian operator. To this purpose, it is necessary to keep in mind that we have to deal both with the states with a continuous degeneracy, labelled by the eigenvalues of the transverse momentum, and with the single state characterizing the self-adjoint extension. Furthermore, it is convenient from now on to measure energies in units of quantum gravitational energy scale $`E_g`$. As a consequence, we define the dimensionless Hamiltonian operators as
$$E_g^1H(\alpha _0)\mathrm{H}(\alpha _0)=\stackrel{ˇ}{\mathrm{H}}[\nu (\alpha _0)]+\widehat{\mathrm{H}}(\alpha _0),$$
(26)
where $`\stackrel{ˇ}{\mathrm{H}}`$ is the part of the spectral decomposition of the Hamiltonian that involves the states with a continuous degeneracy of the transverse momentum, and $`\widehat{\mathrm{H}}`$ is the one that involves the characteristic state of the self-adjoint extension. The explicit form of the spectral decompositions can be written as
$`\stackrel{ˇ}{\mathrm{H}}[\nu (\alpha _0)]={\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑ϵϵ{\displaystyle _0^{\mathrm{}}}𝑑\lambda {\displaystyle \underset{l𝐙}{}}|ϵ,\lambda ,l><ϵ,\lambda ,l|,`$ (27)
$`\widehat{\mathrm{H}}(\alpha _0)={\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑ϵϵϵ,\alpha _0,0><ϵ,\alpha _0,0.`$ (28)
Let us come back to the evaluation of the diagonal kernel. For a matter of simplicity, we shall calculate separately the part of the kernel corresponding to the continuous degeneracy in the transverse momentum, and the other one which characterizes the extension: namely,
$`G(\alpha _0,\gamma ;𝐫)`$ $`=`$ $`\stackrel{ˇ}{G}(\alpha _0,\gamma ;𝐫)+\widehat{G}(\alpha _0,\gamma ;𝐫),`$ (29)
$`\stackrel{ˇ}{G}(\alpha _0,\gamma ;𝐫)`$ $`=`$ $`e^{\gamma z}𝐫\mathrm{exp}\{\gamma \stackrel{ˇ}{\mathrm{H}}[\nu (\alpha _0)]\}𝐫,`$ (30)
$`\widehat{G}(\alpha _0,\gamma ;𝐫)`$ $`=`$ $`e^{\gamma z}𝐫\mathrm{exp}\{\gamma \widehat{\mathrm{H}}(\alpha _0)\}𝐫,`$ (31)
where $`𝐫=(r,\varphi ,z)`$ is the position in the dimensionless coordinate space associated to the system whereas $`\gamma =\beta E_g`$ and $`\beta 1/k_BT`$, with $`k_B`$ the Boltzmann’s constant. The reason for the factor $`\mathrm{exp}\{\gamma z\}`$ in eq.s (30-31) is that one has to recover translation invariance in the absence of contact interaction, i.e., in the limit $`\alpha _00`$. The diagonal kernel $`\stackrel{ˇ}{G}`$ is that one with a continuous degeneracy in the transverse momentum. From eq.s (30) and (27) we can write
$$\stackrel{ˇ}{G}(\alpha _0,\gamma ;𝐫)=_{\mathrm{}}^{\mathrm{}}𝑑ϵe^{\gamma (zϵ)}_0^{\mathrm{}}𝑑\lambda \underset{l𝐙}{}\left|\mathrm{\Psi }_{l,\lambda ,ϵ}(r,z,\varphi )\right|^2$$
(32)
Substituting the wave functions (18)-(19) into the expression (32) we obtain
$`\stackrel{ˇ}{G}(\alpha _0,\gamma ;𝐫)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta e^{\gamma \eta }Ai^2(\eta ){\displaystyle _0^{\mathrm{}}}𝑑\lambda \lambda e^{\gamma \lambda ^2}`$ (33)
$`\times `$ $`\{1\mathrm{cos}^2(\pi \nu )[J_0^2(\lambda r)N_0^2(\lambda r)]`$ (34)
$`+`$ $`\mathrm{sin}(2\pi \nu )J_0(\lambda r)N_0(\lambda r)\}.`$ (35)
The integral over the Airy’s function can be further elaborated and eventually set in the more suitable form where
$$I(\gamma )_{\mathrm{}}^+\mathrm{}𝑑\eta e^{\gamma \eta }Ai^2(\eta )=I_1+I_2+I_3,$$
(36)
with
$`I_1`$ $`={\displaystyle \frac{1}{18}}\sqrt{{\displaystyle \frac{3}{\pi }}}\left({\displaystyle \frac{3}{2}}\right)^{\frac{1}{3}}E_3[{\displaystyle \frac{\gamma }{3}}\left({\displaystyle \frac{3}{2}}\right)^{\frac{2}{3}},{\displaystyle \frac{7}{6}}],`$ (37)
$`I_2`$ $`={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\xi {\displaystyle \frac{\gamma }{\sqrt{\xi }(\gamma ^2+\xi ^2)}}\mathrm{cos}{\displaystyle \frac{\xi ^2}{12}},`$ (38)
$`I_3`$ $`={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\xi {\displaystyle \frac{(\gamma )}{\sqrt{\xi }(\gamma ^2+\xi ^2)}}\mathrm{sin}{\displaystyle \frac{\xi ^2}{12}},`$ (39)
$`E_\rho (z,\mu )`$ being the Mittag Leffer function .
Finally, taking the above integral representation into account, we can rewrite the Heat-Kernel $`\stackrel{ˇ}{G}`$ in the simpler form
$`\stackrel{ˇ}{G}(\alpha _0,\gamma ;r)=I(\gamma )/(8\pi ^2\gamma )\{1{\displaystyle _0^{\mathrm{}}}dte^t\mathrm{cos}^2(\pi \nu )`$ (40)
$`\times [J_0^2\left(s\sqrt{t}\right)N_0^2\left(s\sqrt{t}\right)]\mathrm{sin}(2\pi \nu )J_0\left(s\sqrt{t}\right)N_0\left(s\sqrt{t}\right)\},`$ (41)
$`s={\displaystyle \frac{\sqrt{4\pi (x_1^2+x_2^2)}}{\lambda _T}},\pi \mathrm{tan}(\pi \nu )=\mathrm{ln}(\alpha _0\gamma /\mathrm{ln}t),`$ (42)
where $`\lambda _T`$ is the thermal wave length.
The evaluation of the diagonal kernel $`\widehat{G}`$ can be done in close analogy with the previous case. From eq. (23) with $`\widehat{\lambda }^2`$ replaced by $`\alpha _0`$ and from eq.s (28) and (31) we get
$$\widehat{G}(\alpha _0,\gamma ;r)=\frac{\alpha _0}{2\pi ^2}e^{\gamma \alpha _0}K_0^2(r\sqrt{\alpha _0})I(\gamma ).$$
(43)
It is very interesting now to examine the behavior of the kernel in the limiting cases in which the background fields are switched off. Particularly, we look at the behavior of the diagonal Heat-Kernel when $`g0`$, i.e., when the uniform field is off, and when $`\alpha _00`$, that means without point-like impurity. Finally we will examine the case when both the external fields are turned off and we show that there is commutativity with respect to order of the operations. We do emphasize that, in order to evaluate those limits, it is necessary to pass from the dimensionless Heat-Kernel, to the corresponding one in physical units: namely
$$\stackrel{~}{G}(\alpha _0,g,T;𝐱)=\kappa ^3G(\alpha _0,\gamma ;𝐫)=\left(\frac{\sqrt{4\pi \gamma }}{\lambda _T}\right)^3G(\alpha _0,\gamma ;r),$$
(44)
where the dimensional quantities are labeled with a tilde. From the behavior for small $`\gamma `$, of the integral $`I(\gamma )\sqrt{(\pi /\gamma )}`$ \- see eq.s (36-39) - it immediately follows that
$$\stackrel{~}{G}(g0)=\lambda _T^3.$$
(45)
which shows that the above discussed specific form of the contact interaction does completely depend upon the presence of a non-vanishing uniform field. In order to perform the limit in which the point-like impurity is removed, it is important to keep in mind that $`\alpha _0`$ and $`\nu `$ are strictly related by eq. (13), so that $`\nu `$ goes to $`(1/2)`$ when $`\alpha _0`$ goes to zero. The explicit evaluation yields
$$\stackrel{~}{G}(\alpha _00)=\frac{I(\beta E_g)}{\lambda _T^3}\sqrt{\frac{\beta E_g}{\pi }},$$
(46)
which exhibits manifest translation invariance as it does. From eq.s (45) and (46), after switching off the remaining background fields we get the expected result
$$\underset{g0}{lim}\stackrel{~}{G}(\alpha _0,g,T;𝐱)=\lambda _T^3=\underset{g0}{lim}\underset{\alpha _00}{lim}\stackrel{~}{G}(\alpha _0,g,T;𝐱),$$
(47)
i.e., the Heat-Kernel of a free non-relativistic particle. Furthermore we have shown that the final result is independent from the order with which the external fields are switched off.
The partition function per unit volume of a single molecule is now to be evaluated in the thermodynamic limit. To this aim, let us first consider the case without contact interaction, i.e., $`\alpha _00`$. In this case, owing to translation invariance, we immediately get
$$Z_0(\beta )=_0^{\mathrm{}}𝑑E\rho (E)e^{\beta E}=\pi \left(\frac{2mE_g}{h^2}\right)^{3/2}\frac{I(\beta E_g)}{\beta E_g}.$$
(48)
The corresponding density $`\rho (E)`$ of the one-particle quantum states, i.e., the number of the one-particle quantum states per unit volume and within the energy interval $`E`$ and $`E+dE`$, $`E0`$, can be obtained as the inverse Laplace transform of the above expression (48). Explicit evaluation yields
$$\rho (E)=\pi \sqrt{E_g}\left(\frac{2m}{h^2}\right)^{3/2}\left[\frac{E}{E_g}Ai^2\left(\frac{E}{E_g}\right)+Ai_{}^{}{}_{}{}^{2}\left(\frac{E}{E_g}\right)\right].$$
(49)
As a consequence, we can write the average density of particles in the form
$$\frac{N_T}{V}n_T=z_0^{\mathrm{}}𝑑E\frac{\rho (E)e^{\beta E}}{1ze^{\beta E}}+\frac{z}{V(1z)}$$
(50)
where $`z=e^{\beta \mu }`$. From the above expression it immediately follows that there is no Bose-Einstein condensation in the presence of a uniform gravitational field, owing to the behavior
$$\rho (E=0)=\sqrt{E_g}\left(\frac{2m}{h^2}\right)^{3/2}\frac{3^{2/3}\pi }{[\mathrm{\Gamma }(1/3)]^2}.$$
(51)
Now, it is quite evident that we can rewrite eq. (50) in the suggestive form
$`n_T`$ $`=`$ $`{\displaystyle \frac{3^{2/3}\kappa }{2\pi [\mathrm{\Gamma }(1/3)]^2}}{\displaystyle \frac{g_1(z)}{\lambda _T^2}}+z{\displaystyle _0^{\mathrm{}}}𝑑E{\displaystyle \frac{\rho (E)\rho (0)}{1ze^{\beta E}}}e^{\beta E}`$ (52)
$`+`$ $`{\displaystyle \frac{z}{V(1z)}},`$ (53)
which shows that the lack of Bose-Einstein condensation in the presence of a uniform field in three spatial dimension just corresponds to the very same phenomenon for an ideal gas of free particles in two spatial dimensions. A little though readily drives to gather that also in the case $`\alpha _00`$ the same conclusion holds true, namely condensation does not occur, in contrast with the claim of ref. As a matter of fact, if we set
$`G(\alpha _0,\gamma ;r)`$ $``$ $`G(\alpha _0=0,\gamma ;r)+\mathrm{\Delta }G(\alpha _0,\gamma ;r)`$ (54)
$`=`$ $`{\displaystyle \frac{I(\gamma )}{8\pi ^2\gamma }}+\mathrm{\Delta }G(\alpha _0,\gamma ;r),`$ (55)
it can be proved that integration with respect to the radial dimensionless coordinates leads to the finite result
$$_0^{\mathrm{}}𝑑rr\mathrm{\Delta }G(\alpha _0,\gamma ;r)=_0^{\mathrm{}}𝑑t\frac{(\alpha _0\gamma )^t}{\mathrm{\Gamma }(t+1)}.$$
(56)
It immediately follows that in the thermodynamic limit the only relevant term is the first one in the RHS of eq. (54), all the rest disappearing in that limit, i.e. no Bose-Einstein condensation is possible in the presence of a uniform field and of a point-like impurity. It should be noticed that, from the very same results of ref. , the absence of Bose-Einstein condensation persists even if we describe the impurity by means of an Aharonov-Bohm vortex .
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# The Cometary Blue Compact Dwarf Galaxies Mkn 59 and Mkn 71:
## 1 Introduction
Despite the effort that has been devoted to the investigation of Blue Compact Dwarf Galaxies (BCDs), the origin of their recurrent starburst activity as well as its impact on their spectrophotometric and dynamical evolution is still poorly understood. A thorough investigation of these processes is required before questions pertaining to evolutionary connections among dwarf irregulars (dIs), dwarf ellipticals (dEs) and BCDs (Thuan thuan85 (1985), Davies & Phillipps davies88 (1988), Papaderos et al. 1996b; hereafter P96b , Patterson & Thuan patterson96 (1996), Salzer & Norton salzer98 (1998), Marlowe et al. marlowe97 (1997), marlowe99 (1999)) can be answered. Deep imaging of the low surface brightness (LSB) component of BCDs, discovered first by Loose & Thuan (1985; hereafter LT (85)), disclosed an evolved stellar population underlying the regions of active star formation (SF) (Loose & Thuan lt86 (1986), Kunth et al. kunth88 (1988), Papaderos et al. 1996a; hereafter P96a ). Therefore, the initial hypothesis that BCDs are “extragalactic H ii regions” lacking an older stellar population (Sargent & Searle sargent70 (1970)) had to be dismissed for the majority of these systems.
An exception to this finding is made by a tiny fraction ($``$ 1% ) of young galaxy candidates which recently have been identified among BCDs (Thuan et al. til97 (1997), Papaderos et al. papaderos98c (1998), Thuan & Izotov ti99 (1999)). There is increasing observational evidence that such systems, invariably found among the most metal–deficient ($`Z`$ 1/20 $`Z_{}`$, Thuan & Izotov ti99 (1999), Izotov & Thuan it99 (1999)) members of the BCD class, have started forming the bulk of their stellar component less than 100 Myr ago. This qualifies them as local counterparts of primeval galaxies, i.e. the first galactic building blocks which are thought to have formed at high redshifts.
It seems meaningful to search among the variety of dwarf galaxies for possible successors of young galaxies, i.e. objects in a more evolved stage with a substantial fraction of their masses still in gaseous form. An age sequence built therefrom would allow to study the evolution of nearby dwarf galaxies, which may at the same time give clues to early galaxy evolution processes at high redshift. This is of great interest for the assessment of the starburst–driven evolution of young galaxies, some of which may become the building blocks of larger systems and for the understanding of issues like the faint blue galaxy excess at intermediate redshifts (Babul & Rees babul92 (1992), Koo et al. koo97 (1997), Guzmán et al. guzman98 (1998)).
Loose & Thuan (1985) developed a morphological classification scheme for BCDs and found that $``$ 90% of their sample is made up of iE– and nE–objects. These show, respectively, an irregular or nuclear starburst component superposed on a smooth elliptical LSB host galaxy with optical colours pointing to ages of several Gyr (Loose & Thuan lt86 (1986), P96b , Telles & Terlevich telles97 (1997), Noeske noeske99 (1999)). In contrast, the few members of the rare class of i0 BCDs examined so far have been found to be young galaxy candidates (Thuan et al. 1997, Papaderos et al. papaderos98c (1998), Thuan & Izotov ti99 (1999), Izotov et al. izotov99 (1999)), as they apparently lack a smooth stellar underlying component that would have required previous episodes of SF to be built.
Only the class of iI BCDs, characterized by both irregular LSB components and starburst regions, has not yet been studied systematically. This paper focuses on the subset of iI,C BCDs (LT (85)) where “C” denotes a *cometary* appearance caused by a bright star–forming complex at one end of an elongated stellar galactic body (see Figs. 1 and 6).
In the notation of Salzer et al. (salzer89 (1989)) an iI,C BCD is to be included in the class of H ii hotspot galaxies; these systems seem to systematically differ from other subclasses of H ii galaxies with respect to the metallicities and excitation properties of their ionized gas.
While the main bodies of iI,C systems are reminiscent of low surface brightness dwarf irregulars (dIs) or magellanic irregulars (Kennicutt et al. kennicutt80 (1980), Dottori et al. dottori94 (1994), Wilcots et al. wilcots96 (1996)) with some widespread low–level SF, the bright off–center H ii regions are typical for BCDs with respect to their sizes, H$`\alpha `$ luminosities and electron temperatures. The detection of Super Star Clusters (SSCs) in these spots points to a very recent or still ongoing starburst episode (Kennicutt et al. kennicutt80 (1980), Dottori et al. dottori94 (1994), Barth et al. barth94 (1994)).
Optical surveys of BCDs (LT (85), Kunth et al. kunth88 (1988), Salzer et al. salzer89 (1989), P96a , Marlowe et al. marlowe97 (1997), Doublier et al. doublier97 (1997)) and dIs (e.g. Hunter et al. hunter98 (1998)) suggest that in these systems SF per unit area correlates with the local surface density of the underlying host galaxy. In this respect, the occurrence of a starburst with the amplitude observed in iI,C BCDs at one end of an elongated stellar LSB body appears puzzling and may be interpreted in at least two ways:
(i) iI,C BCDs are in fact dIs or magellanic irregulars caught during a brief stochastic enhancement of their otherwise moderate star formation rate per unit area. This would imply that both the structural properties of their stellar hosts and H I envelopes will be indistinguishable from those of other dIs.
(ii) In iI,C BCDs the conditions necessary for the ignition of a starburst are not fulfilled at the center but in the outskirts of their LSB components. This may be due to a conspiracy of intrinsic and external properties of an evolved BCD (kinematics of the gaseous component and Dark Matter distribution or external perturbation by a companion or intracluster gas) or alternatively be a common feature among less evolved systems.
Here we attempt to test such hypotheses by analysing optical data of two nearby iI,C BCDs, Mkn 59 (NGC 4861) and Mkn 71 (NGC 2366, hosting the H ii region complex NGC 2363). The distance to Mkn 71 was determined to $`D`$=3.44 Mpc through observations of Cepheids (Tolstoy et al. tolstoy95 (1995)). For Mkn 59 different models for infall and small–scale perturbations in the Virgo Cluster environment yield a distance between $``$ 10.7 Mpc (Heckman et al. heckman98 (1998)) and 17.8 Mpc (Tully tully88 (1988)). The principal results of this work do not depend on the exact distance to Mkn 59 (cf. Section 3.1). Throughout this paper, we adopt $`D`$=10.7 Mpc (Heckman et al. heckman98 (1998)).
Using broad– and narrow band images and long–slit spectra from own ground–based observations and archival HST data, we shall investigate the physical state and the chemical composition of the ISM and examine whether the iI,C BCDs under study differ from typical iE/nE BCDs with respect to the ages and structural properties of their LSB hosts. In Section 2, we describe our observations and data analysis. We present our results in Section 3, and discuss them in Section 4. In Section 5, we summarize our results and conclusions.
## 2 Data acquisition and processing
### 2.1 Ground–based imaging
#### 2.1.1 Observations and data reduction
Images were taken on March 7th – 10th 1997 at the 2.2m telescope of the German–Spanish Astronomical Center, Calar Alto, Spain. We used the Calar Alto Faint Object Spectrograph (CAFOS), equipped with a 2048$`\times `$2048 pixel SITe#1d CCD. The focal ratio of f/4.4 in the RC focus and the pixel size of 24 $`\mu `$m/pixel yield an instrumental scale of 0.53″/pixel and a usable field of view of $``$ 15′. At a gain ratio of 2.3 e<sup>-</sup> ADU<sup>-1</sup> the read–out noise was $`<`$ 3 counts (rms). Column (1) of the observing log (Table 1) contains the equatorial coordinates of the targets, cols. (2) and (3) give the filters and exposure times, respectively. Column (4) lists the night of the observing run, starting with night 1 from March 7th to March 8th 1997. The seeing and the mean sky surface brightness during each exposure are given in cols. (5) and (6). During each night, dark–, bias–, and flat–field exposures were taken, and the photometric standard field NGC 2419 (Christian et al. christian85 (1985)) was observed at different zenith angles. Using the ESO MIDAS<sup>4</sup><sup>4</sup>4Munich Image Data Analysis System, provided by the European Southern Observatory (ESO). software package, standard reduction and calibration steps were applied to the raw images. During nights 3 and 4 the conditions were photometric, resulting in calibration errors well below 0.05 mag. A poor atmospheric transparency during night 1 led to strong airmass–dependent terms for the $`B`$ band exposure of Mkn 71. Therefore the calibration was accomplished by comparing aperture photometry of bright isolated point sources on HST– and ground–based images. The photometric uncertainty of the latter $`B`$ band exposure was found to be $``$ 0.1 mag. All exposures for one object were aligned to each other using positions of point sources; the transformed images were generated by re–sampling the original image using a flux–conserving routine. Colour maps were derived after the resolution of different frames had been equalized by convolving the exposure with the better resolution with a normalized gaussian distribution of adequate width.
The continuum level in the H$`\alpha `$ images was inferred by scaling the $`R`$ band images by an empirically determined factor $`C`$ so that the fluxes of field stars were matched between the raw H$`\alpha `$– and the scaled $`R`$ exposures. In turn, the scaled $`R`$ band exposures were subtracted from the raw H$`\alpha `$ images, giving the emission flux $`F_{\mathrm{em}}`$ at each pixel. The latter approach is strictly correct only when the H$`\alpha `$ emission line is contributing a minor fraction of the line–of–sight $`R`$ band flux. In regions where H$`\alpha `$ emission makes a substantial fraction of the photons received in the $`R`$ band the empirical scaling may lead to a slight overestimation of the continuum flux. In these cases the correct H$`\alpha `$ emission $`F_{em}^c`$ is obtained from $`F_{em}`$ as:
$$F_{em}^c=\frac{1}{1CT_R(\mathrm{H}\alpha )}F_{em},$$
(1)
where $`C`$ is the empirical factor described above and $`T_R(\mathrm{H}\alpha )`$ the mean transmission of the $`R`$ band filter at the H$`\alpha `$ wavelength. The resulting narrow band frames were calibrated using continuum–subtracted H$`\alpha `$ exposures of the planetary nebula NGC 2392 for which aperture flux measurements are given by Kaler (kaler83 (1983)).
To derive H$`\alpha `$ equivalent width maps we subtracted first the corrected line emission frame $`F_{em}^c`$ (Eq.1) from the raw H$`\alpha `$ image to obtain a continuum image $`F_{cont}`$, which was then normalized to the continuum flux per 1Å wavelength interval by dividing it by the effective width of the H$`\alpha `$ filter. The emission line image $`F_{em}^c`$ was divided by the transmission of the H$`\alpha `$ filter at the H$`\alpha `$ wavelength. The resulting frame, containing the total H$`\alpha `$ emission flux, was divided by the 1Å– normalized continuum flux frame to obtain the H$`\alpha `$ equivalent width image.
The colour- and H$`\alpha `$ maps were corrected for interstellar extinction following Savage & Mathis (savage79 (1979)) and using the extinction coefficients $`C`$($`\mathrm{H}\beta `$) given in Section 2.3, which translate into $`E(BV)`$=0.06 mag and $`E(BV)`$=0.08 mag for Mkn 59 and Mkn 71, respectively.
#### 2.1.2 Surface photometry
Surface brightness profiles (SBPs) are the result of a transformation of a galaxy’s 2–dimensional flux pattern into a monotonically decreasing 1–dimensional intensity distribution. For a set of methods to compute the photometric radius $`R^{}(\mu )`$, i.e. the radius of a circle with an area equal to the one enclosed by the isophote at the surface brightness level $`\mu [\mathrm{mag}/\mathit{}\mathrm{}]`$, see e.g. Loose & Thuan (1986) and P96a . The method of summing up the area defined by pixels with fluxes exceeding a threshold $`I`$($`\mu `$) (method iii; P96a ) was considered most appropriate for the BCDs analyzed here as both their underlying and starburst components are of irregular morphology. Other techniques, such as isophote integration or ellipse fitting, though well suited for more regular BCDs, are hardly applicable to the optical images of Mkn 59 and Mkn 71, mainly due to the fact that their high surface brightness components split into many star-forming regions occupying a substantial fraction of the host galaxy’s surface.
For each surface brightness level $`\mu `$, the number $`N(\mu )`$ of all pixels within a polygonal aperture with a count level $`I(\mu )`$ was determined; the corresponding photometric radius is
$$R^{}(\mu )=\frac{1}{\pi }\sqrt{N(\mu )A_{\mathrm{pxl}}},$$
(2)
where $`A_{\mathrm{pxl}}`$ is the solid angle per pixel in $`\mathit{}\mathrm{}`$. Numerical simulations showed that for the technique described above Poisson photon noise may cause a systematic flattening of SBPs at a signal to noise $`S/N`$ level $`10`$. This frequently observed artificial profile flattening could be, however, satisfactorily corrected down to a $`(S/N)`$ level of $`1`$ by using the adaptive filtering algorithm implemented in MIDAS (see Richter et al. 1991 for a description). The uncertainties at each photometric radius were calculated taking into account the corresponding intensity, the sky noise of each image, and the number of pixels involved in its calculation (cf. P96a ). Colour profiles were computed by subtracting SBPs from each other; the surface brightness– and colour profiles displayed in Figure 3 are corrected for interstellar extinction (cf. Section 2.1.1).
Unlike stellar systems like globular clusters and giant ellipticals where the mass to light ($`M/L`$) ratio may be regarded nearly constant over the whole system, the integrated luminosity of BCDs originates from two distinct stellar populations with substantially different ages and $`M/L`$ ratios. Although the LSB component contributes on average one half of the total $`B`$ band luminosity of a BCD (P96b , Salzer & Norton 1998), it contains most of the stellar mass of the system. Therefore, its intensity distribution yields information on the inner gravitational potential within which the starburst ignites and evolves. In order to disentangle the light distribution of the older host galaxy from the one of the superimposed starburst component, we applied a simple 2–component decomposition scheme adjusted interactively to each profile, instead of iteratively fitting the full nonlinear 3–component scheme described in P96a . The intensity distribution of the LSB component is approximated by an exponential fitting law of the form
$$I_\mathrm{E}(R^{})=I_{\mathrm{E},0}\mathrm{exp}\left(\frac{R^{}}{\alpha }\right),$$
(3)
or equivalently
$$\mu (R^{})=\mu _{\mathrm{E},0}+1.086\left(\frac{R^{}}{\alpha }\right)$$
(4)
where $`\alpha `$ denotes its exponential scale length in arcsec and $`\mu `$ the surface brightness level in mag/$`\mathit{}\mathrm{}`$. Equation (4) was adjusted to each profile by applying an error weighted linear fit to the data, at radii sufficiently large to be free from the contamination of the starburst light, i.e. where the colour profiles become constant and H$`\alpha `$ emission vanishes. To check at which surface brightness levels gaseous emission becomes negligible, we overlaid the H$`\alpha `$ emission– and $`EW(`$H$`\alpha )`$ maps with broad band isophotes (cf. Figs. 1 and 2) . The fits, obtained at radii $`R^{}`$ 76$`\mathrm{}`$ for Mkn 59 and $`R^{}`$ 160$`\mathrm{}`$ for Mkn 71, yield the extrapolated central surface brightness $`\mu _{E,0}`$ and the exponential scale length $`\alpha `$. The radial surface brightness distribution of the starburst component can be computed from the residual luminosity in excess of the fit Eq. (4).
Table 2 summarizes the results of the profile decomposition. In the same way as described in P96a , cols. (3) and (4) give the central surface brightness and exponential scale length of the LSB component. Columns (5) and (7) list respectively the isophotal radii of the starburst component $`P_{25}`$ (practically the “plateau radius” used in P96a ), and of the stellar LSB host, $`E_{25}`$. Both radii were determined from extinction–corrected SBPs at a surface brightness level of 25 mag/$`\mathit{}\mathrm{}`$. Columns (6) and (8) contain the apparent magnitudes of the latter components determined within $`P_{25}`$ and $`E_{25}`$, respectively, and col. (9) the total apparent magnitude in each band as derived from integration of the corresponding SBP out to the last measured point. Column (10) lists the effective radius $`r_{eff}`$ and the radius $`r_{80}`$, enclosing 80% of the galaxy’s total flux.
### 2.2 HST images
#### 2.2.1 Data reduction
Next we shall discuss colour–magnitude diagrams (CMDs) of the H ii complex NGC 2363 and the main body of Mkn 71 derived from Hubble Space Telescope (HST) archival data. The following analysis is based on integrations of 2$`\times `$800 and 1$`\times `$700 sec in the F439W filter, 2$`\times `$800 sec in the F547M filter and 1$`\times `$900 plus 1$`\times `$600 sec in the F656N filter taken on January 8th, 1996 with the WFPC2 (PI: Drissen, Proposal ID 06096). All images were reduced through the standard pipeline as described in Holtzman et al. (1995a ). Exposures taken in the same filter were co–added and corrected for cosmic ray events using the STSDAS package and IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Astronomy Observatories, which is operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation.. We applied a charge transfer efficiency correction to the data as described in the HST Data Handbook (hst98 (1998)), performed a sky background subtraction and used the synthetic zero points and transformation coefficients given by Holtzman et al. (1995b ) to transform the measured fluxes to magnitudes in the Johnson $`UBVRI`$ system. For the H$`\alpha `$ narrow band filter, F656N, the continuum was subtracted as described in Section 2.1.1; the correction given by Eq. (1) was not necessary, as the H$`\alpha `$ line lies outside the wavelength range covered by the F547M filter. The calibration of the H$`\alpha `$ images was done following the prescriptions by Holtzman et al. (1995b ).
#### 2.2.2 Colour–magnitude diagrams
The H ii region NGC 2363 within Mkn 71 (Figure 1; right) was centered on the PC/WF1 chip (instrumental scale of 0$`\stackrel{}{.}`$046/pixel), allowing the resolution of compact groups of stars with a mean linear separation of $``$ 0.8 pc. We utilized the PC1 data to derive CMDs in order to further constrain the results of the spectral population synthesis analysis (Section 3.3) using the DAOPHOT II stellar photometry package under ESO MIDAS. A model for the point spread function (PSF), necessary for performing multiple–PSF fitting in crowded fields, was computed from several isolated bright stars in each frame. Because of the severe crowding within some regions of NGC 2363 as well as local background variations due to strong gaseous emission, we evaluated the photometric uncertainties and completeness limits of the resulting point-source identification files, taking into account the error growth towards fainter luminosities and the magnitude histograms (cf. Figure 5).
The measured fluxes were calibrated and transformed to the Johnson $`B`$ and $`V`$ filters as described in the previous section. The luminosity functions of the detected point sources (Figure 5, left panel) suggest a reasonable completeness for sources brighter than $``$ 26 mag in both $`B`$ and $`V`$. Faintwards of the latter magnitude cutoff the photometric errors become dominant (see Figure 4). Close to the detection limit, an additional error source is introduced by the small–scale variations of the background due to gaseous emission as well as strong noise peaks which possibly are misinterpreted as faint stars (cf. the excess counts at low luminosities in the upper left panel of Figure 5). We shall therefore consider the data reliable down to a limit of $``$ 25 mag in $`B`$ and $`V`$, close to the one Drissen et al. (drissen99 (1999)) obtain for the same data set. This cutoff is sufficient for our purposes; from a simple comparison of data points and isochrones it becomes evident that, given the photometric errors and poor time resolution of the $`(BV)`$ CMD, no firm conclusions on the SF history can be drawn for ages $``$ 50 Myr. A second CMD was derived in the same way for a region on the WF4 chip apparently free of any appreciable signatures of current SF (marked in Figure 1) to obtain information on the galaxy’s underlying stellar population separately. From an inspection of the right panels of Figure 5, and from the considerations referring to the CMD of NGC 2363 we estimate the confidence limit to $``$24 mag for both $`B`$ and $`V`$. Either CMD is shown in Figure 4 with typical photometric uncertainties for different magnitude intervals along with synthetic isochrones for a metallicity of $`Z`$=0.001 adopted from Bertelli et al. (bertelli94 (1994)).
### 2.3 Spectroscopic observations and data reduction
Spectrophotometric observations of Mkn 59 and Mkn 71 were obtained on March 17th and 18th 1994, using the Ritchey–Chrétien spectrograph of the KPNO 4m Mayall telescope. We used a 2″$`\times `$300″ slit with the KPC-10A grating (316 lines mm<sup>-1</sup>) in first order, with a GG 385 order separation filter which cuts off all second order contamination for wavelengths blueward of 7400Å. This instrumental setup gave a spatial scale along the slit of 0$`\stackrel{}{.}`$69 pixel<sup>-1</sup>, a scale perpendicular to the slit of 2.7 Å pixel<sup>-1</sup>, a spectral range of 3500–7500 Å, and a spectral resolution of $``$ 5 Å. The seeing was 1$`\stackrel{}{.}`$5 FWHM. Mkn 59 was observed at two slit positions (cf. Figure 6), both centered at the brightest star–forming region with P.A. = 59 and P.A. = 15 (close to the direction of the major axis). The total exposure time for the first orientation of the slit was 70 minutes and was broken up into four subexposures. At the second slit orientation a single 15 minutes exposure was taken.
Mkn 71 was observed at a slit orientation of P.A. = 77 (see Figure 1), centered on the brightest star–forming region (NGC 2363 I). The total exposure time was 29 minutes, broken up into five subexposures. Three Kitt Peak IRS spectroscopic standard stars were observed during each night for flux calibration. The airmasses during all observations were $``$ 1.25, therefore a correction for atmospheric dispersion was not necessary. Spectra of He-Ne-Ar comparison lamps were obtained before and after each observation to calibrate the wavelength scale.
The data reduction was done with the IRAF software package. The two–dimensional spectra were bias–subtracted and flat–field corrected. Then the IDENTIFY, REIDENTIFY, FITCOORD, TRANSFORM, BACKGROUND and CALIBRATE routines were used to perform the wavelength calibration, correction for distortion and tilt, night sky subtraction and flux calibration for each frame. One–dimensional spectra were extracted from the flux–calibrated two–dimensional spectra using the IRAF task APALL. For Mkn 59 we chose ten 7″$`\times `$2″ regions along the major axis of the galaxy (i.e the slit at P.A. = 15), the locations of which are indicated in Figure 6. Six of these spectra are shown in Figure 7. From the slit at P.A. = 59 the subspectra were extracted at each pixel row along the slit, so that the resulting areas were 0$`\stackrel{}{.}`$69$`\times `$2″ each. For Mkn 71 (see Figure 1), the subspectra from regions 1 and 2 were extracted from 7″$`\times `$2″regions, while for regions 3 and 4 areas of 14″$`\times `$2″ and 28″$`\times `$2″ were used to obtain a sufficiently high $`S/N`$. The extracted spectra are shown in Figure 8.
The emission line intensities were measured utilizing a Gaussian profile fitting. All spectra were corrected for interstellar extinction, where the extinction coefficient $`C`$(H$`\beta `$) was derived from the hydrogen Balmer decrement using the equations given in Izotov et al. (itl94 (1994)) and the theoretical hydrogen emission line flux ratios from Brocklehurst (brocklehurst71 (1971)). The variations of $`C`$(H$`\beta `$) along different regions (Tables 3 and 4) arise mainly from observational uncertainties (imperfect focussing at the blue and red ends of the spectra), and from different combinations of Balmer lines that could be reliably measured in each subspectrum to determine the local extinction. While each spectrum was corrected using its individual value of $`C`$(H$`\beta `$), images were extinction–corrected adopting a uniform value of $`C`$(H$`\beta `$)=0.09 and $`C`$(H$`\beta `$)=0.12 for Mkn 59 and Mkn 71, respectively (Guseva et al. guseva98b (2000)).
The electron temperature $`T_e`$(O iii) was derived from the observed flux ratio \[O iii\]($`\lambda `$4959+$`\lambda `$5007)/$`\lambda `$4363. In cases where the \[O iii\] $`\lambda `$4363 emission line was not detected, we used the so–called upper branch of the Edmunds & Pagel (edmunds84 (1984)) calibration of the total oxygen emission line flux \[O ii\]$`\lambda `$3727+\[O iii\]($`\lambda `$4959 + $`\lambda `$5007) vs. the electron temperature to determine $`T_e`$ and derived the oxygen abundance following van Zee et al. (1998a ).
The observed ($`F`$($`\lambda `$)) and corrected ($`I`$($`\lambda `$)) emission line fluxes relative to the H$`\beta `$ emission line fluxes for 3 regions in Mkn 59 (slit P.A. = 15) are listed in Tables 3 and 4. The tables contain only the regions where the emission line \[O iii\]$`\lambda `$4363 was detected at a $`S/N`$ that allowed for a reliable flux measurement. Also listed are the extinction coefficient $`C`$(H$`\beta `$), the observed flux of the H$`\beta `$ emission line and its equivalent width $`EW`$(H$`\beta `$) along with the equivalent width of hydrogen absorption lines $`EW`$(abs). The line intensities of the brightest regions (1 and 2) of Mkn 71 are presented in Izotov et al. (itl97 (1997)).
Applying the electron temperature $`T_e`$(O iii), ionic abundances of O<sup>2+</sup>, Ne<sup>2+</sup> and Ar<sup>3+</sup> were derived. The temperature $`T_e`$(O ii) was inferred, according to Izotov et al. (itl94 (1994), izotov96 (1996)), from the relation between $`T_e`$(O ii) and $`T_e`$(O iii) using H ii region photoionization models by Stasińska (stasinska90 (1990)). From $`T_e`$(O ii), the O<sup>+</sup> and N<sup>+</sup> ionic abundances were determined, while the intermediate value of the electron temperature $`T_e`$(S iii) served to derive the ionic abundances of Ar<sup>2+</sup> and S<sup>2+</sup> (Garnett garnett92 (1992)). The \[S ii\]$`\lambda `$6717/$`\lambda `$6731 ratio was used to determine the electron number density $`N_e`$(S ii). Total heavy element abundances were derived after correcting for unseen stages of ionization following Izotov et al. (itl94 (1994), itl97 (1997)) and Thuan et al. (til95 (1995)).
The resulting ionic and heavy element abundances for the 3 H ii regions in Mkn 59 with measured \[O iii\]$`\lambda `$4363 are given in Table 5 along with the adopted ionization correction factors (ICF) while the heavy element abundances in the two brightest knots of Mkn 71 are given in Izotov et al. (itl97 (1997)).
## 3 Results
### 3.1 Structural properties of the host galaxies
Recent studies have established that the underlying host galaxy of iE/nE BCDs does systematically differ with respect to its central surface brightness $`\mu _{E,0}`$ and exponential scale length $`\alpha `$ from other classes of dwarf galaxies such as dIs and dEs (P96b , Patterson & Thuan patterson96 (1996), Marlowe et al. marlowe97 (1997), Papaderos papaderosphd (1998), Salzer & Norton salzer98 (1998), Marlowe et al. marlowe99 (1999)). This structural dichotomy is evident from Figure 9 (adopted from Papaderos papaderosphd (1998)) showing that the central surface brightness and exponential scale length of the LSB component of a BCD with an absolute $`B`$ magnitude $`M_E`$$``$–16 mag are respectively $``$1.5 mag brighter and a factor of $``$ 2 smaller than in a typical dI/dE of equal luminosity. Although there is no sharp limit, a gap around $`\mu _{E,0}`$22 $`B`$ mag$`/\mathit{}\mathrm{}`$ separates the host galaxies of iE/nE BCDs from other classes of dwarf galaxies. The same diagram shows that more luminous BCDs (i.e. systems with a host galaxy brighter than –16 mag in the $`B`$ band) follow the same trend populating systematically different areas in the $`\mu _{E,0}`$$`M_E`$ and log($`\alpha _E`$)–$`M_E`$–planes than dIs and dEs. In Figure 9 (left panel) we show with open squares the positions of the LSB components of the iI,C BCDs studied here (values derived in the $`B`$ band) along with the ones of the iI,C BCDs SBS 1415+437 (Thuan et al. ti99 (1999)), Tol 1214–277 (Fricke et al. fricke99 (2000)) and UM 417 (Cairós et al. cairos99 (2000)). In the right panel of the same figure we include two further iI,C BCDs for which measurements of the exponential scale length of their LSB components are available, UM 133 (Telles & Terlevich telles97 (1997)) and Mkn 1328 (James james94 (1994)). It may be seen that the iI,C BCDs fit into the gap between typical iE/nE BCDs and dIs/dEs except for SBS 1415+437 and Tol 1214–277 which are indistinguishable from typical BCDs in the $`\mu _{E,0}`$$`M_E`$ and log($`\alpha _E`$)–$`M_E`$ parameter space.
The uncertain distance to Mkn 59 (cf. Section 1) does not affect the results stated above. As can be seen in both panels of Figure 9, the vector illustrating a shift to the data points due to a change of the Hubble constant is approximately parallel to the sequences of data points for iE/nE BCDs and other types of dwarf galaxies. For a different distance, Mkn 59 therefore remains within the gap between compact BCDs and dIs/dEs.
### 3.2 Heavy element and line intensity distributions
In Figure 10 we show the spatial intensity distributions of the H$`\alpha `$ and H$`\beta `$ lines along with that of the continuum adjacent to H$`\beta `$, as determined from the one–dimensional spectra extracted along the slit at P.A. = 59 of Mkn 59 (cf. Figure 6). It is obvious that all three distributions are spatially coincident. This is also the case for other nebular emission lines, including He i, \[N ii\], \[O ii\], \[O iii\] and \[Ne iii\], which show the same spatial distribution as the H$`\alpha `$ and H$`\beta `$ emission lines.
Wolf–Rayet stars were detected in region No. 1 of both, Mkn 59 and Mkn 71. For Mkn 71, 5 late nitrogen WR stars (WN stars) and 2 early carbon WR stars (WC stars) stars were found. The number of WC stars in Mkn 59 was determined to be 13 and that of WN stars to be 40 (Guseva et al. guseva98b (2000), values adapted to the distances used here). The spatial flux distributions of the blue $`\lambda `$4650 and red $`\lambda `$5808 WR bumps are shown in the right panel of Figure 10 along with the one of the nebular He ii $`\lambda `$4686 emission line. The red bump results primarily from WC stars, the blue bump from WN stars. We remark that the maxima of the H$`\beta `$ and WR emission are shifted by 1 pixel (0$`\stackrel{}{.}`$ 69) along the slit at P.A. = 59 suggesting that the locations of O and WR stars may slightly differ.
In all one–dimensional spectra of Mkn 59 at P.A. = 59 the \[O III\] $`\lambda `$4363 emission line is detected which allows a reliable determination of the electron temperatures and element abundances. Despite the large spatial variations of the oxygen line intensities and their ratios, the oxygen abundance is practically constant over the starburst region within the errors (Figure 11). The abundances of other heavy elements (neon, nitrogen, argon, sulphur, iron and chlorine) show the same, nearly constant spatial distribution along the slit with P.A. = 59 through the center of the giant H ii region. In spite of the small abundance variations within the central part, in most cases $``$ 0.2 dex, we remark that a weak gradient is present for the abundances of all heavy elements along the slit from –2″ to +3″, i.e. on scales of $``$ 260 pc. This suggests a local heavy element enrichment and may be related to the possible displacement of WR stars relative to O stars (see above).
In Figure 12 we show the spatial distribution of the oxygen abundance in Mkn 59 along the major axis with P.A. = 15. Three regions where the electron temperature was derived from the \[O iii\] $`\lambda `$4363/($`\lambda `$4959 + $`\lambda `$5007) flux ratio are indicated by stars, while open circles mark the regions where the \[O iii\] $`\lambda `$4363 line was not detected, thus oxygen abundances were derived following van Zee et al. (1998a ). The latter abundances appear systematically higher compared to those where the electron temperature could be directly constrained utilizing the \[O iii\] $`\lambda `$4363 line. However, the differences are still comparable to the intrinsic uncertainty of the order of $`\pm `$0.2 in log(O/H) ascribed by van Zee et al. (1998a ) to their empirical calibration. Taking this fact into account, as well as the weakness of the \[O iii\] $`\lambda `$4363 line in the outer parts of the galaxy, we shall consider the oxygen abundance as constant over the main body of Mkn 59 with an average value of 12+log(O/H) = 8.0.
The same analysis for Mkn 71 yields for the two brightest regions (regions 1 and 2; Figure 1, right), in which the \[O iii\]$`\lambda `$4363 line is observed, heavy element abundances of 12+log(O/H) = 7.79 and 7.77, respectively. The oxygen abundance in the regions 3 and 4 derived from the empirical calibration method described in van Zee et al. (1998a ) is larger by $``$ 0.3 dex.
### 3.3 Population synthesis models
For an analysis of the stellar content of the young ionizing clusters, we used the spectral energy distributions (SEDs) calculated by Schaerer & Vacca (schaerer98 (1998)) for heavy element mass fractions and ages in the range between $`Z`$ = 0.001–0.02 and $`t`$ = 0.1 – 10 Myr, respectively. For ages $`t`$ $``$ 10 Myr, we calculated a grid of SEDs for stellar populations with ages ranging from 10 Myr to 10 Gyr in time steps of $`\mathrm{\Delta }\mathrm{log}t[`$yr$`]=`$0.1 and a heavy element mass fraction of $`Z`$ = 0.002, using isochrones from Bertelli et al. (bertelli94 (1994)) and the compilation of stellar atmosphere models from Lejeune et al. (lejeune98 (1998)). An initial mass function (IMF) with a Salpeter slope (2.35) and lower and upper mass limits of 0.6 $`M_{}`$ and 120 $`M_{}`$was adopted.
To study the age of the stellar populations in the galaxies and to compare the results of the spectral analysis with broad band photometric data, it is necessary to take into account both the stellar and ionized gaseous emission in the spectra. For this purpose, the stellar SED has to be separated from the gaseous emission following the procedure described by Guseva et al. (guseva98a (1998)) and Papaderos et al. (papaderos98c (1998)). We added stellar SEDs calculated for instantaneous bursts with different ages (“single stellar populations”; SSPs) to the observed gaseous emission SED to match the total SED. The contribution of the gaseous emission was scaled to the stellar emission by the ratio of the observed equivalent width of the H$`\beta `$ emission line to the equivalent width of H$`\beta `$ expected for pure gaseous emission. To calculate the gaseous continuum SED at each region along the slit, the observed H$`\beta `$ flux and the electron temperatures were derived from the respective spectrum. The contributions of bound–free, free–free, and two–photon processes to the continuum emission were then calculated for the spectral range from 0 to 5 $`\mu `$m (Ferland ferland80 (1980), Aller aller84 (1984)).
### 3.4 The ages of the stellar populations
The number of O stars within the giant H ii region complexes at the southeastern tips of Mkn 59 and Mkn 71 was derived from the H$`\beta `$ flux to $``$ 4740 and 2010, respectively, following the prescriptions by Guseva et al. (guseva98b (2000)). As stated in Section 3.2, WR stars were detected in both BCDs, which, together with the large number of O stars, contribute strongly to the total SED. Taking into account the average metallicities, the ages of the brightest H ii regions in which WR stars were detected cannot exceed 4 to 5 Myr (e.g. Schaerer & Vacca 1998).
In Figs. 7 and 8 the modelled stellar SEDs, the gaseous SEDs and their coadded spectral distributions (marked as “total”) are shown superposed on the observed spectra for different regions of the BCDs. For nearly all spectra, a good agreement is achieved between the observed and modelled SEDs. Only in the case of region 2 of Mkn 71 (see Figure 1, right and Figure 8), the agreement is not as good due probably to an imperfectly focussed telescope.
Assuming that the observed gas emission is not caused by local SF, but due to photoionization or shock ionization by SF outside the observed region, the observed SEDs are in some cases reproducible with one SSP of intermediate age (few 10<sup>8</sup> yr) only; mostly, however, reproducing the observed SEDs required a superposition of SSPs with different ages and mass fractions. Tables 6 and 7 list the ages and relative mass fractions of the old and young SSPs which provide the best fits to the observed SEDs at different slit positions. The models imply respective ages of 4 and 2.9 Myr for the stellar population formed in the bright H ii region complexes in Mkn 59 and Mkn 71 (slit regions 1) while for other slit positions the age of the younger stellar continuum was inferred to few 10 Myr. The ages of the old SSPs could be constrained to $``$ 2 Gyr. Within the uncertainties of the methods applied, an upper limit of $``$ 3–4 Gyr is possible (see the following subsections), whereas even higher ages appear untenable. One has, however, to keep in mind that present spectrophotometric dating methods cannot definitely rule out the presence of a small fraction of even older stars, which due to their high $`M/L`$ would barely contribute to the SED.
#### 3.4.1 Uncertainties of the population synthesis models
The goodness of the spectral fits was found to sensitively depend on the mass fraction and ages of the adopted young and old SSPs. For instance, for region 1 of Mkn 59, increasing the age of the young SSP from 4 to 5 Myr yielded a satisfactory fit, but a 6 Myr old SSP resulted in systematic residuals between the observed and modelled SED. The age of the old SSPs could be constrained with an accuracy of $``$ 1 Gyr at slit regions where young stellar sources provide a minor contribution to the light (regions 3, 7 and 10 in Mkn 59 and region 3 in Mkn 71). In e.g. region 7 of Mkn 59, old SSPs up to 2 Gyr could still reproduce the observed SED, while already a 3 Gyr old SSP produced obvious residuals to the observed spectrum. A comparison of the observed and modelled $`EW(`$H$`\beta )`$ yielded further constraints to resolve ambiguities between different SED solutions. The assumption of high ages for the old SSP requires an increased contribution by the young SSP to reproduce the blue continuum, which in turn results in too high $`EW(`$H$`\beta )`$ as compared to the observed value.
Given that the metallicities of the stellar populations in BCDs may be lower than those of the ionized gas (Calzetti calzetti97 (1997), Guseva et al. guseva98a (1998), Mas–Hesse & Kunth mashesse99 (1999)), we computed a set of models varying the metallicity of the SSPs around $`Z=`$0.002 by 0.5 dex, the smallest stepsize our model libraries allowed. Models computed on the assumption of a metallicity other than $`Z=`$0.002 failed to adequately reproduce the observed spectrum, suggesting that the metallicity of the ionized gas is similar to that of the stellar population.
Alternatively, to constrain the age of the underlying galaxies, not only an instantaneous burst was considered, but also an extended episode of moderate SF at a constant rate. This way, the presumably complex formation histories of the galaxies were bracketed between two limiting cases.
In the same way as described in Section 3.3, a grid of model SEDs was calculated for stellar populations that were formed in extended SF episodes which started and ended at different times. The best–fitting solutions for e.g. Mkn 71 correspond to a formation episode of the underlying galaxy which started $``$ 2 Gyr ago and ended $``$ 10 Myr ago. Within the uncertainties, a SF episode that began up to 3 Gyr ago is possible; higher ages of the underlying galaxies would require a very recent ($`<10`$ Myr ago) termination of their formation episode to reproduce the observed SEDs. Such a recently ongoing SF activity in the underlying galaxy in Mkn 71 can be rejected from the colour–magnitude diagrams (Figure 4); in the left panel (a), obtained at the starburst region NGC 2363, stars are absent between the actual burst (younger than $``$ 10 Myr) and the last SF episode of the underlying galaxy which seems to have occurred more than $``$ 20 Myr ago (note the red supergiants at the tip of the 25 Myr isochrone). Also in the right panel (b), no stars much younger than $``$ 20 Myr are present. This indicates that, in the galaxy underlying the starburst, no significant SF occured during the last $``$ 20 Myr, in agreement with the ground–based CMDs by Aparicio et al. (1995). A similar analysis for Mkn 59 indicates, within the unvertainties, an upper age limit of $``$ 4 Gyr.
#### 3.4.2 Consistency with imaging data; further implications
The $`EW`$(H$`\beta `$) of region 1 in Mkn 59 and Mkn 71 amount to 150Å (Table 3) and 316Å(Izotov et al. 1997), respectively. For these regions the broad band colours calculated from the synthetic SEDs, including and ignoring gaseous emission (Tables 6 and 7) differ from each other by up to 0.5 mag. This demonstrates that a reliable dating of stellar populations within a starburst environment requires a correction for the colour shift induced by ample gaseous emission.
The consistency of the observed broad band colours with the spectral population synthesis results was checked by extracting the areas covered by the spectrograph slits from the colour maps of Mkn 59 and Mkn 71 (cf. Section 2.1.1). Figure 13 shows the modelled and observed colours along the slits after the latter were transformed to the Johnson–Cousins system (Bessell bessell90 (1990)).
Generally, the observed and synthesized colours agree well with each other (within 1$`\sigma `$ uncertainties), given that i) the synthesized colours are averaged values over regions covering several arcsec length along the slit, and ii) a minor displacement between the positions of a region in the colour map and of the slit may significantly alter the observed colours (see Figure 6 and the 1$`\sigma `$–variations of the observed colours in Figure 13), as the H ii regions are typically very compact.
Transforming the colours of the LSB component as derived from the colour profiles (Figure 3) to the Johnson–Cousins system yields a $`BR`$ $``$0.6 for Mkn 59 and Mkn 71. This is in good agreement with the upper limits of the colours derived from the spectra (cf. Tables 6 and 7). The constancy of the colour profiles for photometric radii where the starburst light contribution becomes negligible (i.e. for $`R^{}P_{25}`$) indicates that in the underlying galaxies population gradients are absent. Therefore, the upper age limits as being derived for either BCD along the spectrograph slits can be considered to be valid for the entire underlying galaxy.
## 4 Discussion
### 4.1 The metallicity distribution
Compared to BCDs in general, which show a metallicity distribution that peaks around $``$ 1/10 $`Z_{\mathrm{}}`$ and steeply decays towards lower metallicities (Kunth & Sargent kunth86 (1986)), the metal abundances of Mkn 59 ($`1/8Z_{\mathrm{}}`$) and Mkn 71 ($`1/14Z_{\mathrm{}}`$) are not exceptional. The nearly constant metallicity along the major and minor axes of Mkn 59, as well as along the minor axis of Mkn 71 (Figs. 11 and 12; see also González–Delgado et al. gonzalez94 (1994)), suggests that large–scale mixing processes in the ISM of both BCDs were at work on time scales of a few 10<sup>6</sup> yr. Otherwise, one would expect measurable metallicity enhancements in the vicinity of regions of ongoing or recent SF. Furthermore, the transport of metal–enriched warm and hot gas on scales of up to $``$ 1 kpc from the starburst region, an observational signature of which is the formation of supergiant shells, is commonly seen in the ISM of star–forming dwarf galaxies (cf. e.g. Marlowe et al. marlowe95 (1995), Hunter & Gallagher hunter97 (1997), Bomans et al. bomans97 (1997), Brinks & Walter brinks98 (1998), Papaderos & Fricke papaderos98a (1998), Strickland & Stevens strickland99 (1999); cf. also Figure 2, this work). The dilution of heavy elements in the vicinity of starburst regions may even be powered by galactic outflows. These were predicted by numerical models to develop within dwarf galaxies (Vader vader86 (1986), De Young & Gallagher deyoung90 (1990), De Young & Heckman deyoung94 (1994), Mac Low & Ferrara maclow99 (1999)) and proved necessary to account for the much lower metallicities observed in BCDs with respect to those predicted from closed–box evolutionary synthesis models (Krüger krueger92 (1992), Lisenfeld & Ferrara lisenfeld98 (1998)). The development of such a large–scale perturbation of the ambient gaseous component does not appear unreasonable given the high–velocity ($``$ 10<sup>3</sup> km sec<sup>-1</sup>) gaseous motions revealed spectroscopically by Izotov et al. (izotov96 (1996)) and Roy et al. (roy91 (1991), see also Figure 2, this work) in Mkn 59 and Mkn 71, respectively.
The small variations of the heavy element abundances along the giant H ii region complex of Mkn 59 (Section 3.2), as well as the possible spatial offset of the O– and WR stars, may be explained by a relocation of SF processes within a time span of few Myr. Propagating SF on a linear scale of $``$ 400 pc within the last 10 Myr has also been suggested by Drissen et al. (drissen99 (1999)), to account for the age differences of the young star clusters NGC 2363 I and II within Mkn 71.
### 4.2 Morphology vs. structural properties and age
The two iI,C BCDs investigated here appear similar with respect to the age and structure of their exponentially distributed host galaxies. This raises the question of whether the entire class of iI,C BCDs shares, besides a morphological resemblance, a set of common physical properties, such as an intermediate age and structural properties bridging the gap between iE/nE BCDs and dIs/dEs. A literature search for data on other nearby cometary BCDs does indeed provide some support to this hypothesis.
Mkn 1328$`/`$VCC 1374: NIR surface photometry (James james94 (1994)) yields an exponential scale length of $`\alpha `$=0.88 kpc for its LSB component. With its integrated absolute $`B`$-magnitude of –16 mag, and taking into account that a burst raises the $`B`$ luminosity of a BCD by typically $``$ 0.75 mag (P96b , Salzer & Norton salzer98 (1998)), Mkn 1328 resides presumably in between the parameter spaces populated by BCDs and dIs in the log($`\alpha _E`$)–$`M_E`$–plane (cf. Figure 9). From the integral $`(BH)`$ colour of $``$ 0.6 mag, the age of its host galaxy may be estimated to $``$ 6 Gyr, following the predictions by Krüger (krueger92 (1992)) and Krüger et al. (krueger95 (1995)) when a burst parameter $`b<`$ 0.01 is adopted. The metallicity was determined to 1/8 $`Z_{\mathrm{}}`$ (Kinman & Davidson kinman81 (1981)).
UM 133: For the observed metallicity (1/17 $`Z_{\mathrm{}}`$, Telles telles95 (1995)), the integral $`(RI)`$ colour of 0.4 mag derived by Telles & Terlevich (telles97 (1997)) is consistent with an age of $``$ 4 Gyr. With $`M_V`$=–18.25, and an exponential scale length of $`\alpha `$=1.39 kpc for its LSB component (Telles & Terlevich telles97 (1997)), this object falls also into the gap between BCDs and dIs.
UM 417 ($`Z`$$``$ 1/13 $`Z_{\mathrm{}}`$; Campos–Aguilar at al. campos93 (1993)). Surface photometry of the LSB component yields an intermediately compact structure, too ($`\alpha _E=`$ 0.45 kpc, $`M_E=`$ –14.4 $`B`$ mag and $`\mu _{E,0}=`$ 22.4 $`B`$ mag$`/\mathit{}\mathrm{}`$; Cairós et al. cairos99 (2000)). The same authors deduce a $`(BV)`$ color for the LSB component of $``$ 0.5 mag, compatible with an age of few Gyr.
SBS 1415+437 (Thuan, Izotov & Foltz ti99 (1999)): With a probable age of $``$ 100 Myr and an extremely low metallicity (1/21$`Z_{\mathrm{}}`$), this galaxy is a young galaxy candidate. Its LSB component, however, does not fit within the structural gap between BCDs and dIs; with $`\mu _{E,0}`$=21.0 $`V`$ mag$`/\mathit{}\mathrm{}`$, $`\alpha _E`$=0.30 kpc and $`M_E`$=–14.95 $`V`$ mag, it is rather comparable to a compact BCD. On the other hand, as this object apparently still undergoes its first major episode of SF, it is likely that its stellar LSB component has not yet been fully built but, as suggested by the age gradient along its major axis, delineates the trail along which SF has occurred.
Tol 1214–277: Similar to SBS 1415+437, this is an extremely low–metallicity iI,C BCD ($`ZZ_{\mathrm{}}/`$23, Terlevich et al. terlevich91 (1991)). Recent surface photometry studies with the VLT (Fricke et al. fricke99 (2000)) yield $`\alpha _E`$ 0.48 kpc and $`M_E`$ –16 $`B`$ mag. The average colours of the LSB component, $`(UB)`$ –0.4 mag and $`(BR)`$ +0.3 mag, suggest an unevolved galaxy with an age of few 10<sup>8</sup> yr.
The examples given above call for a further investigation of the following two hypotheses:
(i) Cometary BCDs are relatively young objects, with ages not exceeding a few Gyr, thus systematically younger than the majority of BCDs classified iE/nE.
(ii) Except for the candidate young extremely metal–deficient galaxies ($`\tau 10^8`$ yr) among them, cometary BCDs show structural properties of their host galaxies being typically intermediate between those of iE/nE BCDs and dIs/dEs.
Provided that the above hypotheses will be strengthened by an investigation of a larger sample of nearby iI,C BCDs, one may expect the iI,C morphology to occur more frequently in dwarf objects at higher redshifts. Indeed, examples of high-$`z`$ galaxies displaying a comet–like morphology have frequently been reported (cf. e.g. Dickinson dickinson96 (1996) for a galaxy cluster at $`z`$=1.15). Furthermore, a few medium-redshift objects with cometary morphology were included in the sample of Gúzman et al. (guzman98 (1998), Figure 1), notably the Faint Blue Galaxy (FBG) HERC 13088, at a redshift $`z`$=0.436. With a total luminosity of $`M_B`$=–21.5 mag and an exponential scale length $`\alpha `$ 2.8 kpc, this galaxy appears an upscaled version of the nearby iI,C BCDs studied here. It is worth noting that the rest–frame $`(BV)`$ colours of the LSB host of the latter FBG, $``$ 0.35 mag, also suggest an age of $``$ 5 Gyr.
These findings provide support to the hypotheses proposed in Section 1, namely that a strong off–center burst ontop a dwarf LSB component as observed in iI,C BCDs is primarily not a stochastic event but occurs in systems considerably younger than iE/nE BCDs, probably differing from the latter with regard to their structural properties. These results emphasize the need for a detailed investigation of the processes leading to the development of iI,C morphology in dwarf galaxies. We shall briefly remark on that issue in the next section.
### 4.3 An elongated structure of high gas density?
In the majority of BCDs, the centrally concentrated H ii complexes where massive SF takes place are typically found to almost coincide with maxima in the surface density of an extended H i envelope (Taylor et al. taylor94 (1994), Simpson simpson95 (1995), van Zee et al. 1998b ). Since a high gas density is necessary to sustain SF, such a condition is expected to be fulfilled near the end of the stellar body of an iI,C BCD. Age gradients of the star–forming regions along the stellar body’s major axis seem to be frequent attributes of such a system (Barth et al. barth94 (1994); Aparicio et al. aparicio95 (1995); Thuan et al. ti99 (1999)). These may be plausibly explained as due to propagation of SF activity (cf. Thuan et al. thuan87 (1987)) along the major axis of an elongated H i body.
Given these observational results, it seems that the H i halos of iI,C BCDs have elongated central concentrations with major axes that coincide with those of the stellar bodies. Indeed, low–resolution H i maps of Mkn 59 (Wilcots et al. wilcots96 (1996)) and Mkn 71 (Wevers et al. wevers86 (1986)) show the stellar components to follow prolate central concentrations within elongated H i clouds. A similar distribution of the gaseous halo does not seem rare among magellanic irregulars (Wilcots et al. wilcots96 (1996)) and has been interferometrically mapped also in BCDs in different evolutionary stages (see e.g. II Zw 40, van Zee et al. 1998b , and I Zw 18, van Zee et al. 1998c ). Moreover, the young BCD SBS 0335–052 (Izotov et al. izotov90 (1990)) is forming within an elongated gas cloud with a probably primordial chemical composition (Thuan & Izotov ti97 (1997), Lipovetsky et al. lipovetsky99 (1999)).
It might be argued that stellar bodies and H i concentrations of iI,C BCDs merely represent edge–on disks, as suggested by the slow rotation found for the H i envelopes of Mkn 59 and Mkn 71 as well as for the ionized gas of SBS 1415+437 (Thuan et al. ti99 (1999)). The generally exponential surface brightness profiles of their stellar LSB hosts (cf. Figure 3) do not serve as a proof of a disk structure (Freeman freeman70 (1970)), since exponential SBPs are also observed for e.g. small spheroidal systems (Binggeli et al. binggeli84 (1984)) and bars in late–type spirals (Elmegreen et al. elmegreen96 (1996)). However, the assumption of edge–on disks would demand the presence of numerous less–inclined iI,C BCDs, i.e. relatively blue disks with a single starburst in their outskirts, in contradiction to morphological studies of dwarf galaxies.
#### 4.3.1 iI,C systems vs. old iE/nE BCDs
Considering that only little data is available for BCDs with morphological types other than iE/nE, the construction of a tentative age–morphology sequence for gas–rich dwarf galaxies appears premature. The evolution of the different types of dwarf galaxies as well as their possible relations among each other are still sketchy, especially the triggering and effects of starbursts as well as the role of the Dark Matter (see Thuan thuan85 (1985), Davies & Phillipps davies88 (1988), P96b , Meurer meurer98 (1998), Swaters swaters98 (1998), Marlowe et al. marlowe97 (1997), marlowe99 (1999); van Zee et al. 1998b , Gil de Paz et al. gildepaz99 (1999)).
However, we argue that their relatively low ages would make iI,C BCDs a possible link between the extremely young ($`\tau `$ 100 Myr) dwarf galaxy candidates, most of which belong to the i0 class, and the evolved iE/nE BCDs, for which a lower age limit of 2 Gyr has been invariably derived (Krüger & Fritze–v.Alvensleben krueger94 (1994), Noeske noeske99 (1999)). In the young BCDs investigated thus far, a propagation of the SF activity has been observed (Papaderos et al. papaderos98c (1998), Izotov et al. izotov99 (1999), Thuan et al. ti99 (1999)). Thereby, it is conceivable that after a period of several 10<sup>8</sup> yr, these objects will gradually develop a cometary morphology with the youngest and most active star–forming regions located at one end of an elongated stellar body.
## 5 Summary and conclusions
Aiming at a better understanding of the different morphological subclasses of Blue Compact Dwarf Galaxies (BCDs) in the context of the evolution of star–forming dwarf galaxies, we have started a study of the “cometary” (iI,C ) subclass of BCDs. Contrary to the majority of BCDs where SF is confined to the inner part of an old circular or elliptical stellar LSB host, iI,C BCDs exhibit intense starburst activity close to one end of an elongated irregular LSB component. This intriguing morphology prompts two hypotheses: (i) iI,C BCDs are in fact dwarf irregulars (dIs) observed during a major stochastic enhancement of their otherwise moderate SF activity and (ii) a set of physical properties of the gaseous and stellar components favours the ignition of a burst with an amplitude comparable to that typically observed in other BCDs at the outskirts of an older LSB component.
For the two nearby cometary BCDs Markarian 59 and Markarian 71, a variety of deep ground–based and HST photometric and spectrophotometric data are available. To investigate the structural properties of the LSB component we derived surface brightness profiles (SBPs), corrected for systematic effects which can be important at low $`S/N`$ levels. Photometric properties of the underlying host galaxy and the superposed starburst component were derived by fitting a simple decomposition scheme to the SBPs; the radial extent of the starburst component as derived from this profile decomposition was verified using H$`\alpha `$– and colour maps. The spatial distributions of heavy element abundances in the starburst regions and over the main stellar body were derived using long–slit spectra. After the superposed emission by ionized gas had been semi–empirically modelled and subtracted from the original spectra, a population synthesis analysis was carried out at several positions along each slit to derive the properties of the underlying young and older stellar continuum. For Mkn 71 we also derived $`(BV)`$ colour–magnitude diagrams from HST data. The main findings of our analysis may be summarized as follows:
1. The azimuthally averaged intensity distribution of the underlying LSB host galaxy of both iI,C BCDs can be approximated by an exponential fitting law with a central surface brightness and a linear scale length intermediate between those typically inferred for iE/nE BCDs and dIs/dEs.
2. Spectral synthesis modeling of the starburst region and the main body implies the presence of an older population with a most probable age $``$ 2 Gyr for the underlying host galaxies of Mkn 59 and Mkn 71. Ages up to 4 Gyr for Mkn 59 and 3 Gyr for Mkn 71 are tenable within the model uncertainties, but higher values are unlikely.
3. The average oxygen abundances were determined to be 12+log(O/H)=8.0 (1/8 $`Z_{\mathrm{}}`$) for Mkn 59 and 12+log(O/H)=7.8 (1/14 $`Z_{\mathrm{}}`$) for Mkn 71, which are typical among BCDs. In addition, the metallicity distribution as derived for various elements in the vicinity of the starburst regions and along the major axis of the LSB-body shows only small scatter ($``$ 0.2 dex), suggesting that mixing of heavy elements has been efficient.
The similarity of the two objects with respect to the ages and structural properties of their LSB components motivated a search for published data on other cometary BCDs. We found that five objects for which the required data are available share similar properties with Mkn 59 and Mkn 71 with respect to the results (1) – (2). These findings suggest that the specific starburst morphology observed in iI,C BCDs comes along with distinct physical properties of their LSB host galaxies, i.e. is not attributable to stochastic processes only. Hypotheses which we consider worth investigating are:
(i) Cometary BCDs are relatively young objects, with ages not exceeding a few Gyr, thus systematically younger than “classical” BCDs of type iE/nE. If true, the development of iI,C morphology may represent a late evolutionary stage of an i0 BCD before it gradually assumes iE/nE–characteristics.
(ii) The underlying host galaxies of iI,C BCDs with an age $``$ 1 Gyr are moderately compact, in the sense that they show central surface brightnesses and exponential scale lengths intermediate between those typically derived for iE/nE BCDs and dIs/dEs.
The strongly extranuclear location of the starburst regions and the signatures of propagating star–forming activity along the main stellar body of an iI,C BCD suggest that the surface density distribution of cold H i gas in these systems resembles the optical morphology; published H i maps support this assumption. Spatially resolved interferometric studies will be of major importance for assessing the intrinsic processes regulating SF in iI,C BCDs and exploring possible evolutionary links to “classical” iE/nE BCDs.
###### Acknowledgements.
N.G.G. and Y.I.I. thank the Universitäts–Sternwarte of Göttingen for warm hospitality. We acknowledge the support of Volkswagen Foundation Grant No. I/72919. Research by K.G.N, P.P. and K.J.F. has been supported by Deutsche Forschungsgemeinschaft (DFG) grant FR 325/50–1, Deutsche Agentur für Raumfahrtangelegenheiten (DARA) GmbH grants 50 OR 9407 6 and 50 OR 9907 7. We thank the referee, Dr. R.C. Dohm–Palmer, for helpful comments and suggestions and C. Möller for providing us with a calibration spectrum. P. Papaderos thanks K. Bischoff for his assistance during the observations at Calar Alto. T.X. Thuan thanks the partial financial support of NSF grant AST–9616863.
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# Untitled Document
hep-th/0007226
Symmetry Restoration and Tachyon Condensation in Open String Theory
Rajesh $`\mathrm{Gopakumar}`$, Shiraz $`\mathrm{Minwalla}`$, and Andrew $`\mathrm{Strominger}`$
Jefferson Physical Laboratory
Harvard University
Cambridge, MA 02138
Abstract
It has recently been argued that D-branes in bosonic string theory can be described as noncommutative solitons, outside whose core the tachyon is condensed to its ground state. We conjecture that, in addition, the local $`U(1)`$ gauge symmetry is restored to a $`U(\mathrm{})`$ symmetry in the vacuum outside this core. We present new solutions obeying this boundary condition. The tension of these solitons agrees exactly with the expected D-brane tension for arbitrary noncommutativity parameter $`\theta `$, which effectively becomes a dynamical variable. The restored $`U(\mathrm{})`$ eliminates unwanted extra modes which might otherwise appear outside the soliton core.
1. Introduction
The general theory of relativity follows largely from the demand that the laws of physics take the same form in all coordinate systems. In string theory, the massless boson associated to this coordinate invariance - namely the graviton - is just one mode of an infinite tower of mostly massive string states. Associated to this infinite tower of modes is a stringy generalization of coordinate invariance. In the usual perturbative string vacuum, almost all of the string modes are massive and almost all of this stringy symmetry is accordingly spontaneously broken .
One may expect that string theory itself largely follows from the demand of stringy symmetry. However, despite the spectacular developments of the last five years, the nature of this stringy symmetry remains enigmatic.<sup>1</sup> A recent discussion can be found in . In this paper we investigate this issue of (open) stringy symmetry restoration in the context of a recent circle of ideas involving tachyon condensation, D-branes and noncommutative geometry. We will consider only the classical<sup>2</sup> Quantum effects could well be important in tachyon condensation, but we will not consider them. open bosonic string.
Following the work of Sen \[3,,4\], it is widely believed that the endpoint of the condensation of the open string tachyon is the closed string vacuum. There is by now compelling evidence for this conjecture from diverse points of view, including numerical computations \[5,,6,,7\] using Witten’s open string field theory . Moreover, Sen has argued (see also \[9,,10,,11\]) that D-branes in bosonic string theory can be viewed as solitons of the open string tachyon. Outside the core of the soliton the tachyon is in its ground state, and the theory is in the closed string vacuum with no open string excitations.
Recently Harvey, Kraus, Larsen and Martinec and Dasgupta, Mukhi and Rajesh have shown that turning on a large $`B`$ field enables an elegant realization of D-branes as tachyonic solitons. Techniques from noncommutative field theory can be used to construct the D-brane soliton in the $`\theta \mathrm{}`$ limit of large noncommutativity. The soliton and D-brane tensions agree exactly in this limit. A simple and beautiful explanation of the non-abelian structure of D-branes is found , with a natural embedding into string field theory .
However, even with these improvements several puzzles remain. In order to eliminate unwanted propagating open string states far outside the D-brane soliton core (i.e. in the closed string vacuum), one must assume that the coefficients in the tachyon-Born-Infeld Lagrangian take special values together with a special choice of field variables. Even with these assumptions, unwanted propagating modes persist inside the core in the bifundamental of $`U(N)\times U(\mathrm{}N)`$, where $`N`$ is the number of D-branes. Although plausible mechanisms \[16,,17\] for the elimination of these modes have been proposed, it is unsatisfying that these depend on unknown higher stringy corrections and cannot be seen directly from the Lagrangian employed in the analysis. In addition it is difficult to understand why $`\frac{1}{\theta }`$ corrections would not spoil the exact agreement found in between the soliton and D-brane tensions.
In this paper we consider the open bosonic string theory in the presence of a maximal rank $`B`$ field. We propose that in the process of tachyon condensation, as the tachyon rolls to its minimum, the noncommutative gauge field simultaneously rolls to a maximally symmetric configuration, about which the noncommutative gauge symmetry is fully unbroken, and becomes a linearly realized $`U(\mathrm{})`$.<sup>3</sup> We will refer to a configuration as having unbroken gauge symmetry if all fields are left invariant by the gauge transformations. Note that, with this usage, the usual perturbative vacuum of a gauge theory breaks local gauge invariance as $`\delta A0`$ for non-constant gauge transformations. The propagation of open string modes in this ‘nothing’ state is forbidden by the $`U(\mathrm{})`$ symmetry, and there is no need to invoke higher-order stringy corrections or special values of coefficients for their elimination.
We also modify the proposed identification of D-branes as noncommutative tachyon solitons by demanding that far from the core of the soliton, the solution approaches the nothing state, in which the $`U(\mathrm{})`$ symmetry of the noncommutative field theory (which is broken to a local $`U(1)`$ on the D-brane) is completely unbroken. We construct exact soliton solutions of the noncommutative tachyon-Born-Infeld Lagrangian obeying the modified boundary conditions, without expanding in $`\frac{1}{\theta }`$. It is further argued that these are exact-to-all-orders solutions of classical open string theory. The soliton tension exactly matches the expected D-brane tension. Furthermore, the propagation of open string modes far from the core, (i.e. in the nothing state) is forbidden as above by the $`U(\mathrm{})`$ symmetry. We regard these successes as evidence for the conjecture that the vacuum outside the D-brane core is the state of fully unbroken open string symmetries. On the other hand the situation for the bifundamentals is somewhat improved, but not fully resolved, as will be discussed in section 4.1.
One way of understanding the $`\theta `$-independence of the D-brane tension is that, in the context of tachyon condensation, $`\theta `$ is effectively a dynamical variable. In a sense (to be made precise herein), our proposal is that $`\theta `$ effectively relaxes to $`\mathrm{}`$ at the boundary.
While some puzzles are resolved in our approach, a significant new puzzle arises. In addition to the solutions corresponding to D-branes, there are a number of other spurious solutions obeying the same boundary conditions for which we have no physical interpretation. These must be understood or somehow excluded before the picture presented here can be regarded as complete.
2. The Action in Shifted Variables
2.1. The Action
The Euclidean action for $`U(1)`$ open bosonic string theory contains the terms \[18,,4\] (see also \[19,,20,,21\])
$$S=\frac{1}{G_o^2\alpha _{}^{}{}_{}{}^{13}(2\pi )^{25}}d^{26}x\left(V(T)\sqrt{\mathrm{det}(G+2\pi \alpha ^{}F)}+\frac{\alpha ^{}}{2}f(T)D_\mu TD^\mu T\sqrt{\mathrm{det}G}+\mathrm{}\right).$$
$`(2.1)`$
The tachyon potential $`V`$ has a maximum at $`T=T_{max}`$ corresponding to the unstable perturbative string vacuum and a minimum at $`T=T_{min}`$ which should contain no perturbative open string excitations. According to the minimum obeys
$$V(T_{min})=0,$$
$`(2.2)`$
and the maximum is determined from the D25-brane tension to be, in our conventions
$$V(T_{max})=1.$$
$`(2.3)`$
The universal coefficient of the potential term in (2.1) was demonstrated with worldsheet methods in . In addition it has been conjectured that $`f(T_{min})=0`$ \[19,,20\]. <sup>4</sup> Note that if $`f`$ is smooth at $`T_{min}`$ it can in any case be set to one by a field redefinition. This will not play an essential role in our analysis, although it is required in .
We wish to study the open bosonic string theory in the background of a constant $`B`$ field. According to \[22,,23\], the Euclidean action in this background continues to be given by (2.1), except that:
a. Space becomes noncommutative, i.e. all products in (2.1) are replaced by star products, with a noncommutativity tensor $`\mathrm{\Theta }`$, whose value is given below.
b. The parameters that appear in (2.1); the open string metric $`G_{\mu \nu }`$, the noncommutativity tensor $`\mathrm{\Theta }^{\mu \nu }`$ and the open string coupling $`G_o`$ are related to closed string moduli by the formulae
$$\begin{array}{cc}\hfill 2\pi \alpha ^{}& G^{\mu \nu }+\mathrm{\Theta }^{\mu \nu }=\left(\frac{2\pi \alpha ^{}}{g+2\pi \alpha ^{}B}\right)^{\mu \nu },\hfill \\ \hfill G_o^2& =g_{str}\sqrt{\frac{\mathrm{det}(g+2\pi \alpha ^{}B)}{\mathrm{det}g}}.\hfill \end{array}$$
$`(2.4)`$
Here $`g`$ and $`g_{str}`$ are the usual constant closed string metric and coupling.
We are thus led to study a $`U(1)`$ noncommutative gauge theory, interacting with a scalar field (the tachyon) that transforms in the adjoint of the gauge group. Note that we have not taken the $`\alpha ^{}0`$ scaling limit, so Born-Infeld corrections are retained. The noncommutative action (2.1) together with (2.4) is identical to that considered in (prior to taking the $`\theta \mathrm{}`$ limit). An alternate form of the action, used for example in \[13,,19,,20,,21\] differs by higher derivative tachyon terms which would not affect our conclusions.
We choose $`B_{\mu \nu }`$ so that space is maximally noncommuting, i.e. $`\mathrm{\Theta }`$ has maximal rank. We parameterize space with complex coordinates $`z^m`$, $`m=1,\mathrm{}13`$ obeying
$$[z^m,\overline{z}^{\overline{n}}]=i\mathrm{\Theta }^{m\overline{n}}.$$
$`(2.5)`$
2.2. Brief Review of the Operator Formalism
In this subsection we recall certain facts about noncommutative field theories, and especially noncommutative gauge theories, that we will use in our construction below. See, for instance, for more details.
The algebra of functions on a 26 dimensional noncommutative space is represented by operators on the Hilbert space of a thirteen dimensional particle. On this space, we define thirteen annihilation operators $`a_{\overline{m}}`$ and an equal number of creation operators $`a_m^{}`$
$$a_{\overline{m}}=i\mathrm{\Theta }_{\overline{m}n}^1z^n,a_m^{}=i\mathrm{\Theta }_{m\overline{n}}^1\overline{z}^{\overline{n}}.$$
$`(2.6)`$
These operators obey the commutation relations
$$[a_m^{},a_{\overline{n}}]=i\mathrm{\Theta }_{m\overline{n}}^1.$$
$`(2.7)`$
Several useful relations in translating from functions to operators are
$$\begin{array}{cc}\hfill d^{2n}x& (2\pi )^n\sqrt{()^n\mathrm{det}\mathrm{\Theta }}\mathrm{Tr},\hfill \\ \hfill _m& [a_m^{},],\hfill \\ \hfill _{\overline{m}}& [a_{\overline{m}},].\hfill \end{array}$$
$`(2.8)`$
We now consider a noncommutative gauge theory written in the operator language. The covariant derivative of a field $`\phi `$ that transforms in the adjoint of the noncommutative gauge group may be cast in the form
$$D_m\phi =_m\phi +i[A_m,\phi ]=[C_m,\phi ];D_{\overline{m}}\phi =_{\overline{m}}\phi +i[A_{\overline{m}},\phi ]=[C_{\overline{m}},\phi ]$$
$`(2.9)`$
where
$$C_m=iA_m+a_m^{},C_{\overline{m}}=iA_{\overline{m}}+a_{\overline{m}}.$$
$`(2.10)`$
The noncommutative field strength is
$$F_{m\overline{n}}=i[C,\overline{C}]_{m\overline{n}}\mathrm{\Theta }_{m\overline{n}}^1$$
$`(2.11)`$
where $`[C,\overline{C}]_{m\overline{n}}=[C_m,C_{\overline{n}}]`$ and $`\mathrm{\Theta }_{m\overline{n}}^1\mathrm{\Theta }^{\overline{n}p}=\delta _m^p.`$ The fields $`C_m`$, $`C_{\overline{m}}`$ transform homogeneously under gauge transformations. In particular, the field configurations $`C_m=C_{\overline{m}}=0`$ leave the gauge symmetry unbroken.
2.3. The Action and Equations of Motion
The noncommutative action (2.1) (for open string modes in the presence of a $`B_{\mu \nu }`$ field) can be rewritten in operator language as
$$\begin{array}{cc}\hfill S=& \frac{\sqrt{\mathrm{det}\mathrm{\Theta }}}{G_o^2\alpha _{}^{}{}_{}{}^{13}(2\pi )^{12}}\mathrm{Tr}[V(T)\sqrt{\mathrm{det}(G+2\pi \alpha ^{}(i[C,\overline{C}]\mathrm{\Theta }^1))}\hfill \\ & +\alpha ^{}f(T)[C_p,T][T,C^p]\sqrt{\mathrm{det}(G)}+\mathrm{}].\hfill \end{array}$$
$`(2.12)`$
Operators in (2.12) are appropriately ordered so as to reproduce string amplitudes; the precise ordering of operators in this action will not be important for us.
The tachyon equation of motion that follows from (2.12) is
$$2\alpha ^{}f(T)[C_m,[C^m,T]]\sqrt{\mathrm{det}G}\alpha ^{}f^{}(T)[C_m,T][C^m,T]+V^{}(T)\sqrt{\mathrm{det}M}=0,$$
$`(2.13)`$
where we have defined
$$M=G+2\pi \alpha ^{}(i[C,\overline{C}]\mathrm{\Theta }^1).$$
$`(2.14)`$
The equation for $`C_m`$ is
$$\frac{1}{2}\alpha ^{}[T,[C^m,T]f(T)]\sqrt{\mathrm{det}G}+i\pi \alpha ^{}[C_{\overline{n}},V(T)\sqrt{\mathrm{det}M}(M^1)^{m\overline{n}}]=0.$$
$`(2.15)`$
3. The Nothing State
In the variables (2.10), the usual vacuum with a single D25-brane is
$$T=T_{max},C_m=a_m^{}.$$
$`(3.1)`$
In it was conjectured that the ‘nothing’ state with no D25-branes is
$$T=T_{min},C_m=a_m^{}.$$
$`(3.2)`$
We would like to propose instead that the nothing state is
$$T=T_{min},C_m=0.$$
$`(3.3)`$
Due to (2.2), (3.3) and (3.2) are energetically degenerate.
Under a local $`U(1)`$ gauge transformation
$$\delta A_m=_mϵ+i[A_m,ϵ],$$
$`(3.4)`$
it follows from 2.7 that $`C_m`$ transforms as
$$\delta C_m=i[C_m,ϵ].$$
$`(3.5)`$
Hence, as remarked above, the nothing state (3.3) is fully invariant under this symmetry. The local $`U(1)`$ symmetry is restored to an unbroken $`U(\mathrm{})`$ symmetry of unitary transformations on the quantum mechanical Hilbert space.
We will now argue that fluctuations about the fully symmetric state (3.3) have no perturbative propagating open string degrees of freedom. (2.1) describes a noncommutative gauge theory whose matter fields all transform in the adjoint of the gauge group. Gauge invariance dictates that derivatives and gauge fields $`A`$ appear in the action only in the combination $`C_m`$. Thus fluctuations about any background with $`C_m=0`$ are governed by an action with no derivatives beyond those that appear in the star product. In particular, quadratic terms in the action have no derivatives (as the star product acts trivially on such terms); this statement is unchanged by nonlinear field redefinitions. Further, $`U(\mathrm{})`$ invariance ensures that explicit derivative terms are not dynamically generated. Thus, perturbatively, open string modes do not propagate about the background (3.3).
3.1. Nothing in Ordinary Variables
In it was shown that there is a nonlocal field redefinition which relates the noncommutative field strength $`F`$ to an ‘ordinary’ field strength, which we shall denote $`F^{\mathrm{ord}}`$, appearing in the commutative formulation of the same theory. Under this Seiberg-Witten map , a constant noncommutative field $`F`$ maps to a constant ordinary field strength, whose value is given by
$$F^{\mathrm{ord}}=F\frac{1}{1+\mathrm{\Theta }F}.$$
$`(3.6)`$
In (3.3) the noncommutative field strength $`F`$ takes the constant value $`\mathrm{\Theta }^1`$, and so corresponds to a divergent ordinary field strength<sup>5</sup> We are grateful to Jeff Harvey for explaining this point.. Schematically, (3.3) in ordinary variables, and the gauge $`B=0`$ takes the form
$$T=T_{min}F_{\mu \nu }^{\mathrm{ord}}=\mathrm{}.$$
$`(3.7)`$
By the second equation in (3.7) we mean that, in the vacuum state, $`F^{\mathrm{ord}}`$ is a rank 26 tensor, all of whose eigenvalues diverge.
Thus the conjecture of this section may be worded in ordinary variables as follows: The perturbative open string vacuum state with $`T=T_{max}`$, and a finite constant $`F^{\mathrm{ord}}`$ (in the gauge $`B=0`$) is unstable and decays to the nothing state, $`T=T_{min}`$ and infinite $`F^{\mathrm{ord}}`$. This is in part possible because at $`T=T_{min}`$, $`V(T)`$ vanishes and hence there is no energy cost to changing $`F^{\mathrm{ord}}`$.
3.2. Nothing as $`\mathrm{\Theta }=\mathrm{}`$
We have argued above that in ordinary variables the nothing state is given by (3.7), or equivalently by
$$T=T_{min},F_{\mu \nu }^{\mathrm{ord}}=0,B_{\mu \nu }\mathrm{}.$$
$`(3.8)`$
In order to analyze this state, we move to yet another set of variables; the gauge field whose noncommutativity is set by the large $`B`$ of (3.8). In terms of the new noncommutative $`F`$ (whose background value is zero in the nothing state), the action takes the form (2.1) with parameters
$$G_{\mu \nu }=(b\frac{1}{g}b)_{\mu \nu },\mathrm{\Theta }^{\mu \nu }=2\pi \alpha ^{}(\frac{1}{b})^{\mu \nu },G_o^2=g_{str}\sqrt{\frac{\mathrm{det}b}{\mathrm{det}g}},$$
$`(3.9)`$
where $`b_{\mu \nu }=2\pi \alpha ^{}B_{\mu \nu }`$. Notice that $`\mathrm{\Theta }^2\mathrm{Tr}(\mathrm{\Theta }G\mathrm{\Theta }G)=(2\pi \alpha ^{})^2\mathrm{Tr}(\frac{1}{g}b\frac{1}{g}b)\mathrm{}`$ in the limit under consideration. Thus, focusing on energies for which noncommutative phases are finite, explicit derivatives in the action (2.1) may be dropped. Defining a rescaled gauge field $`H_\mu =g_{\mu \alpha }(\frac{1}{b})^{\alpha \nu }A_\nu `$ (2.1) takes the form
$$\begin{array}{cc}\hfill S=\frac{1}{G_o^2\alpha _{}^{}{}_{}{}^{13}(2\pi )^{25}}d^{26}x& \sqrt{detG}(V(T)\sqrt{\mathrm{det}(\delta _\mu ^\nu +2\pi i\alpha ^{}[H_\mu ,H_\alpha ]g^{\alpha \nu })}\hfill \\ \hfill +& \frac{\alpha ^{}f(T)}{2}[H_\mu ,T][T,H_\nu ]g^{\mu \nu }+\mathrm{})).\hfill \end{array}$$
$`(3.10)`$
In the operator language
$$S=\frac{2\pi }{g_{str}}\mathrm{Tr}\left(V(T)\sqrt{\mathrm{det}(\delta _\mu ^\nu +2\pi i\alpha ^{}[H_\mu ,H_\alpha ]g^{\alpha \nu })}+\frac{\alpha ^{}f(T)}{2}[H_\mu ,T][T,H_\nu ]g^{\mu \nu }+\mathrm{}\right),$$
$`(3.11)`$
where we have used
$$\frac{\sqrt{\mathrm{det}\mathrm{\Theta }}\sqrt{detG}}{(2\pi )^{13}\alpha _{}^{}{}_{}{}^{13}G_o^2}=\frac{1}{g_{str}}.$$
$`(3.12)`$
Thus fluctuations about the nothing state are governed by the action (3.11), the dimensional reduction of the infinite $`N`$ open string field theory to a spacetime point. Note that $`B`$ does not enter into (3.11), consistent with the expectation that the end product of tachyon condensation is insensitive to the initial value of $`\mathrm{\Theta }`$.
4. D23 Branes
4.1. The Soliton Solution
We wish to find a soliton solution which is translationally invariant in 24 directions and approaches the nothing state in the complex transverse $`z^1`$ direction away from the core. For these purposes we take
$$\begin{array}{cc}\hfill G^{1\overline{i}}=\mathrm{\Theta }^{1\overline{i}}& =0,i=2,\mathrm{}13,\hfill \\ \hfill \mathrm{\Theta }^{1\overline{1}}& =\theta G^{1\overline{1}}.\hfill \end{array}$$
$`(4.1)`$
Consider the field configuration
$$\begin{array}{cc}\hfill TT_{min}& =(T_{max}T_{min})P_{N_1},\hfill \\ \hfill C_i& =P_{N_2}a_i^{},i=2,\mathrm{}13,\hfill \\ \hfill C_1& =0,\hfill \end{array}$$
$`(4.2)`$
where $`P_{N_k}`$ is a rank $`N_k`$ projection operator in the Hilbert space constructed from $`a_1^{}`$. For example we could take $`P_{N_k}`$ to be the projection onto the first $`N_k`$ states of the harmonic oscillator. Then the right hand side of (4.2) vanishes exponentially outside the soliton core, and the solution is asymptotic to the nothing state (3.3).<sup>6</sup> We note that in the $`\theta \mathrm{}`$ limit this solution is of the general form required for the string field theory construction described in . (In contrast, the approximate solutions found in \[13,,12\] have the same tachyon field but $`C_m=a_m^{}`$, and are asymptotic to (3.2). )
It is easy to check that (4.2) solves the equations of motion (2.13), (2.15). The first term in the tachyon equation (2.13) vanishes if we require
$$[P_{N_1},P_{N_2}]=0.$$
$`(4.3)`$
The second term vanishes because $`V^{}(T_{max})=0`$. (4.3) also implies the separate vanishing of both terms in the $`C^m`$ equation (2.15).
Of course, the true equations of motion that follow from the action (2.12) have an infinite number of terms (from the $`\mathrm{}`$ in (2.12)) that we have not considered here. However, each of these terms contains at least one factor of a covariant derivative of either $`T`$ or $`F`$. Since all covariant derivatives of $`T`$ and $`F`$ given in (4.2) vanish, additional terms in the equation of motion also vanish to all orders in the $`\alpha ^{}`$ expansion. Non-perturbative effects could alter the situation. It is rather surprising that an exact-to-all-orders solution can be constructed without even knowing what the Lagrangian is! Usually such constructions are possible only with supersymmetry: here it is a consequence of the magic of noncommutativity.
4.2. The Soliton Action
We will interpret solutions of the form (4.2) with
$$P_{N_1}=P_{N_2}=P_N$$
$`(4.4)`$
as $`N`$ coincident D-branes. Solutions of the Lagrangian (2.1) with $`P_{N_1}P_{N_2}`$ certainly exist but they do not correspond to conventional D-branes (the spectrum is wrong). The role of these solutions - or a rationale for their exclusion - must be understood before the picture presented here can be regarded as satisfactory. For now we consider (4.4). Using (4.1), (2.12) reduces to
$$\begin{array}{cc}\hfill S& =\frac{V(T_{max})\mathrm{Tr}P_N}{G_o^2\alpha _{}^{}{}_{}{}^{13}(2\pi )^{12}}\sqrt{\mathrm{det}\mathrm{\Theta }}\sqrt{\mathrm{det}G}\sqrt{1+\left(\frac{2\pi \alpha ^{}}{\theta }\right)^2}.\hfill \end{array}$$
$`(4.5)`$
We wish to rewrite this in terms of the coupling ($`G_o^{}`$), measure ($`\sqrt{\mathrm{det}G^{}}`$) and noncommutativity parameter ($`\mathrm{\Theta }^{{}_{}{}^{}i\overline{j}}`$) with respect to the 24 longitudinal dimensions. It follows from (2.4) that these are related to the 26-dimensional quantities by
$$\begin{array}{cc}\hfill G_o^2& =G_{o}^{}{}_{}{}^{2}\sqrt{1+\left(\frac{\theta }{2\pi \alpha ^{}}\right)^2},\hfill \\ \hfill \sqrt{\mathrm{det}G}& =G_{1\overline{1}}\sqrt{\mathrm{det}G^{}},\hfill \\ \hfill \sqrt{\mathrm{det}\mathrm{\Theta }}& =\theta G^{1\overline{1}}\sqrt{\mathrm{det}\mathrm{\Theta }^{}}.\hfill \end{array}$$
$`(4.6)`$
The trace gives
$$\mathrm{Tr}P_N=\frac{NV_{24}}{\sqrt{\mathrm{det}G^{}}(2\pi )^{12}\sqrt{\mathrm{det}\mathrm{\Theta }^{}}},$$
$`(4.7)`$
where $`V_{24}=d^{24}y\sqrt{\mathrm{det}G^{}}`$. Substituting into (4.5) and using (2.3) yields
$$S=\frac{NV_{24}}{G_o^{}_{}{}^{}2\alpha _{}^{}{}_{}{}^{12}(2\pi )^{23}}.$$
$`(4.8)`$
All $`\theta `$ dependence has disappeared, and this is exactly the action of $`N`$ parallel D23-branes.
4.3. The Spectrum
We now describe the spectrum of the solution (4.2) (4.4). We choose a basis in which
$$P_N=\underset{a=1}{\overset{N}{}}|aa|.$$
$`(4.9)`$
$`U(N)`$ Adjoint Fields
$`U(N)`$ adjoint fluctuations in the tachyon field can be expanded as
$$\delta T=\underset{a,b=1}{\overset{N}{}}T_{ab}(y)|ab|,$$
$`(4.10)`$
where $`y`$ is a longitudinal 24-dimensional coordinate and $`T_{ab}`$ is hermitian. As in , substituting into (2.12) reveals 24-dimensional tachyons in the adjoint of $`U(N)`$. A similar expansion gives $`U(N)`$ gauge fields. This is exactly the low-lying spectrum of $`N`$ bosonic D23-branes. Higher mass open string states on the D25-brane similarly descend to adjoint fields on the D23-branes, as in \[12,,15\].
$`U(\mathrm{}N)`$ Adjoint Fields
Derivative terms in modes of the form $`T_{jk}(y)|jk|+h.c.`$, where $`j,k>N`$ are projected out of the quadratic action because $`C`$ is proportional to $`P_N`$. Hence there are no propagating adjoint $`U(\mathrm{}N)`$ fields. In \[12,,13\] the gauge field does not have a transverse profile (as is consistent with the boundary condition (3.2)) and $`C`$ is proportional to the identity instead of $`P_N`$. In order to eliminate propagation of these modes, the additional assumption $`f(T_{min})=0`$ is required. Even then, if $`f`$ is quadratic or otherwise smooth about the minimum it may be set to one with a field redefinition. In these variables - which are the natural ones for studying propagation - propagating $`U(\mathrm{}N)`$ tachyons reappear. In any case, with the solution (4.2) the absence of such propagating modes is a natural consequence of the symmetries and no such additional assumptions about the $`f`$ prefactor or restrictions on field variables are necessary.
$`U(\mathrm{}N)\times U(N)`$ Bifundamental Fields
We may also consider $`U(\mathrm{}N)\times U(N)`$ bifundamental modes of the form $`T_{ak}(y)|ak|+h.c.`$ where $`aN,k>N`$.<sup>7</sup> In fact these modes can be gauged away or are eaten by the Higgs mechanism , but then similar comments pertain to the fluctuations of the gauge field. We consider the tachyon here for notational simplicity. Again, because of the projection operators in $`C`$, these modes do not acquire ordinary kinetic terms. They do however have a nonvanishing quadratic action involving fixed matrices. Substituting into (2.12) we get
$$S_{eff}(T_{ak})\mathrm{Tr}\left[a^ja_j^{}T_{ak}^2\right](j=2\mathrm{}13).$$
$`(4.11)`$
(4.11) is the action for a charged particle in a magnetic field of strength $`\frac{1}{\theta }`$. It has a discrete spectrum with spacing of order $`\frac{1}{\theta }`$, rather than a spectra of continuous momenta. In particular, there are no bifundamental excitations with energies below $`\frac{1}{\theta }`$.
Formally these modes disappear as the longitudinal noncommutativity $`\theta `$ is taken to zero, however higher order corrections to the action (2.1) appear to be suppressed by powers of $`\frac{\alpha ^{}}{\theta }`$, and hence cannot be ignored at small $`\theta `$. Hence we cannot make firm conclusions about the spectrum at small $`\theta `$.
5. Discussion
In this paper we have proposed answers to two puzzles relating to the condensation of the open bosonic string tachyon (in the presence of a $`B`$ field)
a. Why are there no open string excitations at the bottom of the tachyon well for any value of $`\mathrm{\Theta }`$ ?
b. Why is the condensed state at the bottom of the well independent of $`\mathrm{\Theta }`$?
We propose that as the tachyon rolls to its minimum, the gauge field also dynamically rolls to its maximally symmetric value (with nonzero field strength), and the fully unbroken gauge invariance prohibits perturbative propagation. This rolling is in part possible because, exactly at the bottom of the tachyon well, the coefficient of the Born-Infeld term in the action (2.1) vanishes, and there is no energy cost for changing a constant field strength. Using the Seiberg-Witten change of variables, this maximally symmetric configuration with nonzero field strength and finite $`\mathrm{\Theta }`$ can be reinterpreted as one with zero field strength and $`\mathrm{\Theta }=\mathrm{}`$.
Restated, we propose that $`\mathrm{\Theta }`$ (as set by the value of the commutative $`=F^{\mathrm{ord}}+B`$ at infinity) is effectively a dynamical variable that, regardless of its initial value, rolls to infinity in the process of tachyon condensation. The $`\mathrm{\Theta }`$ independence of the tension and spectrum of our soliton is a consequence of this dynamical nature of $`\mathrm{\Theta }`$. If this proposal is indeed correct, it would be very interesting to understand in detail the dynamics that sends $`\mathrm{\Theta }`$ to infinity, rather than any other (seemingly degenerate) value, as the tachyon rolls to its minimum.
In closing we note that the $`U(\mathrm{})`$ symmetry restoration described here is obviously closely related to the symmetry restoration in the cubic formulation of Witten’s open string field theory. It would be of interest to understand this connection in more detail.
Acknowledgements
We are grateful to J. Harvey, P. Kraus, F. Larsen, J. Maldacena, E. Martinec and A. Sen for useful discussions. Observations related to those of this paper have been independently made by the authors of reference (unpublished). This work was supported in part by DOE grant DE-FG02-91ER40654.
References
relax See for example, W. Seigel, “Covariantly Second Quantized String.2; 3,” Phys. Lett. B149, 157, 162 (1984); B151, 391, 396, (1984); T. Banks and M. Peskin, “Gauge Invariance of String Fields” Nucl. Phys. B264, 513 (1986). relax S. Shenker, “What are Strings Made of?” talk given at “String Theory at the Millenium”, http://quark.theory.caltech.edu/people/rahmfeld/Shenker/fs1.html relax A. Sen, “Descent Relations Among Bosonic D-branes,” Int.J.Mod.Phys. A14 (1999) 4061-4078, hep-th/9902105. relax A. Sen, “Universality of the tachyon potential,” JHEP 9912 (1999) 027 \[hep-th/9911116\]. relax A. Sen and B. Zwiebach, “Tachyon condensation in string field theory,” JHEP 0003 (2000) 002. relax N. Moeller and W. Taylor, “ Level truncation and the tachyon in open bosonic string field theory,” hep-th/0002237 . relax L. Rastelli and B. Zwiebach, “Tachyon potentials, star products and universality,” hep-th/0006240. relax E. Witten, “Noncommutative Geometry And String Field Theory,” Nucl. Phys. B268, 253 (1986). relax J. A. Harvey and P. Kraus, “D-branes as unstable lumps in bosonic open string field theory,” JHEP 0004, 012 (2000) \[hep-th/0002117\]. relax R. de Mello Koch, A. Jevicki, M. Mihailescu and R. Tatar, “Lumps and p-branes in open string field theory,” Phys. Lett. B482, 249 (2000) \[hep-th/0003031\]. relax N. Moeller, A. Sen and B. Zwiebach, “D-branes as tachyon lumps in string field theory,” hep-th/0005036. relax J. A. Harvey, P. Kraus, F. Larsen and E. J. Martinec, “D-branes and strings as non-commutative solitons,” hep-th/0005031. relax K. Dasgupta, S. Mukhi and G. Rajesh, “Noncommutative tachyons,” JHEP 0006, 022 (2000) \[hep-th/0005006\]. relax R. Gopakumar, S. Minwalla and A. Strominger, “Noncommutative solitons,” JHEP 0005, 020 (2000) \[hep-th/0003160\]. relax E. Witten, “Noncommutative Tachyons and String Field Theory”, hep-th/006071. relax V. A. Kostelecky and S. Samuel, “The Static Tachyon Potential In The Open Bosonic String Theory,” Phys. Lett. B207, 169 (1988). relax D. Ghoshal and A. Sen, “Tachyon condensation and brane descent relations in p-adic string theory,” hep-th/0003278. relax A. Sen, “Supersymmetric world-volume action for non-BPS D-branes,” JHEP 9910, 008 (1999) \[hep-th/9909062\]. relax M. Garousi, “Tachyon couplings on non-BPS D-branes and Dirac-Born-Infeld action”, hep-th/0003122. relax E. A. Bergshoeff, M. de Roo, T. C. de Wit, E. Eyras and S. Panda, “T-duality and actions for non-BPS D-branes,” JHEP 0005, 009 (2000) \[hep-th/0003221\]. relax J. Kluson, “Proposal for non-BPS D-brane action,” hep-th/0004106. relax V. Schomerus, “D-branes and deformation quantization,” JHEP 9906, 030 (1999) \[hep-th/9903205\]. relax N. Seiberg and E. Witten, “String Theory and Noncommutative Geometry”, hep-th/9912072. relax G. T. Horowitz, J. Lykken, R. Rohm and A. Strominger, “A Purely Cubic Action For String Field Theory,” Phys. Rev. Lett. 57, 283 (1986).
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# Phenomenological Evidence for Gluon Depletion in 𝑝𝐴 Collisions
## Abstract
The data of $`J/\psi `$ suppression at large $`x_F`$ in $`pA`$ collisions are used to infer the existence of gluon depletion as the projectile proton traverses the nucleus. The modification of the gluon distribution is studied by use of a convolution equation whose non-perturbative splitting function is determined phenomenologically. The depletion factor at $`x_1=0.8`$ is found to be about 25% at $`A=100`$.
PACS number: 25.75.Dw, 24.85.+p
It is conventional in the study of $`J/\psi `$ production in heavy-ion collisions that the gluon distribution before the hard subprocess of $`c\overline{c}`$ production is assumed to be the same as in a free nucleon . The unconventional view that the gluon distribution can be modified in the nuclear medium due to depletion was suggested in . In this paper we focus on $`p`$-$`A`$ collisions and show that the data on $`\alpha (x_F)`$ can be used to infer that gluon depletion in the projectile proton is not negligible.
Charmonium absorption in $`pA`$ collisions has been studied in without finding any satisfactory explanation for the $`x_F`$ dependence of $`\alpha (x_F)`$. In the effect of energy loss of partons is considered, but that is only one aspect of gluon depletion. Here we pay particular attention to the evolution of the gluon distribution of the projectile as it traverses the nucleus. The approximate absence of dilepton suppression and the consequent implication that the quark distribution is nearly unaltered by the nuclear medium lead some to expect that the gluon distribution would be unaltered also. However, such a view is based on the validity of DGLAP evolution of the parton distribution functions . We adopt the reasonable alternative view that the evolution in a nucleus is different from that of pQCD at high $`Q^2`$; indeed, we shall let the data guide us in determining the proper dynamics of the low-$`Q^2`$ non-perturbative process.
The Fermilab E866 experiment measured the $`J/\psi `$ suppression in $`p`$-$`A`$ collisions at 800 GeV/c with a wide coverage of $`x_F`$ . The result is given in terms of $`\alpha (x_F)`$, which is defined by the formula
$`R(x_F,A)=\sigma _A(x_F)/A\sigma _N(x_F)=A^{\alpha (x_F)1},`$ (1)
where $`\sigma _{N,A}`$ is the cross section for $`J/\psi `$ production by a proton on a nucleon $`(N)`$ or on a nucleus $`(A)`$. In a parametrization of $`\alpha (x_F)`$ for $`J/\psi `$ production is given:
$`\alpha (x_F)=0.952(1+0.023x_F0.397x_F^2)`$ (2)
for $`0.1<x_F<0.9`$. It is our aim here to explore the implication of Eq. (2) on the evolution of the gluon distribution.
Since the semihard subprocess of $`g+gc+\overline{c}`$ is common for $`p`$-$`N`$ and $`p`$-$`A`$ collisions, they cancel in the ratio $`R(x_F,A)`$ so the $`x_F`$ dependence can come from three sources: (a) the ratio of the gluon distribution in the projectile passing through a nucleus to that in a free proton, $`G(x_F,A)`$, (b) nuclear shadowing of gluons in the target, $`N(x_F,A)`$, and (c) hadronic absorption of the $`c\overline{c}`$ states after the semihard subprocess, $`H(x_F,A)`$. Putting them together, we have
$`R(x_F,A)=G(x_F,A)N(x_F,A)H(x_F,A).`$ (3)
$`G(x_F,A)`$ and $`N(x_F,A)`$ are ignored in . Since $`x_F<0.25`$ in , there is not much dependence on $`x_F`$ to be ascribed to $`H(x_F,A)`$, but in , where the full range of $`x_F`$ is considered, $`H(x_F,A)`$ is forced to carry the entire $`x_F`$-dependence by a fitting procedure, resulting in an unreasonably short octet lifetime. Our approach by including $`G(x_F,A)`$ and $`N(x_F,A)`$ in Eq.(3) is therefore complementary to the work of .
The nuclear shadowing problem has been studied in detail by Eskola et al. , using the deep inelastic scattering data of nuclear targets at high $`Q^2`$. On the basis of DGLAP evolution they can determine the parton distributions at any $`Q^2>2.25`$ GeV<sup>2</sup>. The results are given in terms of numerical parametrizations (called EKS98 ) of the ratio $`N_i^A(x,Q^2)=f_{i/A}(x,Q^2)/f_i(x,Q^2)`$, where $`f_i`$ is the parton distribution of flavor $`i`$ in the free proton and $`f_{i/A}`$ is that in a proton of a nucleus $`A`$. We shall be interested in the ratio for the gluon distributions only at $`Q^2=10\mathrm{GeV}^2`$, corresponding to $`c\overline{c}`$ production, and denote it by $`N(x,A)`$. From the numerical output of EKS98 we find that a simple formula can provide a good fit to within 2% error in the range $`40<A<240`$ and $`0.01<x<0.12`$; it is
$`N(x,A)=A^{\beta (x)},`$ (4)
where
$`\beta (\xi (x))=\xi (0.0284+0.0008\xi 0.0041\xi ^2),`$ (5)
with $`\xi =3.912+\mathrm{ln}x`$. Thus the $`A`$ dependence is minimal at $`\xi =0`$, corresponding to $`x=0.02`$.
The variable $`x`$ in Eq.(4) is the gluon momentum fraction in a nucleon in the nucleus, usually referred to as $`x_2`$. Both $`x_F`$ in Eq.(1) and $`x_2`$ in Eq.(4) are to be converted to the $`x_1`$ variable for the projectile nucleon, using
$`x_F=x_1x_2,x_1x_2=\tau M_{J/\psi }^2/s,`$ (6)
so that a part of Eq. (3) can be rewritten as
$`R(x_F,A)/N(x_2,A)=A^{\alpha (x_F(x_1))\beta (x_2(x_1))1}.`$ (7)
In our approach we treat $`H(x_F,A)`$ as having negligible dependence on $`x_F`$ for all $`x_F`$. Attempts to find that dependence have failed and led to the suggestion of the existence of an unaccounted mechanism responsible for the enhanced suppression in $`R(x_F,A)`$ at large $`x_F`$. In our view that mechanism is gluon depletion. Of course, if the $`x_F`$ dependence of $`H(x_F,A)`$ were independently known, its incorporation in our analysis is straigtforward. For us here, we identify the $`x_1`$ dependence of $`G(x_1,A)`$ in Eq.(3) with that in Eq.(7), which is completely known, and proceed to the study of the phenomenological implication on gluon depletion.
In the spirit of DGLAP evolution, even though the effect of a nuclear target on the projectile gluon distribution is highly non-perturbative, we now propose an evolution equation on the gluon distribution $`g(x,z)`$, where $`z`$ is the path length in a nucleus. For the change of $`g(x,z)`$, as the gluon traverses a distance $`dz`$ in the nucleus, we write
$`{\displaystyle \frac{d}{dz}}g(x,z)={\displaystyle _x^1}{\displaystyle \frac{dx^{}}{x^{}}}g(x^{},z)Q({\displaystyle \frac{x}{x^{}}}),`$ (8)
where $`Q(x/x^{})`$ describes the gain and loss of gluons in $`dz`$, but unlike the splitting function in pQCD, it cannot be calculated in perturbation theory. Equation (8) is similar to the nucleonic evolution equation proposed in , except that this is now at the parton level. Instead of guessing the form of $`Q(x/x^{})`$, which is unknown, we shall use Eq.(7) to determine it phenomenologically.
To that end, we first define the moments of $`g(x,z)`$ by
$`g_n(z)={\displaystyle _0^1}𝑑xx^{n2}g(x,z).`$ (9)
Taking the moments of Eq. (8) then yields
$`dg_n(z)/dz=g_n(z)Q_n,`$ (10)
where $`Q_n=_0^1𝑑yy^{n2}Q(y)`$. It then follows that
$`g_n(z)=g_n(0)e^{zQ_n},`$ (11)
whose exponential form suggests $`Q_n<0`$ for the physical process of depletion. The gluon depletion function $`D(y,z)`$ is defined by
$`g(x,z)={\displaystyle _x^1}{\displaystyle \frac{dx^{}}{x^{}}}g(x^{},0)D({\displaystyle \frac{x}{x^{}}},z),`$ (12)
where $`g(x^{},0)`$ is the gluon distribution in a free nucleon. From Eq. (12) we have $`g_n(z)=g_n(0)D_n(z)`$, where $`D_n(z)`$ is the moment of $`D(y,z)`$. Comparison with Eq. (11) gives
$`D_n(z)=e^{zQ_n}.`$ (13)
To relate this result to $`R(x_F,A)`$, we first note that $`G(x_F,A)`$ in Eq.(3) is, by definition, $`G(x_F,A)=g(x_1,A)/g(x_1,0)`$, where $`x_F`$ is expressed in terms of $`x_1`$. It then follows from Eq.(3) that
$`J(x_1,A)`$ $``$ $`g(x_1,0)R(x_F(x_1),A)/N(x_2(x_1),A)`$ (14)
$`=`$ $`g(x_1,A)H(A).`$ (15)
In relating $`A`$ to the average path length $`L`$ of the projectile $`p`$ through the nucleus, we use $`L=3R_A/2=1.8A^{1/3}`$fm. We then set $`z=L/2`$ for the average distance traversed at the point of $`c\overline{c}`$ production. Thus when referring to the last expression of Eq.(15), we write $`J(x_1,A)=g(x_1,z(A))H(z(A))`$, where $`g(x_1,z)`$ is to be identified with that in Eq.(12). Note that the $`A`$ dependence of the middle expression in Eq.(15) is, on account of Eq.(7), in terms of ln$`A`$, whereas that of the last expression is in terms of $`z`$, or $`A^{1/3}`$. Since it is known that ln$`AA^{1/3}`$ for $`60<A<240`$, we shall consider the consequences of Eq.(15) only for $`A`$ in that range. We suggest that a revised form of presenting the data, different that in Eq.(1), should be tried in the future.
Taking the moments of $`J(x_1,A)`$, we get using Eq.(11)
$`\mathrm{ln}J_n(A)\mathrm{ln}g_n(0)=zQ_n+\mathrm{ln}H(z).`$ (16)
To determine $`Q_n`$, it is necessary to use as an input the gluon distribution $`g(x_1,0)`$ in a free proton at $`Q^2=10\mathrm{GeV}^2`$. We adopt the simple canonical form
$`g(x_1,0)=g_0(1x_1)^5,`$ (17)
where the constant $`g_0`$ is cancelled in Eq.(16) due to the definition of $`J(x_1,A)`$. In our calculation we set $`g_0=1`$. Indeed, the accuracy of $`g(x_1,0)`$ is unimportant, since it enters Eqs.(15) and (16) in ways that render the result insensitive to its precise form. On the basis of Eqs.(7) and (17), $`J(x_1,A)`$ is therefore known. The LHS of Eq.(16) can then be computed except for a caveat. To calculate the moments of $`J(x_1,A)`$, it is necessary to compute $`_0^1𝑑x_1x_1^{n2}J(x_1,A)`$. However, $`x_1`$ cannot be less than $`\tau `$ in order to keep $`x_21`$ \[see Eq.(6)\]. Furthermore, Eq.(7) does not provide reliable information on $`J(x_1,A)`$ at small $`x_1`$, since the parametrizations of $`\alpha (x_F)`$ and $`\beta (x_2)`$ are for the variables in ranges that exclude the $`x_1\tau `$ limit. Fortunately, that part of the integration in $`x_1`$ can be suppressed by considering $`n3`$. The part of the integration in the interval $`0<x_1<\tau `$ amounts to only about 2% contribution even at $`n=2`$ (if naive extrapolation is used), so its inaccuracy will be neglected. Physically, it is the data at high $`x_F`$ that we emphasize in our analysis, and that corresponds to the high-$`n`$ moments of $`J(x_1,A)`$.
For convenience, let us denote the LHS of Eq.(16) by $`K_n(z)`$, i.e., $`K_n(z)\mathrm{ln}[J_n(z(A))/g_n(0)]`$. For sample cases of $`A=100\mathrm{and}200`$, they are shown as discrete points in Fig. 1 for $`3n20`$. Instead of performing an inverse Mellin transform on $`K_n(z)`$, our procedure is to fit $`K_n(z)`$ by a simple formula that can yield $`Q(y)`$ by inspection. The fitted curves shown by the solid and dashed lines in Fig. 1 are obtained by use of the formula
$`K_n=k_0+{\displaystyle \frac{k_1}{n}}{\displaystyle \frac{k_2}{n+1}}+{\displaystyle \frac{k_3}{n+2}}.`$ (18)
Using $`k_i`$ and $`k_i^{}`$ to denote the values for the cases $`A=100`$ and 200, respectively, we have
$`k_0=1.592,k_1=23.42,k_2=97.66,k_3=89.17`$ (19)
$`k_0^{}=1.831,k_1^{}=27.43,k_2^{}=113.97,k_3^{}=103.80.`$ (20)
Because of Eq.(16), the $`n`$ dependence of $`K_n`$ prescribes the $`n`$ dependence of $`Q_n`$. Let us therefore write
$`Q_n=q_0+{\displaystyle \frac{q_1}{n}}{\displaystyle \frac{q_2}{n+1}}+{\displaystyle \frac{q_3}{n+2}}.`$ (21)
Since Eq.(16) is to be used only for $`A>60`$, we evaluate it at $`A=100`$ and 200, and take the difference. Denoting $`z`$ by $`z_1`$ and $`z_2`$, respectively, for the two $`A`$ values, and with $`\mathrm{\Delta }k_i=k_i^{}k_i,\mathrm{\Delta }z=z_2z_1`$, we have
$`\mathrm{\Delta }k_0=q_0\mathrm{\Delta }z\mathrm{ln}{\displaystyle \frac{H(z_2)}{H(z_1)}},\mathrm{\Delta }k_i=q_i\mathrm{\Delta }z,(i0).`$ (22)
For the hadron absorption factor $`H(z)`$ we write it in the canonical exponential form , $`H(z)=\mathrm{exp}(\rho \sigma z)`$, where $`\rho ^1=(4/3)\pi (1.2)^3`$fm$`{}_{}{}^{3},z=0.9A^{1/3}`$fm, and $`\sigma `$ is the absorption cross section. Putting these in Eq.(22), we get (with $`\mathrm{\Delta }z=1.086`$ fm)
$`q_0+\rho \sigma =0.22,q_1=3.68,q_2=15.01,q_3=13.47`$ (23)
in units of fm<sup>-1</sup>.
There is a reason why $`q_0`$ and $`\rho \sigma `$ enter Eq.(23) as a sum. To appreciate the physics involved, we first note that Eq.(21) implies directly
$`Q(y)=q_0\delta (1y)+q_1yq_2y^2+q_3y^3.`$ (24)
The first and third terms on the RHS above are the loss terms (i.e., gluon depletion), while the second and last terms represent gain (i.e., gluon regeneration). If $`Q(y)`$ consisted of only the first term, then using it in Eq.(8) would give $`dg(x,z)/dz=q_0g(x,z)`$, whose solution is of the same exponential form as that of absorption. With both depletion and absorption present, the exponents lead to a sum, as in Eq.(23). Our $`Q(y)`$ is, however, more complicated. The $`q_2y^2`$ term gives rise to depletion that depends on the shape of $`g(x,z)`$, while the $`q_1y+q_3y^3`$ terms generate new gluons at $`x`$ from all the gluons at $`x^{}>x`$.
Since $`Q_n`$ decreases monotonically with $`n`$, we require $`Q_3<0`$, and exclude $`Q_2`$ from this consideration because of its inaccuracy discussed earlier. Combining Eqs.(18) and (20), we get $`\rho \sigma <q_0+\rho \sigma q_1/3+q_2/4q_3/5=0.05`$ fm<sup>-1</sup>. We thus set $`q_0=0.17`$ fm<sup>-1</sup>.
Since it is not easy to see directly from $`Q_n`$ or $`Q(y)`$ the magnitude of the effect of gluon depletion and regeneration, we can calculate $`g(x_1,z)`$, not from Eq.(12), but by fitting the calculated $`g_n(z)`$ in Eq.(11), using the formula
$`g_n(z)={\displaystyle \underset{i=1}{\overset{3}{}}}a_i(z)B(n1,5+i),`$ (25)
where $`B(a,b)`$ is the beta function. Then the result yields directly
$`g(x_1,z)={\displaystyle \underset{i=1}{\overset{3}{}}}a_i(z)(1x_1)^{4+i}.`$ (26)
For $`A=100(200)`$, i.e., $`z=z_1(z_2)`$, we have $`a_1=0.58(0.485),a_2=0.92(1.118)`$, and $`a_3=0.47(0.56)`$ for $`g_0=1`$ in Eq.(17). The result for $`G(x_1,z)=g(x_1,z)/g(x_1,0)`$ is then
$`G(x_1,z)=a_1(z)+a_2(z)(1x_1)+a_3(z)(1x_1)^2,`$ (27)
which is shown in Fig. 2 for two values of $`A`$. It is now evident that gluon depletion suppresses the gluon distribution at medium and high $`x_1`$, but the unavoidable gluon regeneration enhances the distribution at low $`x_1`$. The cross-over occurs at $`x_10.28`$.
Let us now exhibit our result for $`\alpha (x_F)`$, which is shown in Fig.3. The line is obtained by use of Eq.(27) in Eq.(3) and $`\sigma =6.5`$ mb in $`H(z)`$. Only one line is shown for both $`A=100`$ and 200, their difference being negligible in the plot. Since our method of using the moments cannot be extended to $`n=2`$ due to the problems mentioned after Eq.(17), there is some inaccuracy inherent in our analysis. Thus the fit cannot be expected to be perfect. Our model can reproduce the general trend of the $`x_F`$ dependence, but not the detail structure, for which more terms in Eqs.(18) and (21) would be needed. The overall suppression is achieved by use of a phenomenological value of $`\sigma `$, rather than the bound based on the technical assumption of $`Q_3<0`$.
Our analysis has been based on the assumption that $`H(z,A)`$ is independent of $`x_F`$. If and when that $`x_F`$ dependence can be determined independently, the effect can easily be incorporated in our analysis to modify our numerical result . Since that dependence is not likely to be strong , the modification would be minor. Our study shows that the $`J/\psi `$ suppression observed at large $`x_F`$ in $`pA`$ collisions strongly suggests the presence of gluon depletion in the beam proton at high $`x_1`$. The significance of this finding goes beyond the $`J/\psi `$ suppression problem itself, since it would revise the conventional thinking concerning the role of partons in nuclear collisions.
Since the gluon distribution is enhanced for $`x_10.28`$, the $`J/\psi `$ suppression observed in the $`x_F0`$ region in the heavy-ion collisions at CERN-SPS cannot be due to the gluon depletion effect. The same would be true at RHIC. However, we expect a significant increase in suppression at large $`x_F`$ due to gluon depletion, not to color deconfinement. We further speculate at this point that the gluon enhancement at low $`x`$ may be responsible, at least in part, for the strangeness and dilepton enhancement already observed in heavy-ion collisions.
This work was supported, in part, by the U.S.-Slovakia Science and Technology Program, National Science Foundation under Grant No. INT-9319091 and by the U. S. Department of Energy under Grant No. DE-FG03-96ER40972.
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# Absence of a metallic phase in random-bond Ising models in two dimensions: applications to disordered superconductors and paired quantum Hall states
## I Introduction
Ising models with quenched random bonds have been considered over many years. Negative couplings produce frustration and this is the starting point for the spin glass problem . A large class of models possess a “Nishimori line” in their phase diagram, on which the internal energy is analytic , and the correlation functions of the Ising spins obey certain identities . In two dimensions, the Ising model can be represented as a noninteracting fermion problem, even when the bonds are random . The problem then reduces to something similar to a two-dimensional (2D) tight-binding Hamiltonian with quenched disorder. Properties of the Ising model are then related to those of the fermion system, in particular to the fermion Green’s functions corresponding to the “Hamiltonian”, at a fixed “energy”, namely zero (this “energy” is not directly related to the energy in the sense of the Ising Hamiltonian). Then it is of interest to understand the properties of the fermion eigenstates near this energy, in particular whether they are localized or extended. In this paper, we consider such problems, and in particular argue that a recent proposal that there exists a phase of the Ising model in which the fermion eigenstates at zero “energy” are extended (a “metallic phase”) is ruled out. We also apply the results to paired fermion systems as in superconductors and quantum Hall states, which map onto similar noninteracting fermion problems.
Models of noninteracting fermions can in principle be considered using the methods of localization theory and random matrices. A list of symmetry classes (larger than the standard list due to Dyson) of ensembles of matrices was introduced by Altland and Zirnbauer (AZ) . The work of AZ was motivated by problems of disordered superconductors. Within the mean-field approximation, the fermionic quasiparticles of a superconductor are noninteracting, thus can be described using a single-particle formulation. The latter involves a Hamiltonian which in general contains quenched disorder, and could be a tight-binding Hamiltonian in 2D, for example. The energy levels of this Hamiltonian are the excitation energies of the quasiparticles. Once again, we may ask questions about the nature of the fermion eigenfunctions and eigenvalues. For superconductor problems, the natural zero of energy is a special point in the spectrum (unlike the case of a normal metal, for example) .
Among the symmetry classes found by AZ, one, denoted class D, describes disordered superconductors with broken time-reversal and spin-rotation symmetries. The symmetries are the same as those of the fermion problem in the two-dimensional (2D) random-bond Ising models (RBIM), and “energy” for the fermions of the Ising model corresponds to excitation energy for the fermions in the superconductor. The nonlinear sigma model for class D , which in effect defines this ensemble for dimensions greater than zero, has been shown, in the 2D case, to flow under the renormalization group to weaker values of the coupling constant . The coupling constant is related to the inverse of the thermal conductivity of the superconductor, and this flow implies that there is a phase in which there is a nonzero density of extended fermion eigenstates at zero excitation energy, and a superconductor described by this model would be in a thermal metal phase. We will refer to such a phase simply as a metallic phase. See also Refs. , respectively, for the 1D and 3D cases.
Senthil and Fisher considered possibilities for the application (via the fermion mapping) of results for class D to 2D RBIM’s. One scenario they discussed includes a metallic region in the phase diagram, below the Nishimori line, at relatively strong disorder and low temperature. They suggested that such a phase would have vanishing expectation values for both the Ising spin (“order”) and the dual “disorder” variables. Another scenario was that the metallic phase should be identified with the zero-temperature spin-glass region of a RBIM.
The preceding statements will be formulated more precisely in the course of this paper. Here we will begin by writing the Ising model Hamiltonian,
$$\beta =\underset{ij}{}K_{ij}\sigma _i\sigma _j,$$
(1)
where $`\beta =1/T`$ is the inverse temperature, the Ising spins $`\sigma _i=\pm 1`$, $`i`$, $`j`$ label sites of the lattice, and $`K_{ij}=J_{ij}/T`$ is a convenient notation for the Ising couplings (bonds). We will assume that $`J_{ij}`$ is zero unless $`i`$, $`j`$ are nearest neighbors on (say) the square lattice, and that there is a $`T`$-independent probability distibution for $`J_{ij}`$, such that the different nearest-neighbor bonds are statistically independent and identically-distributed. The statistical assumptions are not crucial and could be relaxed further, but we will see that it is important that the $`J_{ij}`$ are real, not complex. The partition function is then
$$Z=\underset{\{\sigma _i\}}{}\mathrm{exp}(\beta ),$$
(2)
where the sum is over all spin configurations $`\sigma _i=\pm 1`$ for all $`i`$. We will avoid discussing the boundary conditions on the lattice, or the thermodynamic limit, since we are mainly concerned with averages over the disorder of correlations of operators at separations that can be held fixed and far from the boundaries as the system size is taken to infinity after the disorder average.
We now recall a trivial fact, which will be central to the later arguments: the Ising spin correlation function for a fixed set of bonds $`J_{ij}`$,
$$\sigma _i\sigma _j\underset{\{\sigma _k\}}{}\sigma _i\sigma _j\mathrm{exp}(\beta )/Z,$$
(3)
is bounded above by $`1`$ and below by $`1`$:
$$|\sigma _i\sigma _j|1.$$
(4)
The bound is attained in the zero-temperature limit in pure or unfrustrated models, which include the antiferromagnetic models (all $`J_{ij}<0`$) on a bipartite lattice, as well as ferromagnetic (all $`J_{ij}>0`$) models. The bound follows from the Boltzmann-Gibbs probabilities $`\mathrm{exp}(\beta )/Z`$ being positive (and summing to 1), due to the reality of the couplings $`J_{ij}`$.
In this paper, we will discuss the statistics of the correlation functions in the order and disorder operators in a RBIM and in the class D nonlinear sigma model. Our central result is that in the metallic phase, the moments of either correlation function increase as powers of distance, which for the order (Ising spin) correlations eventually violates the upper bound, Eq. (4). This implies that the metallic phase described by the sigma model cannot occur in a RBIM as long as the couplings between the Ising spins are real. Our results apply to both non-zero and zero temperature in the Ising model. We trace the difference between the behaviors to differences in the form of the disorder, and suggest that the metallic phase may not after all occur in spinless or spin-polarized superconductors, or in paired fractional quantum Hall states with disorder.
In the remainder of this paper, we present our results. In Sec. II, we show that the Kadanoff-Ceva disorder correlation function in a RBIM has moments bounded below by one, and that its logarithm is symmetrically distributed, whenever the bonds are symmetrically distributed, as in an Edwards-Anderson (EA) spin-glass model. This relatively simple result will serve to illustrate points in the later discussion. In Sec. III, we obtain our central result, that the logarithms of the squared order and disorder correlations in the metallic phase are normally distributed, with mean zero and variance increasing as the logarithm of the distance, and hence the even moments of the correlations increase as powers of distance. Several steps are involved to set this up. An important point that arises along the way is that the distinctions between ensembles D, B, and BD, introduced in Ref. , are not important for local properties, such as these correlations. In Sec. IV, we consider another model, the O(1) model, and show that both its order and disorder correlations have properties like those in Sec. II. This model is most likely in the metallic phase. The crucial difference between such a model, and the RBIM, is that (in network model language, discussed in Sec. III) the disorder adds $`\pi `$ fluxes or vortices on one sublattice in the RBIM, but on both in the O(1) model; in Ising model language, the O(1) model corresponds to an Ising model with some couplings being complex. We also obtain the exact exponent for the mean order and disorder correlations at the critical point in another network model, the class C, or spin quantum Hall, model of Ref. . In Sec. V, we consider applications of our results to spinless or spin-polarized p-wave superconductors or paired fractional quantum Hall effect (FQHE) states. We show that independent insertion of vortices on a single sublattice corresponds to the RBIM situation, and cannot produce a metallic phase, at least at low densities. We argue that such “vortex disorder” destroys the Ising low-temperature ordered, or weak-pairing phase. For correlated vortices, the latter phase can occur, and there may be transitions in the universality classes found in the RBIM, rather than an intermediate metallic phase. Sec. VI is the conclusion.
## II Disorder Correlations for a Symmetric Distribution of Bonds
Our first result concerns the dual correlations in the EA spin glass case where the mean of $`J_{ij}`$ is zero. The two-point correlation of the Kadanoff-Ceva disorder variable $`\mu _\alpha `$ is defined in the following way (adapted from the pure case ). The disorder variables are associated with sites $`\alpha `$ of the (graph-theoretic) dual lattice, that is plaquettes of the original lattice. Given a choice of two such sites $`\alpha `$, $`\beta `$, we take the Hamiltonian (1) and modify it by reversing the sign of the $`J_{ij}`$’s on the links of the lattice crossed by a path on the dual lattice that runs from $`\alpha `$ to $`\beta `$. We can then construct the corresponding modified partition function $`Z_{\mathrm{mod}}`$. Then we define
$$\mu _\alpha \mu _\beta Z_{\mathrm{mod}}/Z.$$
(5)
This definition is independent of the choice of path from $`\alpha `$ to $`\beta `$, because of $`𝐙_2`$-gauge properties of the Ising model. Note that $`\mu _\alpha \mu _\beta >0`$ when the $`J_{ij}`$’s are real.
Now we consider the statistical properties of the disorder correlation function. We denote the average over the random bonds by an overbar, for example $`\overline{\mu _\alpha \mu _\beta }`$. We again make use of $`𝐙_2`$ gauge properties, this time of the distribution function for $`J_{ij}`$. There is a statistical $`𝐙_2`$ gauge invariance if the distribution is symmetric, $`P(J_{ij})=P(J_{ij})`$ for each $`i`$, $`j`$. However, such reversed bonds were exactly what was used in the definition of the disorder correlation. The set of bonds used in $`Z_{\mathrm{mod}}`$ occurs with the same probability, or probability density, as those in $`Z`$. Also, interchanging the original with the modified bonds exchanges $`Z_{\mathrm{mod}}`$ with $`Z`$. Hence $`\mathrm{ln}\mu _\alpha \mu _\beta `$ is symmetrically distributed, and
$$\overline{(\mathrm{ln}\mu _\alpha \mu _\beta )^m}=0$$
(6)
for $`m`$ odd, while
$$\overline{(\mathrm{ln}\mu _\alpha \mu _\beta )^m}0,$$
(7)
for $`m`$ even. For the correlation function itself, we have
$$\overline{\mu _\alpha \mu _\beta }=\overline{Z_{\mathrm{mod}}/Z}=\frac{1}{2}(\overline{Z_{\mathrm{mod}}/Z}+\overline{Z/Z_{\mathrm{mod}}})1.$$
(8)
The same argument works for any moment of the correlation function,
$$\overline{\mu _\alpha \mu _\beta ^m}1,$$
(9)
for any positive or negative integer $`m`$. The bounds are attained in the high temperature limit, where $`\mu _\alpha \mu _\beta =1`$.
We can predict how the disorder correlation function would behave in some well-known phases. In the paramagnetic phase, where $`\sigma _i\sigma _j0`$ as $`r_{ij}`$ (the distance between $`i`$ and $`j`$) goes to infinity, we expect that the mean disorder correlation goes to a constant at large distances, as in the pure case, and as in the high temperature limit. The constant must be $`1`$, and it appears that it will increase with decreasing temperature. We also expect that the width of the distribution of the logarithm of the correlation goes to a constant. A finite-temperature spin-glass phase is believed not to occur in 2D, but if it did we would predict that the distribution of $`\mathrm{ln}\mu _\alpha \mu _\beta `$ would have a width that goes as $`(C_1r_{\alpha \beta }^\theta +C_2)/T`$ at low temperature, where $`C_1`$, $`C_2`$ are positive constants, and $`\theta `$ is an exponent that characterizes the spin-glass phase as follows. In the spin glass, the insertion of the disorder variables induces a domain wall terminating at $`\alpha `$ and $`\beta `$. The wall is a fractal object, with a fractal dimension less than $`2`$, and its free energy, which is random and can be positive or negative, scales as $`r_{\alpha \beta }^\theta `$ . This exponent is believed to be the same one that enters the effect of reversing the boundary conditions, from periodic to antiperiodic, in one direction in a finite system of size $`L`$; the change in free energy scales as $`L^\theta `$. The exponent $`\theta `$ must be positive if the spin-glass phase is to be stable at finite $`T`$; it is found numerically to be negative, for continuous (e.g. Gaussian) distribution of $`J_{ij}`$, indicating that no finite $`T`$ spin glass phase exists in 2D . For some special discrete distributions, such as the bimodal $`\pm J`$ distribution (which has many degenerate ground states, giving an extensive entropy at $`T=0`$), $`\theta `$ is small and negative, or possibly zero . Finally, for a critical point, $`r_{ij}^\theta `$ in the width should be replaced by $`\mathrm{ln}r_{ij}`$ to a power $`1/2`$, but $`<1`$, when certain conditions hold, or most generally a function of $`r_{ij}`$ that is smaller than $`\mathrm{ln}r_{ij}`$ as $`r_{ij}\mathrm{}`$ (these follow from general results in Ref. ).
The result Eq. (9) is in stark contrast to the Ising (order) correlations. With a symmetric probability distribution, the odd moments of $`\sigma _i\sigma _j`$ vanish, and the even moments are $``$ 1. These opposite inequalities illustrate the extreme lack of duality for a symmetric probability distribution. If a metallic phase did occur in the RBIM, it would have to be due to frustration, as recognized in Ref. . It would then naively be expected to occur for a symmetric distribution of bonds. We have now shown that the idea of a phase in which the mean disorder correlation tends to zero is untenable in any RBIM with a symmetric distribution of bonds.
What has happened to the duality present in the pure 2D Ising model? Kramers and Wannier showed that the Ising model on the square lattice can be reformulated as a dual model on the dual lattice, with Ising spins $`\mu _\alpha =\pm 1`$, and dual couplings $`\stackrel{~}{K}`$. If the disorder variables become Ising spins, why does one not again obtain a correlation less than one? In the pure case, of course, one does. But the general relation between the original couplings $`K_{ij}`$ and the corresponding $`\stackrel{~}{K}_{\alpha \beta }`$ is
$$\mathrm{exp}(2\stackrel{~}{K}_{\alpha \beta })=\mathrm{tanh}K_{ij}.$$
(10)
For $`K_{ij}<0`$, $`\stackrel{~}{K}_{\alpha \beta }`$ has an imaginary part $`i\pi /2`$ (modulo a multiple of $`i\pi `$). The Boltzmann-Gibbs weights of the dual spin configurations become complex in general. However, the weights for a given nearest-neighbor bond $`\alpha `$, $`\beta `$ for the two values of $`\mu _\alpha \mu _\beta =\pm 1`$ differ simply by a sign. The disorder correlation for fixed $`J_{ij}`$’s then becomes a weighted average of $`\mu _\alpha \mu _\beta `$ (for arbitrary $`\alpha `$, $`\beta `$) with weights that sum to one but can be positive or negative. Hence the disorder correlation can be larger than one. Put another way, $`Z`$ may be smaller than $`Z_{\mathrm{mod}}`$, unlike the pure case.
For more general distributions, including those with a nonzero mean for $`J_{ij}`$ (which we can take to be positive without loss of generality), we cannot obtain a general result so easily. It is clear that when the disorder is weak (say, the standard deviation is small compared with the mean), there will be a ferromagnetically-ordered Ising phase, as in the pure Ising model, and in this the disorder correlation goes to zero at large distances. In order to rule out the existence of a metallic phase in the intermediate region with non-zero mean $`J_{ij}`$, another approach is needed.
## III Spin Correlations in the Metallic Phase
Now we turn to our second result, which directly concerns the metallic phase in the nonlinear sigma model for class D. We ask the question: if such a phase occurs in a random Majorana fermion model, what will be the behavior of the order and disorder correlations? We note immediately that the phase, as discussed in Refs. , is intermediate between two localized phases that would be identified with the paramagnetic and ferromagnetic Ising phases, which are still approximately dual to each other as in the pure Ising model. Then the intermediate metallic phase maps to itself under duality, and should treat the order and disorder correlations on an equal footing. The asymptotics of the two correlation functions should be similar.
In Sec. III A, we discuss the representation of the Ising model as a lattice free fermion quantum field theory, the relation of this to a network model, and the representation of order and disorder correlations in this language. In Sec. III B, we describe the nonlinear sigma model that is used to define the metallic phase. We argue that the distinctions between ensembles D, B, and BD , that differ globally, are not important for local correlations. In Sec. III C, we introduce the “twist operators” that represent the order and disorder operators in the nonlinear sigma model. Then in Sec. III D, we show that the statistics of the order correlations is incompatible with a RBIM with real bonds, but compatible with other models that violate the latter condition.
### A Fermion Representation
The metallic phase in the nonlinear sigma model for class D describes Majorana fermions, so we must consider the fermion representation of the Ising model. This can be set up in a variety of ways. The details are not in fact all that important here. The important points are that in fermion language, duality becomes rather self-evident, and both the order and disorder variables are represented as modifying the partition function by inserting an additional $`𝐙_2`$ fluxes or vortices seen by the fermions. A fermion propagating around a vortex picks up a phase factor $`1`$. The difference between the two operators is in the locations on the lattice at which they occur. The duality is most evident if the fermions are considered as moving on the “medial graph” of the original square lattice , as shown in Fig. 1 for a simply connected cluster. The medial graph of a given planar graph possesses two sublattices of plaquettes, on one of which each plaquette encloses a site of the original lattice, and on the other of which each plaquette encloses the center of a plaquette of the original lattice (i.e. a site of the dual lattice). The medial graph of the square lattice is another square lattice, as shown, and we consider only square-lattice clusters from here on. Note that we view the corners outside the cluster as nodes, so that the total number of links in the medial graph is a multiple of four.
The square of the Ising model partition function $`Z`$ can be represented as a six-vertex model on the medial graph, with free-fermion values of the parameters at each node . (There are also many other ways to represent the Ising model as a noninteracting fermion field theory on a decorated version of the square lattice. One such approach was used in pioneering work by Blackman and Poulter on the RBIM.) This free-fermion system is also equivalent to a (second-quantized representation of) the Chalker-Coddington network model , as has been emphasized recently . In other words, the single particle model underlying the fermion field theory is a network model. We omit a complete description of these models since they have been discussed so frequently in recent years, but an outline of the main points is as follows. The links of the medial graph square lattice are viewed as directed with an arrow on each link; the arrows circulate around the plaquettes, which implies that they circulate in opposite ways for the two sublattices of plaquettes (see Fig. 1). The particle propagates on the links of this medial graph in the direction of the arrows, picking up amplitudes that depend on the original random bonds $`J_{ij}`$. The amplitudes for each time step, during which the particle must move “forward”, following the arrows, to an adjacent link consistent with the arrows on the network, are elements of a unitary ($`S`$-) matrix assigned to each node. Thus the time-evolution is described by a unitary matrix $`𝒰`$, that is real in the present case, and has size a multiple of $`4`$. The sign of the product of amplitudes picked up by the particle propagating once around a plaquette determines whether a $`𝐙_2`$ flux or vortex (we use these terms, or $`\pi `$ flux, interchangeably) is present; a vortex is present when the sign is $`1`$. In the pure Ising model, such a vortex (a flux of $`\pi `$) is present on every plaquette. The insertion of negative $`J_{ij}`$ in the Ising model introduces an additional pair of $`𝐙_2`$ fluxes on the plaquettes of the medial graph that enclose the plaquette centers of the original lattice of the two plaquettes that are adjacent to the bond in question . When we speak of adding vortices or fluxes to plaquettes, the fluxes add mod $`2\pi `$, since the net phase picked up by the particle is what really counts; the gauge choices involved will not matter. The effect of negative $`J_{ij}`$’s in the Ising model is thus to add vortices, but only on one of the two sublattices of plaquettes of the medial graph network. By duality, vortices can also be produced in similar pairs on the other sublattice, by adding an imaginary term $`i\pi /2`$ to $`K_{ij}`$ (see Sec. II).
The squared partition function of the Ising model, $`Z^2`$, is now given by the second-quantized version of the network . The partition function for noninteracting fermions is generally a determinant of the inverse fermion propagator; in the present case, the propagator between two links is a sum over paths, given by the corresponding matrix element of $`1+𝒰+𝒰^2+\mathrm{}=(1𝒰)^1`$, so we have $`Z^2det(1𝒰)`$. Note that, unlike many other representations of the Ising model as a fermion field theory, in our case the matrix $`1𝒰`$ is not antisymmetric, so we cannot say that $`Z`$ is the Pfaffian of the same matrix.
Because $`𝒰`$ is unitary, its eigenvalues lie on the unit circle and may be written $`e^{iϵ}`$, where the eigenvalues $`ϵ`$ of $`i\mathrm{ln}𝒰`$ play the role of excitation energy eigenvalues, even though they are defined only mod $`2\pi `$. It is clear that for the long-time properties, such as the partition function, the important part of the spectrum of $`ϵ`$ is near $`ϵ=0`$. Since $`𝒰`$ is real, its complex eigenvalues come in complex conjugate pairs, while $`1`$ and $`1`$ are possible and will usually be nondegenerate. Also, since the network has a two-sublattice property (the particles hop from one type of link to the other alternately), the eigenvalues come in pairs $`e^{iϵ}`$, $`e^{iϵ}`$. This implies that if $`1`$, $`1`$ are present, then so are $`i`$, $`i`$, since the total number of eigenvalues is a multiple of $`4`$. Thus we could restrict attention to the range $`\pi ϵ\pi `$, which represents $`\mathrm{}`$ to $`\mathrm{}`$ in a continuum model. For the RBIM, the pair $`1`$, $`1`$ does not occur, $`det(1𝒰)>0`$, and the square root can be taken to obtain $`Z>0`$ . The case without the quadruplet $`1`$, $`1`$, $`i`$, $`i`$ (i.e. when $`det𝒰=1`$) corresponds to random matrices in class D, while the case with that quadruplet, $`det𝒰=1`$, corresponds to those in what has been termed class B . These random matrix ensembles are of matrices in the Lie algebras of SO($`2N`$), SO($`2N+1`$) (for some $`N`$), respectively . Matrices found in class B possess at least one, and typically only one, exact zero eigenvalue.
We now consider the calculation of the moments of the two-point functions of the order and disorder variables in the metallic phase of the nonlinear sigma model for class D in 2D. In terms of the medial graph or network model, the order and disorder operators are both represented as the ratio of a modified to the unmodified partition function, where order variables are represented by modifying the partition function by inserting vortices on the sublattice of plaquettes that correspond to the sites of the original lattice, and the disorder variables are vortices on the plaquettes that correspond to the plaquettes of the original lattice. Either partition function, when squared, is given by $`det(1𝒰)`$. As a check on the formulas, we can consider the order and disorder correlations in the pure case. An isolated vortex on a site of the original lattice carries a zero eigenvalue $`ϵ=0`$ in the high, but not in the low, temperature phase. For two vortices, the zero modes can mix and split away from zero, by an amount exponential in the separation when the latter is greater than the correlation length. At such large distances, the other eigenvalues tend to nonzero constants, so the behavior of the ratio of products of eigenvalues of $`1𝒰`$ is determined by the eigenvalues that tend to zero. Hence the correlation function tends to zero exponentially with distance in the high $`T`$ phase, but goes to a constant in the low $`T`$ phase. For the disorder correlation, the situation is reversed.
The average over the disorder of the ratio of determinants is performed by using either the replica method, with $`2n`$ copies of the system and $`n0`$, or the supersymmetry method, where $`2n`$ copies of the system are supplemented by $`2n`$ copies of the system with a certain kind of boson in place of the fermions, and no $`n0`$ limit. In the supersymmetry method, the bosons cancel the fermion determinants, as long as they are all unmodified. We will use the replica method, but the same results can easily be obtained with supersymmetry. For technical reasons, it is easiest to consider only the moments with $`m=`$ even of the correlation functions. Then we need to average the ratio of the $`m`$th power of the modified partition function to the unmodified partition function. Therefore we will modify the network for $`m`$ copies of the fermions so that they pick up an additional factor $`1`$ on propagating around either vortex (we can take these on the positions of the original sites, so as to obtain the spin-spin correlation function, but the disorder correlation is similar). The remaining $`2nm`$ fermions are unmodified. When $`n0`$, the partition function of the latter yields the division by $`Z^m`$. Thus in the average, the moment of the correlation function is simply the partition function of the replicated system at $`n=0`$, that is $`=1`$ when unmodified, but is not when $`m`$ of the Majorana fermions have been modified. That is
$$\overline{(\sigma _i\sigma _j)^m}=\overline{(Z_{\mathrm{mod}}/Z)^m}=\underset{n0}{lim}𝒵_{(m)}/𝒵,$$
(11)
where $`𝒵`$ stands for the partition function of the replicated and averaged system, $`lim_{n0}𝒵=1`$, and the subscript indicates that $`m`$ components have been modified.
As we will discuss further below, the nonlinear sigma model for the metallic phase in class D requires us to introduce an infrared regulator $`\eta >0`$ which can be viewed as an imaginary part of the energy (the real part being $`=0`$) at which we calculate the fermion Green’s functions, or as a corresponding shift in the energy eigenvalues. This is necessary in random fermion problems when the mean density of states at $`ϵ=0`$ is nonzero, so in general it can be included as a precaution. In the network model, it can be included by replacing $`𝒰`$ by $`𝒰e^\eta `$. We will need, first, to take moments in a finite size system with $`\eta >0`$, then take the system size to infinity, then let $`\eta 0`$. Some preliminary investigation suggests that for the moments of the ratio of determinants we consider, with finite separation of $`i`$ and $`j`$ (or $`\alpha `$ and $`\beta `$), this will give the same result as taking $`\eta =0`$ from the beginning, which is the strict definition for Ising models. This is for a fixed nonzero Ising temperature $`T`$. However, if one tries to take the temperature to zero before $`\eta 0`$, problems may arise. The reason is that, in the $`T0`$ limit, the fermions circulate around the plaquettes of the original Ising lattice, with amplitudes 1 (for $`\eta =0`$). The eigenvalues of $`𝒰^4`$ are then determined by the flux, either $`0`$ or $`\pi `$, on those plaquettes. Hence the eigenvalues $`ϵ`$ tend to either $`4ϵ=0`$ (mod $`2\pi `$) or $`4ϵ=\pi `$ (mod $`2\pi `$), and when the RBIM has a finite probability for any given plaquette to be frustrated, a finite fraction of eigenvalues $`ϵ`$ (and also the corresponding eigenvalues of $`1𝒰`$) tend to zero as $`T0`$ . In the modified partition function needed to obtain the spin correlation squared, the number of eigenvalues $`ϵ`$ that tend to zero is the same as in the unmodified partition function, since otherwise the spin correlation will go to zero or infinity, which is not the case. It is only these eigenvalues that are important in determining the spin correlation in the $`T0`$ limit. When the partition function is regulated with $`\eta `$, the corresponding eigenvalues of $`1𝒰e^\eta `$ tend to $`\eta `$, independent of $`ϵ`$, and the squared spin correlation goes to 1. This is expected in the case of a continuous distribution of random bonds, but definitely not for a bimodal ($`\pm J`$) distribution, where the $`T=0`$ spin correlation should be nontrivial. Thus the order of limits $`T0`$, $`\eta 0`$, makes a difference in this case.
### B Nonlinear Sigma Model
The claim about the metallic phase is that, in that phase, the partition function $`𝒵`$, and correlation functions, can be represented at large distances by those of the nonlinear sigma model for class D . In replica language, this model contains a field that takes values in the target manifold O($`2n`$)/U($`n`$). This may be parametrized by a $`2n\times 2n`$ complex matrix $`Q`$, which obeys $`Q=Q^{}`$, $`Q^2=I_{2n}`$, and $`Q^t=\mathrm{\Lambda }_xQ\mathrm{\Lambda }_x`$, where <sup>t</sup> denotes transpose, and $`\mathrm{\Lambda }_x=I_n\tau _x`$ ($`\tau _x`$ is a $`2\times 2`$ Pauli matrix). In terms of $`n\times n`$ blocks, the top right block $`V`$ of $`Q`$ is an $`n\times n`$ antisymmetric complex matrix. (A different parametrization is used in Ref. .) The symmetry operations are $`QOQO^{}`$, where in our basis, a matrix $`O`$ is in O($`2n`$) if $`O^1=O^{}=\mathrm{\Lambda }_xO^t\mathrm{\Lambda }_x`$ \[and in SO($`2n`$) if also $`detO=1`$\]. $`Q`$ can be written as $`Q=U\mathrm{\Lambda }_zU^1`$ for $`U`$ in O($`2n`$), where $`\mathrm{\Lambda }_z=I_n\tau _z`$ ($`U`$ should not be confused with $`𝒰`$). This represents the coset space O($`2n`$)/U($`n`$) because $`Q`$ is invariant when $`UUg`$, where $`g`$ is in the U($`n`$) subgroup of SO($`2n`$) parametrized in our basis as $`g=\mathrm{diag}(u,u^{})`$, where $`u`$ is a $`n\times n`$ unitary matrix \[thus, in U($`n`$)\], and $`u^{}`$ is the complex conjugate of $`u`$.
In Ref. , it was emphasized that O($`2n`$)/U($`n`$) has two disconnected components, corresponding to whether $`detU=\pm 1`$. For a zero-dimensional system, the other ensembles, termed classes B and BD, can be obtained by treating the component with $`detU=1`$ differently . These correspond to the existence of a single exact zero mode, $`ϵ=0`$, in a finite size system, with probability 1 (for class B) or $`1/2`$ (for class BD). Some network models (still with real $`𝒰`$) possess such zero modes, namely whenever $`det𝒰=1`$, and this can occur, depending on what fluxes are present, and the boundary conditions. We can avoid them by making appropriate choices of the latter. Even when present, they cancel in the regularized ($`\eta 0`$) ratios of determinants we consider in this paper. The reason is that when we insert two additional $`\pi `$ fluxes on the same sublattice, the determinant of $`𝒰`$ does not change. (However, an exact zero mode could still affect the other eigenvalues through level repulsion effects, for example.) While the presence or absence of such a zero mode may be important in random matrix ensembles for zero dimensional systems, or for global properties in higher dimensions, we do not expect it to play a role in local properties in more than zero dimensions, such as the correlations we consider here. (This applies to localized, as well as metallic, phases.) Therefore we expect that the distinctions between the nonlinear sigma models should not be important, and we will refer to class D/B/BD when this is so.
For dimensions larger than zero, a precise prescription for handling the two components of the target manifold has not been given. One would expect there could be domains of the two “phases” (in which the field $`Q`$ is on one or other of the two components of the target manifold). The domain walls would likely cost some action per unit length, and therefore additional parameters will be needed to specify the model. We would expect that there will then be a regime of parameters in which domain walls are costly and all domains of the “opposite” phase are small. Then the $`Q`$ field would essentially be globally on one component or the other. In that case, calculations can be done without domain walls as in other sigma models, but with a sum over the two phases. In fact, all existing proposals for a metallic phase in class D/B/BD neglect domain walls. Alternative phases where domain walls proliferate may exist, but have not been identified, and may not be metallic. In the absence of such proposals, we will consider the system without domain walls as defining the metallic phase we consider here. We note that the results in Ref. give a way to handle, in effect, the different components of the target manifold in a strong-coupling situation in dimensions $`2`$.
We further argue that the regulator $`\eta `$ which we introduce suppresses the second component. In the nonlinear sigma model, it introduces a term in the action of the form $`\eta d^2r\mathrm{tr}_{2n}\mathrm{\Lambda }_zQ`$, where $`\mathrm{tr}_{2n}`$ denotes a trace over $`2n`$-dimensional space. This term has to be minimized on each component to find the saddle point(s) about which perturbative fluctuations are expanded. We find that at such $`Q`$ values for the two components, where $`Q=\mathrm{\Lambda }_z`$, $`Q=O\mathrm{\Lambda }_zO^1`$, respectively, and $`O`$ represents a reflection in a hyperplane, the second component has relative weight like $`e^{\eta L^2}`$ compared with the first, where $`L^2`$ is the area of the system. Since we take $`L^2\mathrm{}`$ before $`\eta 0`$, we find that the second component is suppressed. This does not change the partition function $`𝒵=1`$ at $`n=0`$ for $`\eta =0`$, since for $`\eta =0`$ the functional integral over the first component gives $`1`$, and that over the second component gives $`0`$. (The use of just the first component, which includes $`U=I_{2n}`$, corresponds strictly to class BD .) Therefore, we drop the second component entirely, and no difference between the metallic phases in classes D, B, and BD will be seen in local correlations. (In the total density of states in the “ergodic” regime discussed in Ref. , smearing by energy resolution $`\eta `$ makes all three classes the same when $`\eta `$ is greater than the level spacing, of order $`L^2`$, consistent with this conclusion.)
### C Twist Operators
From a perturbative point of view, $`Q`$ arises from bilinears in the underlying Majorana fermions, which naturally leads to antisymmetric matrices. To obtain the correct structure of $`Q`$, it is essential that we start from the correct vacuum at weak coupling, represented by $`Q=\mathrm{\Lambda }_z`$, which is invariant under the U($`n`$) subgroup introduced above. The basis in which we gave $`Q`$ corresponds to the use of $`n`$ complex Fermi fields $`\psi `$ in place of the $`2n`$ Majoranas $`\xi `$. The diffusing (Goldstone) modes of the model involve only modes of the form $`\psi \psi `$ or $`\psi ^{}\psi ^{}`$ (the indices are suppressed), which correspond to the two off-diagonal blocks. (Goldstone modes corresponding to $`2n\times 2n`$ real antisymmetric matrices would give class DIII , in which time-reversal is unbroken.) This parametrization can also be arrived at using the O($`1`$) network model, which in first-quantized form is a single particle propagating on the medial graph network with fixed nodes of a standard form , and picking up $`\pm 1`$ factors (with independent probabilities $`1/2`$) on each link. Averaging over the group O(1)$`𝐙_2`$ in a replicated second-quantized representation leads to propagating Goldstone modes, and this model is in class BD, as described in Refs. .
In general, the modified partition function of the model is defined by the presence of a “twist” in the $`Q`$ field. The twist is a boundary condition at the points corresponding to $`i`$, $`j`$, that is obtained from the fact that $`m`$ of the Majorana fermion fields pick up a $`1`$. Since $`m`$ is even, this corresponds to a proper rotation $`O`$ in SO($`2n`$). We can choose the $`m`$ components of the fermions that are modified to be the real and imaginary parts of the first $`m/2`$ of the complex fermions that define our basis for $`Q`$. Then $`O`$ is represented by a matrix in the same U($`n`$) subgroup mentioned above, with $`O=\mathrm{diag}(u,u^{})`$, and $`u=\mathrm{diag}(1,1,\mathrm{},1,1\mathrm{})`$ with $`1`$ appearing $`m/2`$ times. Hence the modified partition function $`𝒵_{(m)}`$ is defined as the usual one but with the condition on the $`Q`$ fields at the points $`i`$, $`j`$, that on making a circuit around these points the $`Q`$ field is not periodic but changes as $`QOQO^{}`$, using the same $`O`$.
### D Result at Weak Coupling
As mentioned above, the nonlinear sigma model for class D/B/BD flows (when $`n=0`$) to weak coupling. Accordingly, we can compute the spin correlation function in the weak-coupling limit. To leading order, the action can be approximated as Gaussian for small $`V`$,
$$S=\frac{1}{2g^2}d^2r\mathrm{tr}_nVV^{},$$
(12)
where the trace is over the $`n\times n`$ matrices $`V`$, and $`g^2`$ here is the coupling constant squared, proportional to the inverse of the thermal conductivity $`\kappa _{xx}`$ . We have neglected the topological ($`\theta `$) term, since it plays no role in the following calculation. We have also omitted the leading nontrivial part of the $`\eta `$ term, $`\eta d^2r\mathrm{tr}_nVV^{}`$ with $`\eta >0`$. The limit $`\eta 0`$ is taken after the thermodynamic limit, because massless scalar fields in an infinite 2D system are problematic. In the weak-coupling limit, the twist operators take a simple form, since the operation described by our choice of $`O`$ acts linearly on $`V`$; $`V`$ transforms as the antisymmetric second-rank tensor representation of U($`n`$). The condition on $`V`$ on going around the points $`i`$, $`j`$ is that the components corresponding to complex fermions that are both modified or both unmodified are periodic, but those corresponding to one modified and one unmodified fermion pick up $`1`$. Thus $`m(2nm)/4`$ distinct (complex) components of $`V`$ pick up a $`1`$ on going around $`i`$ or $`j`$, and the remainder of the total $`n(n1)/2`$ pick up $`+1`$ (are periodic).
We now need the ratio of the modified to unmodified partition functions for the $`V`$ field with $`n0`$. This has the form of a standard problem in conformal field theory (the Gaussian theory is conformal since the coupling $`g`$ does not get renormalized). A twist operator of a single real massless scalar field, defined as a ratio of partition functions as here, has conformal weight $`1/16`$, and so its left-right symmetric correlation function decays as $`r^{1/4}`$. The exponent is doubled for a complex scalar, both of whose components are twisted. Multiplying these for our $`m(2nm)/4`$ complex components, we obtain
$$\overline{(\sigma _i\sigma _j)^m}r_{ij}^{m(2nm)/8},$$
(13)
which is the central result of this paper. Note that this result is independent of the coupling $`g`$. When $`n0`$, we obtain $`r^{m^2/8}`$, a positive power of distance. In the full non-linear sigma model, $`g^2`$ approaches zero logarithmically with distance when $`n=0`$ . From standard perturbative renormalization group arguments, we expect that the nonconstancy of the coupling produces at worst a factor of the form $`\mathrm{exp}[C^{}(m)(\mathrm{ln}r_{ij})^{\alpha (m)}]`$ on the right hand side, where $`\alpha (m)<1`$ is an $`m`$-dependent exponent. If the twist operator does not mix with any other operator in the RG, then the factor is only an $`m`$-dependent power of $`\mathrm{ln}r_{ij}`$.
In general at a random critical point, the logarithm of any correlation function is expected to have mean and variance depending logarithmically on the distance; the coefficients of these logarithmic dependences are universal. This arises because each extra factor of (say) 2 in distance is expected to contribute identically-distributed, essentially independent factors to the correlation function. The central limit theorem then applies to the distribution as $`r_{ij}\mathrm{}`$. (Here we assumed the moments exist. If the distribution of the logarithm of the factors in the correlation function is too broadly distributed for this to hold, then there is still a limiting distribution with universal properties, in particular the mean (or center) of the distribution varies as $`\mathrm{ln}r_{ij}`$, and the width increases as a universal power, between $`1/2`$ and $`1`$, of $`\mathrm{ln}r_{ij}`$, both with universal coefficients. For a more general discussion of the scaling forms, not assuming the product ansatz, see Ref. .) In our weakly-coupled Gaussian field theory, the moments in Eq. (13) have the form we would obtain by assuming the log of the squared correlation function is Gaussian-distributed; the mean and variance we would obtain are
$`\overline{\mathrm{ln}(\sigma _i\sigma _j)^2}`$ $`=`$ $`𝒪([\mathrm{ln}r_{ij}]^\alpha ^{}),`$ (14)
$`\overline{[\mathrm{ln}(\sigma _i\sigma _j)^2]^2}`$ $`=`$ $`\mathrm{ln}r_{ij}+𝒪([\mathrm{ln}r_{ij}]^{\alpha ^{\prime \prime }}),`$ (15)
where $`\alpha ^{}`$, $`\alpha ^{\prime \prime }`$ (both $`<1`$) are again some exponents. Thus these resemble the results for a random critical point, if we ignore the possible subleading corrections. Although it is well-known that the log-normal distribution is not uniquely defined by its moments, it is plausible that in the present problem the distribution is indeed asymptotically log-normal. Some consideration of diagrams for directly disorder-averaging powers of the logarithm of the ratio of determinants in some models, using the self-consistent Born approximation to obtain weak-coupling, also suggests that this is correct (the normal distribution is uniquely defined by its moments). Note that strictly we considered the limit $`g^20`$ (or $`r_{ij}\mathrm{}`$) for each $`m`$; this suffices to obtain “weak convergence” of the distribution. At fixed $`g^2`$ or $`r_{ij}`$, high moments, or the tails of the distribution, may not conform to the (log-) normal form.
The fact that Eq. (13) eventually exceeds $`1`$ implies that this behavior is impossible in any RBIM with positive Boltzmann-Gibbs weights. The metallic phase in class D/B/BD cannot occur in such a model. Instead, there can presumably be only gapped or localized phases and critical points between them (and possibly critical phases, meaning regions with scale invariance but described by a non-trivial fixed-point field theory, not a weakly-coupled nonlinear sigma model)—unless some other, so far unknown, stable metallic phase with the symmetries of the RBIM exists, that avoids the contradiction found here. This applies to the zero, as well as the nonzero, temperature region. As we saw, the regulated spin correlation goes to one as $`T0`$, for any distribution of bonds. Even though this is not the same as the correct, $`\eta =0`$, correlation for certain bond distributions, it is still in disagreement with the metallic phase.
We emphasize that results of a similar form can be obtained for the moments of the twist correlations in a variety of other metallic regimes in different ensembles, since these are by definition regions of diffusive behavior that can be described by a nonlinear sigma model at weak coupling. This is true even in systems that do not renormalize towards weak coupling, on length scales shorter than that for the crossover to strong coupling. Two other cases, class DIII and the symplectic (e.g. spin-orbit scattering) case of the Wigner-Dyson ensembles, both of which possess Kramers degeneracy due to time-reversal symmetry, flow to weak coupling in 2D like the class D/B/BD case considered here. However the physical significance of the Ising order correlation is less clear in these systems. Another case of interest is a family of nonlinear sigma models with target space SO($`2n+1`$)/U($`n`$), which with $`n0`$ arose in connection with the Nishimori line . The $`m=`$ even moments of $`\sigma _i\sigma _j`$ can be considered in this case also. The twist operator has the same form, but the total number of Goldstone modes is different: because SO($`2n+1`$)/U($`n`$) is the same, as a manifold, as SO($`2n+2`$)/U($`n+1`$), it is as above but with $`n1`$, not $`0`$. The family of models with SO($`2n+1`$) symmetry has two coupling constants in place of $`g^2`$ , but these do not enter the twist correlation at the Gaussian level. Thus the above result (13) can be used with $`n1`$, and the moments go as $`r_{ij}^{m(m2)/8}`$. For $`m>2`$, these increase with $`r_{ij}`$, eventually exceeding 1, requiring that $`\sigma _i\sigma _j^2>1`$ with nonzero probability. Thus the weak-coupling region of this family of sigma models is inaccessible in a RBIM with real couplings. There might in principle be metallic regimes in other weakly-coupled nonlinear sigma models in which the original Ising correlations are represented by a different sort of twist operator that gives a different result, but we are unaware of any at present.
According to recent work, certain network models are believed to possess a metallic phase . The model of Cho and Fisher is equivalent to an Ising model with couplings $`\pm K`$ on horizontal links (see Fig. 1), and $`K`$, $`K+i\pi /2`$ on vertical links , with independent probabilities $`1p`$, $`p`$ ($`K`$ is positive, and $`p`$ was denoted $`W`$ in Ref. ). From our remarks in Sec. III A, this can also be rephrased by saying that in the Cho-Fisher model, $`\pi `$ fluxes are added randomly in pairs, one above the other in Fig. 1, on both sublattices of plaquettes . On the $`p=1/2`$ line, the Cho-Fisher model is equivalent by a gauge transformation to the O(1) model described in Sec. III C above . \[The equivalence holds in the bulk but breaks down when we consider the boundary conditions; for certain boundary conditions, the Cho-Fisher model has the symmetries strictly of class D ($`1𝒰`$ has no exact zero eigenvalues).\] In Ref. , the Cho-Fisher model was re-examined numerically, and metallic behavior was found in a region including the whole of the $`p=1/2`$ line. We expect therefore that that model flows to the weak-coupling regime of the class D/B/BD sigma model, and then the above result applies. Hence we see that this does not contradict our claim that no metallic phase can occur in RBIMs with positive Boltzmann-Gibbs weights. The result for the O(1) model is not really surprising, in view of the behavior seen above in the dual of the RBIM, in which the Kramers-Wannier spins $`\mu _\alpha `$ have couplings $`\stackrel{~}{K}_{\alpha \beta }`$ with imaginary parts, and the moments of their correlations can be larger than one.
## IV Correlations in the O(1) Model
Our final result is for the O(1) model, already introduced in Sec. III C. We will argue that it is never in the ordered or disordered phases of the Ising model, by showing that the moments of the squared order and disorder correlations are both bounded below by 1. We also point out that the latter behavior is found in the other network models, in classes A, C. In the class C (spin quantum Hall) case, we find the exact exponent for the mean order and disorder correlations at criticality.
In the O(1) model, or the Cho-Fisher model at $`p=1/2`$, each plaquette of the network model (medial graph of the Ising model square lattice) encloses a flux of either $`0`$ or $`\pi `$ with independent probabilities $`1/2`$ (up to some boundary effects). We consider the order or disorder correlation functions, defined as before in fermion language. Like the disorder correlations in the RBIM with symmetric distribution of $`J_{ij}`$’s, the logarithm of either squared correlation is symmetrically distributed, and the even moments of the correlations are bounded below by 1. Note that this behavior is consistent with our results in Eq. (15), if the error term in $`\overline{\mathrm{ln}(\sigma _i\sigma _j)^2}`$ is zero. This means that if the O(1) model really does flow to the metallic phase, then these universal correction terms, and all higher-order analogs, in the nonlinear sigma model, must be zero.
Now we compare this with the behavior expected in the localized phases. Such phases occur at weak disorder (small $`p`$) in the RBIM and Cho-Fisher models. Like the two phases of the pure Ising (massive Majorana field theory) model, one or other mean (and mean square) correlation is supposed to decay to zero, and the other to go to a constant, at large distance. Hence the O(1) model is definitely not in either such phase. It is tempting to conclude that it must therefore be in the metallic phase, though this is not really proved; some other localized phase may not be ruled out. As mentioned above, numerical work led to the hypothesis of metallic behavior everywhere in the O(1) model.
The reason for caution about the last point is obtained by considering other network models for other ensembles. The two models in question are defined similarly to the O(1) model, but in the first, the particles pick up independent, uniformly-distributed U($`1`$) phases on the links, and in the second they pick up SU($`2`$) matrices instead (the latter requires two-component wavefunctions for the particles). These are respectively the Chalker-Coddington model for the integer quantum Hall transition (class A) , and the Kagalovsky et al. model for the spin quantum Hall transition (class C) . Both models possess localized phases away from their critical points. We now consider twist (“order” or “disorder”) correlations, defined as the ratio of modified to unmodified partition functions (fermion determinants) as before. First we note that for the class A model, $`𝒰`$ is a $`4N\times 4N`$ unitary matrix, its eigenvalues come in pairs $`e^{iϵ}`$, $`e^{iϵ}`$, and $`det(1𝒰)`$ is in general complex. For class C, $`𝒰`$ is an $`8N\times 8N`$ symplectic matrix, and its eigenvalues come in quadruplets, $`e^{iϵ}`$, $`e^{iϵ}`$, $`e^{iϵ}`$, $`e^{iϵ}`$, similar to class D; hence $`det(1𝒰)`$ is real and positive. This applies to both the modified and unmodified partition functions, and we see that in the U($`1`$) (class A) case we should consider the modulus square correlations, while for the SU($`2`$) (class C) case we can consider the correlations themselves, which are real and positive. The uniform distributions imply independent uniform distributions for the flux (in U(1) or SU(2) respectively) through each plaquette in these models. Multiplying these fluxes (as group elements) by $`1`$ (for the twist insertion operation) leaves the distributions unchanged, and hence again the logarithm of the \[modulus squared in the U($`1`$) case\] order or disorder correlations in these models are symmetrically-distributed, and the moments of the \[mod-squared for U($`1`$)\] correlations are bounded below by 1, even in the localized phases.
Thus it appears that these localized phases behave differently from those in the RBIM and Cho-Fisher models, and cannot be distinguished by Ising order or disorder variables. This may be connected with the continuous distributions for the flux in the plaquettes in these models, as opposed to the discrete distributions for the flux (which was either 0 or $`\pi `$) in the O(1) model. However, we may also point out that the localized phases in the RBIM and Cho-Fisher models, like the localized phases of a Majorana Fermi field with a weakly random mass, are expected to have vanishing density of states at $`ϵ=0`$, at least at weak disorder. In contrast, there is a possibility of a localized phase with the statistics of class D/B/BD, and the mean local density of states (which is independent of system size) near $`ϵ=0`$ would be expected to be nonvanishing \[as in the localized phase in the U(1) network model, which is in the unitary (class A) ensemble, though not the SU(2) model which is in class C\] and smoothly varying. Further, it would probably have an $`ϵ`$ dependence like that for class D in Ref. , including a peak at $`ϵ=0`$, but with the energy scale for such structure proportional to the inverse-square localization length. Although the existence of such a localized phase was predicted in, for example, Ref. , it has not so far been demonstrated to occur in practice in any model. In fact, given the symmetries of the problem, it is not clear why such behavior does not occur in the localized phases of the RBIM, or for the Majorana fermion with random mass. Perhaps such a phase would be consistent with the above form of probability distribution, and have neither order nor disorder correlations decaying to zero. If so, it may, like the metallic phase, be inconsistent with the Ising correlations in a RBIM.
Since we have been discussing the class C (spin quantum Hall) network model , we also include here a result for the critical correlations of the Ising order and disorder operators in that model. We can obtain a result only for the mean values, $`\overline{\sigma _i\sigma _j}`$ and $`\overline{\mu _\alpha \mu _\beta }`$. Each of these is defined by a twist of each of the two components of the wavefunction in a single copy of the system. It will be necessary here to use the supersymmetry method, in which the division by the unmodified partition function for a single copy is represented by a single two-component boson field . The partition function of the unmodified supersymmetrized system has supersymmetry osp($`2|2`$) $``$ sl($`2|1`$) , and is equal to 1. The mean of either correlation is represented by a modified partition function, in which two twist operators have been inserted. The partition function, like the unmodified one, has a graphical expansion as a sum over coverings by nonintersecting loops on the network model (medial) graph, with certain weights. States of the fermions and bosons flow around the loops; there are only three possible states, which can be labeled by the number of fermions they contain, either $`0`$, $`1`$, or $`2`$ . In the unmodified partition function, each loop is weighted with a factor 1, because the singly occupied state contributes $`1`$, and the other two states $`+1`$ each. (In Ref. , this mapping was constructed and used to show that several exponents for the spin quantum Hall transition are given by percolation, which has the identical loop expansion.) Because the original twist weights a fermion of either component that propagates once around the twist with a factor of $`1`$, the singly-occupied state picks up a $`1`$ and the others are unchanged. That is, a loop that encircles one of the twist insertions and not the other is now weighted by a factor $`3`$, not 1; the other factors which occur at the nodes are unchanged. It follows that either mean correlation is greater than one, as we proved by another method already. Since the maximum number, and the typical number, of such loops will increase with the separation, without limit in the critical case, we expect that either correlation increases as $`r^{2x}`$, where $`x<0`$ is the scaling dimension of the twist operator. In the loop model description, the twist operator has exactly the form recently considered by Cardy for percolation and other problems. Making use of his Eq. (1), $`x=(\chi ^2\chi ^2)/(2g)`$, with $`g=2/3`$, $`\chi =1/3`$ for percolation, and $`2\mathrm{cos}\pi \chi ^{}=3`$ for our twist, we obtain
$$x=\frac{1}{12}\frac{3\mathrm{ln}^2[(3+\sqrt{5})/2]}{4\pi ^2}0.154.$$
(16)
It is implicit in this result that we chose the branch for the logarithm such that $`x`$ is real. With this choice, Cardy’s general formulas imply that $`x`$ is negative whenever the factor for loops that enclose exactly one twist operator is larger than that, $`2\mathrm{cos}\pi \chi `$, for the unmodified loops.
## V Applications to superconductors and paired FQHE states
Paired states of fermions with complex (time-reversal violating) pairing of spinless or spin-polarized particles, or systems with broken time-reversal symmetry and spin-orbit scattering, have the same symmetries as class D/B/BD . We will consider only one-component systems such as spinless or spin-polarized fermions with p-wave pairing, which are the simplest, and begin with the pure case. We will then argue that their phase diagrams are more like that of the RBIM than has previously been recognized. For a RBIM in which frustrated plaquettes (of the Ising model lattice) are introduced independently, with some density, we argue that the Ising ordered phase is destroyed at $`T>0`$ for an arbitrarily small density of frustrated plaquettes (vortices). In the FQHE, this implies that the weak-pairing (nonabelian statistics) phase is destroyed by weak disorder.
It was important in Ref. for the discussion of nonabelian statistics that vortices carry a Majorana fermion zero mode when they occur in the so-called weak-pairing phase, but not in the strong-pairing phase. These phases occur on the two sides of the transition at which the mass of the Majorana fermions changes sign; the weak-pairing phase corresponds to the Ising ordered (low $`T`$) phase. The notion of a dual (in the Ising sense) vortex with the opposite properties—i.e. carrying a fermion zero mode only in the strong pairing phase—was implicitly discarded. The two types of vortices correspond in the network model to the two sublattices of plaquettes on which vortices (or fluxes) may be added to those already present in the pure model. The first type corresponds to adding a vortex on the network model plaquettes that correspond to plaquettes of the Ising model. In a continuum model, only the first type of vortex was considered because it was argued that in the physical situation a vortex should effectively have a region of strong-pairing phase, or vacuum (which was argued to be effectively the same thing in a “topological” sense), at its core.
The argument can be made essentially rigorous by considering a tight-binding Bogoliubov Hamiltonian on a lattice with a finite number $`N`$ of sites, with an edge, not periodic boundary conditions. Such a Hamiltonian corresponds to a $`2N\times 2N`$ matrix, which is in the Lie algebra of SO($`2N`$), and its eigenvalues come in pairs $`ϵ`$, $`ϵ`$, so that it never has an odd number of exact zero eigenvalues. A vortex can be generally defined as an object which, far from its core, approaches a singular gauge transformation that describes the insertion of a flux $`\pi `$ into the system without the vortex; this means that both the phase of the gap function, and the vector potential exhibit the winding by $`\pi `$. If we insert one of the postulated dual vortices in the strong-pairing phase with no other vortices present, then it is supposed to carry a zero-energy mode. There are no other zero modes with which it can mix, as we know because far from the vortex, we can use our understanding of the low-energy properties. In particular, there is no chiral spectrum of edge excitations in the strong-pairing phase. Hence it must be an exact zero mode, which is impossible for this Hamiltonian. A similar argument can be given in the weak-pairing phase, where the dual vortex does not carry a zero mode, but induces one on the edge that encloses it. We conclude then that no such dual vortices can exist, and there is only one type of vortex. Clearly we may take a continuum limit and draw the same conclusion. This means that Kramers-Wannier duality does not apply in such Hamiltonian models.
Turning to quenched disorder, it was shown in Ref. that randomly-inserted vortices in independent positions are a relevant perturbation of the pure Ising (Majorana) fixed point theory . It was pointed out that such disorder always occurs in the applications to FQHE systems, where underlying potential disorder can induce vortices, because they are charged. (In the RBIM, it is well-known that the vortices are correlated in pairs.) It seemed natural to expect such disorder to cause a flow to the metallic phase in class D/B/BD. Now if we assume that all the vortices are of the first type defined above, then we can construct network models of this situation by adding $`\pi `$ fluxes independently, all on the same sublattice, with some density $`p`$. This can be described as an Ising model with bond-disorder that is not independent for each bond, but the bonds are still real, and of fixed magnitude. This will accurately model the low-energy properties if the vortices are dilute (note that we assume the penetration depth and coherence length of the pure system stay finite at the transition). The results of this paper show that such a model cannot have the metallic phase in class D/B/BD.
In fact, such a model may not have a phase transition either. Introducing frustrated plaquettes into the Ising model independently tends to destroy Ising long-range order, though the discussion is complicated by the gauge choices needed in placing the strings of negative bonds that are needed to produce the frustrated plaquettes. We can avoid this difficulty by considering the spin-glass order parameter or correlations instead. Because the bonds are $`\pm J`$, there will be ground state degeneracy. Ground states can be represented by lines of frustrated bonds that join the frustrated plaquettes in pairs, chosen so as to minimise the energy. Distinct ways of dividing the frustrated plaquettes into pairs will frequently be exactly degenerate, and reconnecting the lines of frustrated bonds means reversing the spins in some domain. A condition that plausibly is necessary, but may not be sufficient, for the absence of long-range order at $`T=0`$, and hence of a finite $`T`$ phase transition, is that, in a ground state in the thermodynamic limit, any given spin lies, with probability 1, in a finite domain that can be flipped with zero energy cost. Heuristic considerations of sufficiently large domains suggest that any spin does lie in at least one such flippable domain (the probability for a domain, formed by reconnections of lines of frustrated bonds between nearby frustrated plaquettes, to have zero energy cost decreases as a power of the length of its boundary, while the number of such domains is exponential in this length). Hence we suspect that there is no ordering even at zero temperature in this model, for any nonzero density of the frustrated plaquettes (however, the zero-temperature state may be critical, as that of the usual $`\pm J`$ EA model may be also). This means that the model at $`T>0`$ is presumably in the paramagnetic phase, and all fermion eigenstates are presumably localized. The model can be generalized by introducing a continuous distribution for the magnitudes of the $`J_{ij}`$, and will then order at $`T=0`$, but a similar argument for the free energy at finite $`T`$ suggests that it will still not order at $`T0`$.
A stronger argument can also be given. The mean ground state energy is a function of the density $`p`$ of frustrated plaquettes, and is extensive, and varies smoothly with $`p`$ except at $`p=0`$. Increasing $`p`$ slightly means frustrating a small number of previously unfrustrated plaquettes. Thus frustrating one additional plaquette changes the mean total energy by an amount of order one. Yet this forces a domain wall from the plaquette to the edge of the system, along which the lines of frustrated bonds have been reconnected. The only reasonable conclusion is that the mean energy of the minimum energy domain wall is zero, except for a finite effect from around the added frustrated plaquette. This is the same behavior as in an EA spin-glass model. A similar conclusion holds for two added frustrated plaquettes, and may be compared with the discussion in Sec. II (but note that there the operation also unfrustrated originally-frustrated plaquettes, so that for a symmetric distribution of bonds the mean free energy change was exactly zero). We have no information about the width of the distribution of the domain wall energies, but we can expect that, like the usual short-range EA 2D spin glass models, there will be no finite $`T`$ transition. (A finite $`T`$ transition can occur in a 2D spin glass with sufficiently long-range power-law random bonds, but we have no reason to expect this to be realized here.)
Intuitively, the independently-inserted vortices appear similar to a random magnetic field, though the relation is not exact. However, it is a field that couples to the dual variables $`\mu _\alpha `$, and further it has a uniform component. The latter is the most important effect. A uniform magnetic field in a ferromagnetic Ising model destroys the transition, and the resulting phase has correlations like those in the ordered phase for the spin to which the field couples. In the present case, it would lead to the high temperature paramagnetic phase, in agreement with our conclusion.
The conjecture that vortex disorder destroys the Ising ordered phase has a dramatic consequence for applications to spinless or spin-polarized FQHE paired states. The paramagnetic phase corresponds to the disordered version of the strong-pairing phase; the weak-pairing phase has been destroyed. The weak-pairing or Moore-Read phase was supposed to be the basis for nonabelian statistics . We are arguing that this behavior, including the chiral Majorana fermion edge modes, is destroyed by any weak vortex (i.e. potential) disorder.
In models, such as the tight-binding Bogoliubov Hamiltonian, in which the vector potential and gap function each has only short range correlations, vortices will tend to be produced only in pairs, and again can be of only one type. Thus these one-component models appear similar to the RBIM, and may have a similar phase diagram, in which the weak-pairing phase occurs at weak but not at strong disorder. At the transition, the critical behavior may be that of either the pure Ising model (up to logarithms) or the low $`T`$ phase boundary in the (frustrated) RBIM. Thus the latter universality class could be realized in one-component superconductors. We conclude (in contrast to Refs. ) that, in at least some models of superconductors or FQHE paired states with the symmetries of class D/B/BD, and with only one type of vortex present, there may not be a metallic phase after all, but there may instead be a transition in a distinct universality class from the pure system for some types of disorder.
Clearly, similar possibilities should be explored in connection with other ensembles, which can occur when more components are present, but will not be considered further here. We point out, however, that the case of pairing of spin-1/2 fermions, with spin-orbit scattering and a general random mass, which has the symmetries of class D/B/BD, corresponds to multicomponent models considered in Ref. , where it is argued that a metallic phase is produced. Ref. also argued that the metallic phase would not occur in the one-component case in the absence of vortex disorder, but did not consider vortex disorder as fully as we have. Also, models of superconductors as disordered grains (each described by a random matrix from class D), coupled by weak hopping, appear similar to multicomponent models, and may have a metallic phase, as a mapping to a weakly-coupled nonlinear sigma model would suggest.
## VI Conclusion
The main point to emerge from this study is that the important difference between the RBIM and the (one-component) models which possess a metallic phase is that in the former the added vortices, in network model language, occur on one sublattice only, but on both in the latter. This appears to be a necessary condition for the existence of the metallic phase. If the vortices are correlated in pairs, sufficiently strong disorder will be required to produce the metallic phase. This result casts some doubt on whether the metallic phase will occur in applications to one-component 2D superconductors and paired FQHE systems, because these possess only one type of vortex, corresponding to those on only one sublattice in the network. On the other hand, if vortices of only one type are present, but are uncorrelated, this may lead to the destruction of one of the phases, and hence of the phase transition. If such vortices are correlated in pairs, the phase diagram may resemble that of the RBIM. It would be interesting to test this numerically, both on models defined by a Hamiltonian, such as tight-binding models, and on network models, and also for other symmetry classes.
###### Acknowledgements.
We are grateful to I. A. Gruzberg, J. T. Chalker, M. R. Zirnbauer, and D. S. Fisher for discussions. N.R. is grateful to David Gross and the staff of the Institute for Theoretical Physics, University of California, Santa Barbara, for their hospitality while this paper was being written. This work was supported by the NSF under grant no. DMR-98-18259 (NR). Work of N.R. at the ITP was partially supported by the NSF, under grant no. PHY94-07194.
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# Critical dynamics of two-replica cluster algorithms
## I Introduction
The Swendsen-Wang (SW) algorithm and related cluster methods SwWa ; Wolff ; KaDo ; ChMa97a ; ChMa98a ; WaSw90 ; NeBa99 have greatly improved the efficiency of simulating the critical region of a variety of spin models. The original SW algorithm can be modified to work for spin systems with internal symmetry breaking fields DoSeTa . Spin models of this kind include the Ising antiferromagnet in a uniform field, the random field Ising model and lattice gas models of adsorption in porous media DuMaSaAu . The modification proposed in Ref. DoSeTa, is to assign Boltzmann weights depending on the net field acting on the cluster to decide whether the cluster should be flipped. Unfortunately, the modified SW algorithm is not efficient. The problem is that large clusters of spins usually have a large net field acting on them and are prevented from flipping by these fields. An algorithm for Ising systems with fields that avoids this problem was introduced by Redner, Machta, and ChayesReMaCh ; ChMaRe98b . In this two-replica cluster algorithm large clusters are constructed from two replicas of the same system and have no net field acting on them so that they may be freely flipped. The two-replica cluster algorithm has been applied to study the phase transition of benzene adsorbed in zeolites DuMaSaAu and is more efficient than the conventional Metropolis algorithm for locating and simulating the critical point and the phase coexistence line. Combined with the replica exchange method of Swendsen and Wang SwWa86 , the two-replica method has been applied to the random field Ising model MaNeCh00 . The two-replica method is closely related to the geometric cluster Monte Carlo method DrKr ; HeBl96 ; HeBl98 .
In this paper, we report on a detailed investigation of the dynamics of the two-replica cluster (TRC) algorithm as applied to the two-dimensional Ising ferromagnetic in a staggered field (equivalently, the Ising antiferromagnet in a uniform field). The TRC algorithm introduced in Ref. ReMaCh, has two components that are not required for detailed balance and ergodicity. We studied the contribution to the performance of the algorithm of these optional components. We find that the complete TRC algorithm has a very small dynamic exponent $`z<0.4`$. However, we also find that this small value of $`z`$ requires one of the optional components and that this component depends on a special symmetry of Ising model in a staggered field. This observation leads to the question of whether cluster methods exist for efficiently simulating more general Ising models with fields. We investigated other optional components for the algorithm but these do not lead to acceleration when fields are present.
This paper is organized as follows. In Sec. II we introduce the Ising model in a staggered field and describe the TRC algorithm. In Sec. III we define the quantities to be measured and how errors are computed. In Sec. IV we present the results. The paper closes in Sec. V with a discussion.
## II The Model and Two-Replica Algorithm
### II.1 Ising Model in a Staggered Field
The Hamiltonian for the Ising model in a staggered field is
$$\beta [\sigma ]=K\underset{<i,j>}{}\sigma _i\sigma _j\underset{i}{}H_i\sigma _i$$
(1)
where the spin variables, $`\sigma _i`$ take the values $`\pm 1`$. $`K`$ is the coupling strength and $`H_i`$ is the magnetic field at site $`i`$. The summation in the first term of Eq. (1) is over nearest neighbors on an $`L\times L`$ square lattice with periodic boundary conditions and $`L`$ even. The second summation is over the sites of the lattice. The staggered field is obtained by setting $`H_i=H`$ if $`i`$ is in the even sublattice and $`H_i=H`$ if $`i`$ is in the odd sublattice. The staggered field breaks the up-down symmetry($`\sigma _i\sigma _i`$) of the zero field Ising model, however two symmetries remain. The Hamiltonian is invariant under even translations:
$$\sigma _{i+r_0}\sigma _i\text{ for all }i$$
(2)
with $`r_0`$ any vector in the even sublattice. The Hamiltonian is also invariant under odd translations together with a global flip:
$$\sigma _{i+r_1}\sigma _i\text{ for all }i$$
(3)
with $`r_1`$ any vector in the odd sublattice.
Figure 1 shows the line of critical points, $`K_c(H)`$ for this model. We carried out simulations at three points on the critical line taken from the high precision results of Ref. BlWu, ,
$$K_c(0)=0.4406867952$$
$$K_c(2)=0.7039642053$$
$$K_c(4)=1.1717153065$$
The basic idea of the two-replica cluster algorithm is to simultaneously simulate two independent Ising systems, $`\sigma `$ and $`\tau `$, on the same lattice and in the same field. Clusters of pairs of spins in this two-replica system are identified and flipped. In order to construct clusters, auxilliary bond variables are introduced. The bond variables {$`\eta _{ij}`$} are defined for each bond $`<i,j>`$ and take values 0 and 1. We say that $`<i,j>`$ is occupied if $`\eta _{ij}=1`$. A bond $`<i,j>`$ is satisfied if $`\sigma _i=\sigma _j`$ and $`\tau _i=\tau _j`$. Only satisfied bonds may be occupied.
The two-replica algorithm simulates a joint distribution of the Edwards-Sokal EdSo type for {$`\sigma _i`$} and {$`\tau _i`$}, and {$`\eta _{ij}`$}. The statistical weight $`X[\sigma ,\tau ,\eta ]`$ for the joint distribution is
$$X[\sigma ,\tau ,\eta ]=e^{G[\sigma ,\tau ]}\mathrm{\Delta }[\sigma ,\tau ,\eta ]B_p[\eta ]$$
(4)
where
$$G=K\underset{<i,j>}{}\sigma _i\tau _i\sigma _j\tau _j\underset{i}{}H_i(\sigma _i+\tau _i),$$
(5)
$`B`$ is the standard Bernoulli factor,
$$B_p[\eta ]=p^{|\eta |}(1p)^{N_b|\eta |}$$
(6)
$`|\eta |`$ = # $`\{<i,j>|\eta _{ij}=1\}`$ is the number of occupied bonds and $`N_b`$ is the total number of bonds of the lattice. The $`\mathrm{\Delta }`$ factor enforces the rule that only satisfied bonds are occupied: if for every bond $`<i,j>`$ such that $`\eta _{ij}=1`$ the spins agree in both replicas ($`\sigma _{ij}=\sigma _{ij}`$ and $`\tau _i=\tau _j`$) then $`\mathrm{\Delta }[\sigma ,\tau ,\eta ]=1`$; otherwise $`\mathrm{\Delta }[\sigma ,\tau ,\eta ]=0`$. It is straightforward to show that integrating $`X[\sigma ,\tau ,\eta ]`$ over the bond variables, $`\eta `$ yields the statistical weight for two independent Ising model in the same field,
$$e^{\beta [\sigma ]\beta [\tau ]}=const\underset{\{\eta \}}{}X[\sigma ,\tau ,\eta ]$$
(7)
if the identification is made that $`p=1e^{4K}`$.
### II.2 Two-Replica Cluster Algorithms
The idea of the two-replica cluster algorithm is to carry out moves on the spin and bond variables that satisfy detailed balance and are ergodic with respect to the joint distribution of Eq. (4). The occupied bonds $`\eta `$ define connected clusters of sites. We call site $`i`$ an *active site* if $`\sigma _i\tau _i`$ and clusters are composed either entirely of active or inactive sites. If a cluster of active sites is flipped so that $`\sigma \sigma `$ and $`\tau \tau `$ the factor $`G`$ is unchanged.
A single Monte Carlo sweep of the TRC algorithm is composed of the following three steps:
1. Occupy satisfied bond connecting active sites with probability $`p=1e^{4K}`$. Identify clusters of active sites connected by occupied bond (including single active sites). For each cluster $`k`$, randomly and independently assign a spin value $`s_k=\pm 1`$. If site $`i`$ is in cluster $`k`$ then the new spin values are $`\sigma _is_k`$ and $`\tau _is_k`$. In this way all active sites are updated.
2. Update each replica separately with one sweep of the Metropolis algorithm.
3. Translate the $`\tau `$ replica by a random amount relative to the $`\sigma `$ replica. If the translation is by an odd vector, all $`\tau `$ spins are flipped.
Step 1 of the TRC is similar to a sweep of the SW algorithm except that clusters are grown in a two-replica system rather than in a single replica and only active clusters are flipped. Note also that the bond occupation probability is $`p=1e^{4K}`$ for the TRC algorithm and $`p=1e^{2K}`$ for the SW algorithm. It is straightforward to show that Step 1 of the TRC algorithm satisfies detailed balance with respect to the joint distribution Eq. (4). Since only active sites participate in Step 1 of the algorithm, the Metropolis sweep, Step 2, is required for ergodicity. Step 3 contains the optional components of the algorithm: an even translation or an odd translation plus flip of one replica relative to the other. These moves are justified by the symmetries of the Ising model in a staggered field stated in Eqs. (2) and (3). When we refer to the TRC algorithm without further specification, we mean the algorithm described by the Steps 1-3 above. In the foregoing we also study the TRC with only even translations or with only odd translations.
In the TRC algorithm we flip only active clusters but it is also possible to flip inactive clusters if a weight factor associated with the change in $`G`$ is used. We call a flip of an active cluster to an active cluster ($`+`$ to $`+`$ or $`+`$ to $`+`$) an *active flip*. The TRC algorithm with inactive flips is obtained by replacing Step 1 with the following:
1. Occupy satisfied bonds with probability $`p=1e^{4K}`$. Identify clusters connected by occupied bonds (including single sites). For each cluster $`k`$, taken one at a time, randomly propose two new spin values values, $`s_k=\pm 1`$ and $`t_k=\pm 1`$ for the $`\sigma `$ and $`\tau `$ spins respectively. Compute $`\delta G`$, the change in $`G`$ that would occur if the spins in the $`k^{th}`$ cluster are changed to the proposed values leaving spins in other clusters fixed. If $`\delta G0`$ accept the proposed spin values (set $`\sigma _is_k`$ and $`\tau _it_k`$ for all sites $`i`$ in cluster $`k`$), otherwise, if $`\delta G>0`$ accept the proposed spin values with probability $`e^{\delta G}`$.
Step 1 is by itself ergodic however it may be useful to add Metropolis sweeps and translations.
## III Methods
We measured three observables using the TRC algorithm: the absolute value of the magnetization of a single replica, *m*; the energy of a single replica, $``$; and the absolute value of the net staggered magnetization for both replicas, *s*, where the definition of *s* is
$$\text{s}=|(\underset{iodd}{}\underset{ieven}{})(\sigma _i+\tau _i)|.$$
(8)
Note that the staggered magnetization is conserved by all components of the TRC algorithm except Metropolis sweeps and inactive flips. For each of these observables we computed expectation values of the integrated autocorrelation time, $`\tau _{int}`$ and the exponential autocorrelation time, $`\tau _{exp}`$. From $`\tau _{int}`$, we estimated the dynamic exponent $`z`$.
The autocorrelation function for $`\varphi `$, $`\mathrm{\Gamma }_{\varphi \varphi }(t)`$ is given by,
$$\mathrm{\Gamma }_{\varphi \varphi }(t)=\underset{l\mathrm{}}{lim}\frac{_{t^{}=1}^{lt}(\varphi (t^{})\widehat{\varphi })(\varphi (t^{}+t)\widehat{\varphi })}{_{t^{}=1}^l(\varphi (t^{})\widehat{\varphi })^2}.$$
(9)
The integrated autocorrelation time for observable $`\varphi `$ is defined by
$$\tau =\frac{1}{2}+\underset{t^{}\mathrm{}}{lim}\underset{t=1}{\overset{t^{}}{}}\mathrm{\Gamma }_{\varphi \varphi }(t)$$
(10)
and the exponential autocorrelation time for an observable $`\varphi `$ is defined by SaSo97
$$\tau _{exp,\varphi }=\underset{t\mathrm{}}{lim}\frac{|t|}{\mathrm{log}\mathrm{\Gamma }_{\varphi \varphi }(t)}.$$
(11)
In practice the limits in Eqs. (9), (10) and (11) must be evaluated at finite values. The length of the Monte Carlo runs determine $`l`$ and are discussed below. Following Ref. SaSo97, , we define
$$\tau _{int,\varphi }=\frac{1}{2}+\underset{t=1}{\overset{t^{}}{}}\mathrm{\Gamma }_{\varphi \varphi }(t)$$
(12)
and choose the cutoff $`t^{}`$ to be the smallest integer such that $`t^{}\kappa \tau _{int,\varphi }`$, where $`\kappa `$ = 6. We used the least-squares method to fit $`\mathrm{log}\mathrm{\Gamma }_{\varphi \varphi }(t/\tau _{int,\varphi })`$ as a function of $`t`$ to obtain the ratio of $`\tau _{int,\varphi }/\tau _{exp,\varphi }`$ and chose a cut-off at $`t/\tau _{int,\varphi }=5`$.
We used the blocking method NeBa99 ; SaSo97 to estimate errors. The whole sample of $`n`$ MC measurements was divided into $`m`$ blocks of equal length $`l=n/m`$. For each block $`i`$ and each measured quantity $`A`$, we computed the mean $`\widehat{A}_i`$ . Our estimates of $`\widehat{A}`$ and its error $`\delta A`$ are obtained from:
$$\widehat{A}=\frac{1}{m}\underset{i=1}{\overset{m}{}}\widehat{A}_i$$
(13)
$$\delta \widehat{A}^2=\frac{1}{m(m1)}\underset{i=1}{\overset{m}{}}(\widehat{A}\widehat{A}_i)^2$$
(14)
In our simulations, we divided the whole sample into $`m`$ blocks where $`m`$ is between 10 and 30.
For the data collected using the TRC algorithm, each block has a length $`l10^3\tau _{int}`$. For the data collected using modifications of the TRC algorithm, each block has a length $`l10^2\tau _{int}`$. Data were collected for $`H=0`$, 2 and 4 and for size $`L`$ in the range 16 to 256.
## IV Results
### IV.1 Integrated Autocorrelation Time
Table 1 gives the integrated autocorrelation time using the TRC algorithm for the magnetization, energy and staggered magnetization. Table 1 is comparable to the Table in Ref. ReMaCh, but the present numbers are systematically larger, especially at the larger system sizes. This discrepancy may be due to the sliding cut-off $`t^{}`$ used here instead of a fixed cut-off at 200 employed in Ref. ReMaCh, .
Table 2 gives the integrated autocorrelation times for magnetization using the TRC with only even or only odd translations. The comparison of TRC algorithm with only even translations and with only odd translations in Tables 2 shows that odd translations together with global flips of one replica relative to another are essential to achieve small and slowly growing autocorrelation times when the staggered field is present.
Table 3 shows the magnetization autocorrelation times using different algorithms for system size $`L=80`$. The Swendsen-Wang (SW) algorithm has the smallest $`\tau _{int,m}`$ in the absence of fields. However, when fields are present and the SW algorithm is then modified according to the method of Ref. DoSeTa, the performance is worse even than that of the Metropolis algorithm. The slow equilibration of the SW algorithm in the presence of the staggered field is due to small acceptance probabilities for flipping large clusters. On the other hand, the presence of staggered fields does not significantly change the performance the two-replica algorithm so long as odd translations are present. Inactive flips are helpful when there is no staggered field but when the staggered field is turned on, the autocorrelation time is not substantially improved by inactive flips. The explanation for the ineffectiveness of inactive flips when the staggered field is present is that the probability of accepting an inactive flip is small. For example, this probability is $`1.4\%`$ for $`L=80`$ and $`H=4`$.
The CPU time per spin on a Pentium III 450 MHz machine was also measured for the various algorithms and is listed in Table 3 for $`L=80`$ . By considering a range of system sizes we found that the CPU time for one MC sweep of the TRC algorithm increases nearly linearly with the number of spins. The TRC algorithm is a factor of 3 slower than the Metropolis algorithm but this difference is more than compensated for by system size $`80`$ by the much faster equilibration of the TRC algorithm. Even without odd translations, the TRC algorithm outperforms Metropolis for size 80.
### IV.2 Exponential Autocorrelation Time
The ratio of the integrated to exponential autocorrelation times was found to be nearly independent of system size over the range $`L=16`$ to $`L=256`$. We found that over this size range $`\tau _{int,m}/\tau _{exp,m}`$ varied from $`0.448\pm 0.008`$ to $`0.425\pm 0.008`$ for $`H=0`$; from $`0.44\pm 0.01`$ to $`0.43\pm 0.01`$ for $`H=2`$; and from $`0.448\pm 0.009`$ to $`0.409\pm 0.009`$ for $`H=4`$. The ratio $`\tau _{int,s}/\tau _{exp,s}`$ is also nearly independent of $`L`$ and $`H`$ and is about 0.45. The ratio $`\tau _{int,}/\tau _{exp,}`$ is nearly independent of $`L`$ but decreases slowly with $`H`$ ranging from 0.29 to 0.25 as $`H`$ ranges from 0 to 4. The almost constant $`\tau _{int,\varphi }/\tau _{exp,\varphi }`$ for different sizes suggests that the integrated and exponential autocorrelation times are governed by the same dynamic critical exponent.
### IV.3 Dynamic Exponent
Figures 2 and 3 show the magnetization integrated autocorrelation time for the TRC plotted on log-log and log-linear scales, respectively. Figures 4 and 5 show the energy integrated autocorrelation time for the TRC plotted on log-log and log-linear scales, respectively. Figures 6 and 7 show the staggered magnetization integrated autocorrelation time for the TRC plotted on log-log and log-linear scales, respectively.
For the whole range of $`L`$, logarithmic growth appears to give a somewhat better fit than a simple power law, particularly for the magnetization. Therefore, our results are consistent with $`z=0`$ for the TRC algorithm. Under the assumption that the dynamic exponent is not zero, we also carried out weighted least-squares fits to the form $`AL^z`$ and varied $`L_{min}`$, the minimum system size included in the fit. Figures 8, 9 and 10 show the dynamic exponent $`z`$ for the magnetization, energy and staggered magnetization, respectively, as a function of $`L_{min}`$ using the TRC algorithm. Figures 11 and 12, show the dynamic exponent as a function of $`L_{min}`$ for the magnetization for the TRC with only even translations and only odd translations, respectively. In all cases except $`z_{int,m,even}`$, the dynamic exponent is a decreasing function of $`L_{min}`$. For the magnetization, $`z_{int,m}`$ appears to extrapolate to a value between 0.1 and 0.2 as $`L_{min}\mathrm{}`$ while for the energy and staggered magnetization, the dynamic exponent appears to extrapolate to a value between 0.3 and 0.4. The small value of the dynamic exponent requires that odd translations and flips are included in the algorithm. From Fig. 11 it is clear that the dynamic exponent is near 2 for the TRC algorithm with only even translations.
Table 4 gives results of the weighted least squares fits for $`z`$ for the smallest values of $`L_{min}`$ for which there is a reasonable confidence level. Since there is a general downward curvature in the log-log graphs, these numbers are likely to be overestimates of the asymptotic values. Thus, we can conclude that the asymptotic dynamic exponent for the TRC algorithm is likely to be less than $`0.4`$ and is perhaps exactly zero. The dynamic exponent is apparently independent of the strength of the staggered field. For the case of the SW algorithm applied to the two-dimensional Ising with no staggered field the best estimate is $`z=0.25\pm 0.01`$BaCo ; CoBa but the results are also consistent with logarithmic growth of relaxation times. The numbers for dynamic exponent for the SW appear to be smaller than for the TRC algorithm but this difference may simply reflect larger corrections to scaling in the case of the TRC .
## V Discussion
We studied the dynamics of the two-replica cluster algorithm applied to the two-dimensional Ising model in a staggered field. We found that the dynamic exponent of the algorithm is either very small ($`z0.4`$) or zero ( $`\tau \mathrm{log}L`$) and that the dynamic exponent does not depend on the strength of the staggered field. A precise value of $`z`$ could not be determined because of large corrections to scaling. We tested the importance of various optional components of the algorithm and found that an odd translation and global flip of one replica relative to another is essential for achieving rapid equilibration. Without this component, $`z`$ is near 2 so there is no qualitative improvement over the Metropolis algorithm. An odd translation and global flip of one replica relative to the other allows for a large change of the total magnetization of the system with an acceptance fraction of $`100\%`$. Large changes in the global magnetization may also occur in the Swendsen-Wang algorithm in a field or via inactive flips in the TRC algorithm but these flips have a small acceptance fraction due to the staggered field. Unfortunately, the odd translation and flip move is allowed because of a special symmetry of the Ising model in a staggered field. For more general Ising systems with translationally invariant fields, we expect performance similar to the TRC with even translations only. In this case, the autocorrelation time is significantly less than for the Metropolis algorithm but the dynamic exponent is about the same. While the two-replica approach is useful for these more general problems of Ising systems with fields, it does not constitute a method that overcomes critical slowing down except when additional symmetries are present that allow one replica to be flipped relative to the other. Development of general methods for efficiently simulating critical spin systems with fields remains an open problem.
###### Acknowledgements.
This work was supported in part by NSF grants DMR 9978233.
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# Local Axisymmetric Instability Criterion in the Thin, Rotating, Multicomponent Disk.
## 1 Introduction
It was first shown by Safronov (1960) that a thin gaseous rotating disk can become unstable against local axisymmetric perturbations under the action of its own gravity. Quantitatively the instability criterion is expressed in terms of Toomre’s $`Q`$:
$$Q\frac{\kappa c}{\pi G\mathrm{\Sigma }},$$
(1)
with $`\kappa `$ being the epicycle frequency, $`c`$ being the speed of sound, and $`\mathrm{\Sigma }`$ being the surface density of the disk. Instability arises when $`Q<1`$.
Disks composed of stars behave very similarly to the gaseous disks for long wavelength because epicyclic motion of stars is the same as that of the gas. But collisionless disks of stars have no pressure to resist short wavelength density perturbations and thus the density waves in this case are different from simple sound waves present in the gaseous disks in this regime (Toomre 1964). When the wavelength becomes shorter than the typical amplitude of the star’s epicyclic motion the gravitational effect of the perturbation averages out and stars perform almost pure epicyclic motion with frequency $`\kappa `$. Instability can occur though for some intermediate wavelength as in the case of the gaseous disks. Toomre (1964) had shown that in the stellar disks self-gravity can support exponentially growing axisymmetric modes when the condition $`Q<1`$ is satisfied with $`\pi `$ replaced by $`3.36`$ and $`c`$ replaced by the stellar velocity dispersion $`\sigma `$ in the radial direction in definition (1). Except for the qualitative difference in behavior at small wavelength, stellar and gaseous disks behave almost identically for large wavelength which leads to a close similarity of the corresponding instability conditions.
However, real disks of spiral galaxies do not consist of only one component which could be characterized by a single value of the velocity dispersion. Indeed, a significant part of the ISM in galactic disks is usually concentrated in the form of cold gas with characteristic sound velocity $`510`$ km s<sup>-1</sup>. At the same time stars usually have larger velocity dispersions $`2035`$ km s<sup>-1</sup>.
The problem of the treatment of such two-component disks was first addressed by Jog & Solomon (1984a, JS). They assumed galactic disk to consist of two fluid components: one hot, playing the role of the stellar population, and one cold representing the usual ISM, and they neglected the fact that in real galaxies hot component is collisionless and provides no damping. Though this could be important shortcoming of the analysis for the short wavelength it turns out to be a good approximation for the stability purposes as we will see later in this paper. Jog & Solomon (JS) have demonstrated that even relatively small amounts of the cold gas ($`10\%`$ by mass) can very effectively destabilize the whole system.
The two-fluid stability criterion was applied to a variety of disks by Jog & Solomon (1984b), Elmegreen (1995), and Jog (1996). They showed that our Galaxy is stable against local axisymmetric perturbations to the extent of the accuracy of the parameters and mass distribution models assumed. Theory can be further refined by including the effects of the finite thickness of the disk on its stability (Toomre 1964) and by taking into account the vertical motions in the disk (Romeo 1992). It was shown that the instability condition remains the same provided the unperturbed surface density is multiplied by an appropriate reduction factor.
Real galaxies always contain more than two isothermal components. It is now widely believed that ISM in the galactic disk is subdivided into a number of components with different temperatures and, thus, different dynamical properties (McKee & Ostriker 1977). The same is true about the stellar population. Stars in the galactic disk are being constantly scattered by giant molecular clouds and transient spiral arms which steadily increase their velocity dispersions (Spitzer & Schwarzschild 1951, 1953; Barbanis & Woltjer 1967; Carlberg & Sellwood 1985). It means that the stars of different ages have different velocity dispersions and thus dynamically should be treated separately.
Recent results of the $`Hipparcos`$ satellite (ESA 1997) provided us with a wealth of information about the local stellar kinematics. The proper motions of different star populations were accurately measured and the resulting velocity ellipsoids were constructed for them enabling one to study different kinematic constituents of our Galaxy separately (Dehnen & Binney 1998; Mignard 2000). This mission also gave us a lot of information about distances to the stars which provided us with the local densities of various types of the stars (Holmberg & Flynn 2000).
In this paper we present the analysis of the axisymmetric gravitational stability of the thin rotating disk which consists of a number of components, each of them being characterized by its own temperature and surface density. In doing so we distinguish between the two types of components: collisional and thus having pressure forces such as usual gas, or collisionless such as stellar component, which needs a kinetic treatment. The derivation of the dispersion relation is presented in §2 and 3. In §4 we compare our stability condition with the one derived in JS to test the difference arising when one of the disk components is treated as collisionless which better represents real galaxies. Finally, in §5 we apply our results to the case of our Galaxy in the Solar neighborhood using recent data from $`Hipparcos`$ satellite.
## 2 Basic equations
We work in non-rotating cylindrical system, $`r,\varphi ,z`$, such that $`z`$-axis coincides with the rotation axis of the disk and angle $`\varphi `$ increases in the direction of rotation. We start first with the equations for the fluid components and then make some correction to take collisionless components into account. We will usually use index $`i`$ to describe gas components and index $`j`$ for stellar components.
All the components are spatially concentrated in the thin disk and the effects of the finite disk thickness will usually be disregarded. All the motions are assumed to occur only in the plane of the disk.
We suppose that there are $`n_g`$ gaseous components contributing to the mass of the Galaxy, each characterized by the sound velocity $`c_i`$ and surface density $`\mathrm{\Sigma }_{gi}`$, and $`n_s`$ collisionless components with $`\mathrm{\Sigma }_{sj}`$ being the surface density of the $`j`$-th component. Each collisionless component in the unperturbed state is assumed to have a Schwarzschild distribution function, that is for $`j`$-th component
$`f(v_r,v_\varphi )={\displaystyle \frac{\mathrm{\Sigma }_{sj}}{2\pi \sigma _{rj}\sigma _{\varphi j}}}\mathrm{exp}\left\{{\displaystyle \frac{v_r^2}{2\sigma _{rj}^2}}{\displaystyle \frac{[v_\varphi v_c]^2}{2\sigma _{\varphi j}^2}}\right\},`$ (2)
where $`v_c`$ is the circular velocity at current distance from the galactic center and $`\sigma _{rj}`$ and $`\sigma _{\varphi j}`$ are the velocity dispersions in $`r`$ and $`\varphi `$ directions correspondingly. These velocity dispersions are related by $`\sigma _\varphi ^2/\sigma _r^2=4B^2/\kappa ^2`$, where $`B`$ is Oort’s B constant (Binney & Tremaine 1987). We will later use simply $`\sigma `$ to denote the velocity dispersion in the r-direction $`\sigma _r`$.
For the description of collisional components we assume usual hydrodynamical description with isotropic pressure. To describe the star-like components we use a different approach which is rooted in the kinetic treatment of the collisionless systems as described elsewhere (Toomre 1964; Binney & Tremaine 1987). We will see in §4 that this accurate treatment shows that fluid approach really does quite a good job in describing the stability of the collisionless systems with multiple components as it does in one-fluid case, though some quantitative differences exist.
Equations governing the motion of the gas components in our coordinates are Euler’s and continuity equations:
$`{\displaystyle \frac{𝐯_i}{t}}+(𝐯_i)𝐯_i={\displaystyle \frac{1}{\mathrm{\Sigma }_i}}P_i\mathrm{\Phi },`$ (3)
$`{\displaystyle \frac{\mathrm{\Sigma }_i}{t}}+(\mathrm{\Sigma }_i𝐯_i)=0,`$ (4)
for $`i`$-th component.
We assume that pressure of each component $`P_i=K\mathrm{\Sigma }_i^\gamma `$ and introduce specific enthalpy
$$h_i=\frac{\gamma }{\gamma 1}K\mathrm{\Sigma }_i^{\gamma 1}.$$
(5)
We linearize equations (3) and (4) by assuming that $`v_{ri}=u_i,v_{\varphi i}=v_c+v_i,h_i=h_{0i}+h_{1i},\mathrm{\Sigma }_i=\mathrm{\Sigma }_{0i}+\mathrm{\Sigma }_{1i},`$ and $`\mathrm{\Phi }=\mathrm{\Phi }_0+\mathrm{\Phi }_1`$.
Then equations (3) and (4) reduce in the first order to (Binney & Tremaine 1987)
$`{\displaystyle \frac{u_i}{t}}+\mathrm{\Omega }{\displaystyle \frac{u_i}{\varphi }}2\mathrm{\Omega }v_i={\displaystyle \frac{}{r}}(\mathrm{\Phi }_1+h_{1i}),`$ (6)
$`{\displaystyle \frac{v_i}{t}}+\mathrm{\Omega }{\displaystyle \frac{v_i}{\varphi }}2Bu_i={\displaystyle \frac{1}{r}}{\displaystyle \frac{}{\varphi }}(\mathrm{\Phi }_1+h_{1i}),`$ (7)
$`{\displaystyle \frac{\mathrm{\Sigma }_{1i}}{t}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}(r\mathrm{\Sigma }_{i0}u_i)+\mathrm{\Omega }{\displaystyle \frac{\mathrm{\Sigma }_i}{\varphi }}+{\displaystyle \frac{\mathrm{\Sigma }_{0i}}{r}}{\displaystyle \frac{v_i}{\varphi }}=0,`$ (8)
where
$$B=\frac{1}{2}\left[\mathrm{\Omega }+\frac{(\mathrm{\Omega }r)}{r}\right]$$
(9)
is the Oort’s $`B`$ constant.
Since here we are interested in axisymmetric perturbations only it is always supposed that $`/\varphi =0`$ in our analysis. We assume that all the first order dependent quantities vary like $`\mathrm{exp}[i(kr\omega t)]`$, where $`\omega `$ is the angular frequency and $`k=2\pi /\lambda `$ is the wavenumber. For the local analysis we employ the WKB approximation (or tight-winding approximation) which requires that $`kr1`$ and allows us to neglect terms proportional to $`1/r`$ compared to the terms proportional to $`k`$. With all these simplifications equations (6)-(8) reduce to
$`i\omega u_i2\mathrm{\Omega }v_i=ik(\mathrm{\Phi }_1+h_{1i}),`$ (10)
$`i\omega v_i2Bu_i=0,`$ (11)
$`i\omega \mathrm{\Sigma }_{1i}+ik\mathrm{\Sigma }_{0i}u_i=0,`$ (12)
for $`i`$-th gas component.
These equations have to be supplemented with Poisson equation
$$^2\mathrm{\Phi }=4\pi G\left(\underset{i=1}{\overset{n_g}{}}\mathrm{\Sigma }_{1i}+\underset{j=1}{\overset{n_s}{}}\mathrm{\Sigma }_j\right)\delta (z),$$
(13)
where $`\delta (z)`$ is the Dirac delta function arising from our assumption of infinitely thin disk. With $`\mathrm{\Sigma }_{1i,j}`$ and $`\mathrm{\Phi }_1`$ being proportional to $`\mathrm{exp}[i(kr\omega t)]`$ the relation between the perturbations of the potential and surface densities becomes (Toomre 1964)
$$\mathrm{\Phi }_1=\frac{2\pi G}{k}\left(\underset{i=1}{\overset{n_g}{}}\mathrm{\Sigma }_{1i}+\underset{j=1}{\overset{n_s}{}}\mathrm{\Sigma }_{1j}\right).$$
(14)
The sound speed for each gas component is defined as $`c_j^2=dP_j/d\mathrm{\Sigma }_{0j}`$, and in this case the perturbation of enthalpy reduces to
$$h_{1j}=c_j^2\frac{\mathrm{\Sigma }_{1j}}{\mathrm{\Sigma }_{0j}}.$$
(15)
## 3 Dispersion relation and stability criterion
Equations (10)-(12), (14), and (15) form a closed set of linear equations and we now solve them to get the dispersion relation.
From equations (10) and (11) we can relate perturbation of the radial velocity of each component $`u_i`$ to the potential perturbation $`\mathrm{\Phi }_1`$ using equation (15):
$$u_i=\frac{\omega k}{\mathrm{\Delta }}\left(\mathrm{\Phi }_1+c_j^2\frac{\mathrm{\Sigma }_{1j}}{\mathrm{\Sigma }_{0j}}\right),$$
(16)
where
$$\mathrm{\Delta }=\kappa ^2\omega ^2,$$
(17)
and $`\kappa ^2=4\mathrm{\Omega }B`$ is the epicyclic frequency.
Now we are in position to treat the star-like components. Since they are supposed to be collisionless they cannot create any pressure and it means that for $`j`$-th stellar component sound velocity $`c_j=0`$. But because of the epicyclic motion at any given point in the disk there are stars from different parts of the perturbed structure and it leads to important cancellation effects. They were first calculated by Toomre (1964) and it can be shown that each stellar component is described by the Jeans equations which are pretty similar to the usual hydrodynamic equations with important difference being that instead of equation (16) we have
$$u_k=\frac{\omega k}{\mathrm{\Delta }}\mathrm{\Phi }_1(\frac{\omega }{\kappa },\frac{k^2\sigma _j^2}{\kappa ^2}),$$
(18)
where the reduction factor $`F`$ is given by the expression (Binney & Tremaine 1987)
$`(s,\chi )={\displaystyle \frac{2}{\chi }}(1s^2)e^\chi {\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{I_n(\chi )}{1s^2/n^2}},`$ (19)
and $`I_n`$ are the Bessel functions of order $`n`$.
Now, from (12), (16) and (18) we can eliminate $`u_i`$ to get
$$\mathrm{\Sigma }_{1i}=\frac{k^2\mathrm{\Sigma }_{0i}}{\mathrm{\Delta }+k^2c_i^2}\mathrm{\Phi }_1,$$
(20)
for $`i`$-th gas component and
$$\mathrm{\Sigma }_{1j}=\frac{k^2\mathrm{\Sigma }_{0j}}{\mathrm{\Delta }}\mathrm{\Phi }_1_j,$$
(21)
for $`j`$-th stellar component and $`_j=(\omega /\kappa ,k^2\sigma _j^2/\kappa ^2)`$.
We can then substitute (21) and (20) into (14) to obtain finally the desired dispersion relation:
$`2\pi Gk{\displaystyle \underset{i=1}{\overset{n_g}{}}}{\displaystyle \frac{\mathrm{\Sigma }_{0i}}{\kappa ^2+k^2c_j^2\omega ^2}}+2\pi Gk{\displaystyle \underset{j=1}{\overset{n_s}{}}}{\displaystyle \frac{\mathrm{\Sigma }_{0j}_j}{\kappa ^2\omega ^2}}=1.`$ (22)
If only $`2`$ fluid components are present in the disk this dispersion relation clearly reduces to the one derived by JS. If one of the components is collisionless the dispersion relation is different and it was first considered by Romeo (1992). We discuss the difference between the two cases in §4.
The multicomponent disk is unstable if $`\omega ^2<0`$ because in this case oscillatory behavior of the density waves changes to an exponential growth. The dispersion relation (22) has an infinite number of solutions with $`\omega ^2(k)>0`$ for all $`k`$, which are clearly stable. But there is also one mode of oscillations in the system which has a single solution with $`\mathrm{}<\omega ^2/\kappa ^2<1`$ and this mode could be unstable.
To find the condition for instability we note that as $`\omega ^2\mathrm{}`$ at a fixed wavenumber $`k`$ LHS of the equation (22) tends to zero, and it monotonically increases to $`+\mathrm{}`$ as $`\omega ^2\kappa ^2`$. This property of a steady growth is important because if (22) has no solution for negative $`\omega ^2`$ it implies that the LHS of (22) at $`\omega ^2=0`$ must be less than $`1`$ and vice versa. So the axisymmetric instability arises in multicomponent disk if and only if
$`2\pi Gk{\displaystyle \underset{i=1}{\overset{n_g}{}}}{\displaystyle \frac{\mathrm{\Sigma }_{0i}}{\kappa ^2+k^2c_j^2}}+2{\displaystyle \frac{\pi Gk}{\kappa ^2}}{\displaystyle \underset{j=1}{\overset{n_s}{}}}\mathrm{\Sigma }_{0j}\mathrm{\Psi }_j>1,`$ (23)
where
$`\mathrm{\Psi }_j=(0,{\displaystyle \frac{k^2\sigma _j^2}{\kappa ^2}})={\displaystyle \frac{2\kappa ^2}{k^2\sigma _j^2}}e^{k^2\sigma _j^2/\kappa ^2}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}I_n\left({\displaystyle \frac{k^2\sigma _j^2}{\kappa ^2}}\right).`$ (24)
Since
$$\underset{n=1}{\overset{\mathrm{}}{}}I_n(\chi )=\frac{1}{2}\left[e^\chi I_0(\chi )\right],$$
(25)
(Abramowitz & Stegun 1964) we can rewrite (24) as
$`\mathrm{\Psi }_j={\displaystyle \frac{1e^{\chi _j}I_0(\chi _j)}{\chi _j}}={\displaystyle \frac{1e^{k^2\sigma _j^2/\kappa ^2}I_0(k^2\sigma _j^2/\kappa ^2)}{k^2\sigma _j^2/\kappa ^2}}.`$ (26)
Suppose again that there are only two fluid (collisional) components in the disk then the instability condition (23) reduces to the one derived in JS as it should.
## 4 Realistic two-fluid case
Now we consider the stability in the case when we have only two fluids present in the disk in more details. Our main purpose here is to see how the fact that one of the components is in reality collisionless changes the overall dynamics of the system from the case considered by JS.
In the case of two components, one stellar and one gaseous, we assume $`\mathrm{\Sigma }_g`$ and $`\mathrm{\Sigma }_s`$ to be the unperturbed gas and stellar surface densities, $`c_g`$ the gas sound speed, and $`\sigma _s`$ the stellar velocity dispersion in $`r`$-direction and introduce the following dimensionless quantities
$`Q_s={\displaystyle \frac{\kappa \sigma _s}{\pi G\mathrm{\Sigma }_s}},Q_g={\displaystyle \frac{\kappa c_g}{\pi G\mathrm{\Sigma }_g}},`$ (27)
$`q=k\sigma _s/\kappa ,R=c_g/\sigma _s.`$ (28)
Here $`Q_g`$ is a Toomre’s $`Q`$ parameter for gas, while $`Q_s`$ is different from the Toomre’s $`Q`$ parameter for the collisionless system, which is given by $`\kappa \sigma _s/(3.36G\mathrm{\Sigma }_s)`$ (Toomre 1964). It means that with our definition of $`Q_s`$ the one–component stellar disk is gravitationally unstable when $`Q_s<3.36/\pi =1.07`$.
With these definitions our instability condition (LABEL:2fl) reduces to
$`{\displaystyle \frac{2}{Q_s}}{\displaystyle \frac{1}{q}}\left[1e^{q^2}I_0(q^2)\right]+{\displaystyle \frac{2}{Q_g}}R{\displaystyle \frac{q}{1+q^2R^2}}>1.`$ (29)
If one follows the Jog & Solomon approach and treats stellar component as a fluid with sound speed equal to $`\sigma _s`$, one gets the following instability condition in terms of our dimensionless variables
$`{\displaystyle \frac{2}{Q_s}}{\displaystyle \frac{q}{1+q^2}}+{\displaystyle \frac{2}{Q_g}}R{\displaystyle \frac{q}{1+q^2R^2}}>1.`$ (30)
In Figure 1 we show the dependences of the $`\omega ^2`$ upon $`k`$ for some particular choices of parameters and compare the curves in the case of gas-gas mix considered by Jog & Solomon (labeled as G$`+`$G) and star-gas mix (labeled as S$`+`$G), calculated using dispersion relation (22) for the case of two-component disk. If one adopts $`\kappa =36`$ km s<sup>-1</sup> kpc<sup>-1</sup> to characterize the disk rotation and sound speed of the gas $`c=5`$ km s<sup>-1</sup>, then the model corresponding to the Figure 1a has $`\mathrm{\Sigma }_s=45M_{}`$ pc<sup>-2</sup>, $`\mathrm{\Sigma }_g=9M_{}`$ pc<sup>-2</sup>, and $`\sigma =25`$ km s<sup>-1</sup>, which is a good representation of the typical spiral galaxy like our own.
Second plate of the same Figure correspond to the model with $`\mathrm{\Sigma }_s=17M_{}`$ pc<sup>-2</sup>, $`\mathrm{\Sigma }_g=11M_{}`$ pc<sup>-2</sup>, and $`\sigma =12.5`$ km s<sup>-1</sup> for the same choice of $`\kappa `$ and $`c`$. This might be thought of as a central part of a relatively young spiral and it turns out to be unstable with both instability criteria.
One can immediately see that in general the difference between the curves given by conditions (30) and (29) is small for low enough wavenumbers as it should be because pressure forces are negligible in this regime and it is only the epicyclic motions and the self-gravity which determine the dynamics of the disk. However, as $`k`$ grows, in gas-gas case the frequency describing the oscillations of the whole system grows due to the pressure forces in the gas, while in the star-gas case this does not happen. Despite the presence of the gas in the system it is really stars which determine the natural frequency of oscillations of the system in this case and as $`k\mathrm{}`$ frequency becomes constant – $`\omega \kappa `$ – and it presents a dramatic difference in the disk response to the axisymmetric density perturbations at small wavelengths compared to the fluid-fluid case.
In the two-fluid case JS found that the ratio of the perturbation of the gas surface density to that of the stellar surface density is always a growing function of wavenumber $`k`$ and for large $`k`$ this ratio becomes very large. This is not the case when one of the disk components is collisionless. Indeed, in this case $`\omega \kappa `$ for large $`k`$. Using (20) and (21) we get
$`{\displaystyle \frac{\mathrm{\Sigma }_{1g}}{\mathrm{\Sigma }_{1s}}}{\displaystyle \frac{\mathrm{\Sigma }_g}{\mathrm{\Sigma }_s}}{\displaystyle \frac{\sigma ^2}{c^2}}\left[2e^{q^2}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{I_n(q^2)}{1\omega ^2/(\kappa ^2n^2)}}\right]^10,`$ (31)
because the sum in (31) diverges as $`\omega \kappa `$. So, in realistic gas-stellar case this ratio goes to $`0`$ rather than increases. It is illustrated in Figure 2 where we have shown the ratio of relative perturbed surface densities for two cases: fluid-fluid and fluid-stellar with the same set of parameters. This decrease in the gas density perturbation with growing $`k`$ leads to the dominant role of stellar component in the dynamics of the whole system for large $`k`$, which is different from the JS case.
The relative contribution to the instability per unit surface density of the component considered by JS is given by the ratio of the corresponding terms in the LHS of (30)
$`\gamma ={\displaystyle \frac{\mathrm{\Sigma }_g}{\mathrm{\Sigma }_s}}{\displaystyle \frac{1+q^2}{1+q^2R^2}}`$ (32)
and is always larger than $`\mathrm{\Sigma }_g/\mathrm{\Sigma }_s`$ for $`R<1`$. In the star-fluid case one gets from (29) that
$`\gamma ={\displaystyle \frac{\mathrm{\Sigma }_g}{\mathrm{\Sigma }_s}}{\displaystyle \frac{q^2}{(1+q^2R^2)\left(1\mathrm{exp}(q^2)I_0(q^2)\right)}}.`$ (33)
Careful study shows that for $`\sqrt{3}/2<R<1`$ there is always a range of small $`q`$ (and $`k`$) in which $`\gamma <\mathrm{\Sigma }_g/\mathrm{\Sigma }_s`$. Of course, as $`q`$ grows $`\gamma `$ becomes greater than $`\mathrm{\Sigma }_g/\mathrm{\Sigma }_s`$ for $`R<1`$ ($`\gamma R^2\mathrm{\Sigma }_g/\mathrm{\Sigma }_s`$ as $`q\mathrm{}`$ ), but for $`q1`$ the relative contribution to the instability per unit surface density is greater for stars if $`\sqrt{3}/2<R<1`$. Even in this case though gas contributes only several per cent smaller than the stars and in most astrophysically interesting cases cold material still has a serious impact on the stability of the system.
Each of the criteria (29) and (30) produces some region in parameter space of the disk models in which they are stable against local axisymmetric gravitational perturbations. Unfortunately, it does not seem trivial to construct an effective analytical way of defining some effective value $`Q_{eff}`$ as a function of $`Q_g`$, $`Q_s`$ and other disk parameters so that the system is stable when $`Q_{eff}>1`$, as it was done for the fluid-fluid case by Elmegreen (1995). Instead we follow the approach of Jog (1996) and study the problem seminumerically. We parametrize our models here by $`1/Q_s,1/Q_g`$ and $`R`$. In Figure 3 we compare the stable regions produced by each of the instability criteria.
In the JS case the marginal stability curves given by conditions $`\omega ^2/k=0`$ and $`\omega ^2=0`$ are symmetric with respect to the line $`Q_s=Q_g`$. Indeed, if these conditions are fulfilled and an inequality in (30) changes to equality for $`Q_s=Q_1,Q_g=Q_2`$ and $`q=q_{crit}`$, then one can easily check that this is also true for $`Q_s=Q_2,Q_g=Q_1`$ and $`q=q_{crit}/R`$ provided that the instability condition is given by (30). Of course, this is not the case for the star-gas disk because any such symmetry is absent in the relation (29).
As we increase $`R`$ from $`0`$ to $`1`$ the region of the parameter space occupied by stable models shrinks until $`R=1`$. The further increase of $`R`$ beyond $`1`$ causes reexpansion of the region occupied by stable models. In fact in the fluid-fluid case the marginal curve corresponding to some particular $`R=R_0`$ coincides with the curve corresponding to $`R=1/R_0`$ which can be directly checked using (30). But these models are likely to be uncommon since they have $`c>\sigma `$ which seems to be unusual in real galaxies.
The stable region in the star-gas case is in general smaller than that in the gas-gas case. The difference is especially noticeable at $`1/Q_s1`$ and low $`1/Q_g`$, when the gas influence on the dynamics of the system is smallest and the stability condition is close to the one-fluid stellar stability criterion which is somewhat different from the fluid case. Nevertheless, the difference between two cases is quite small for most of the parameters (especially for small $`Q_s`$) and one can usually use the JS stability criterion in these cases.
## 5 Application to the Solar neighborhood
In this section we apply the results derived in §3 to the neighborhood of the Sun in the Galaxy.
In doing this one should realize that all the stars and gas in the Galaxy cannot be simply put into two distinct groups with some well defined velocity dispersion for stars and sound speed for gas. The reason for that is that the stars of different ages have different velocity dispersions - the older the star the larger its random motion. This random heating of stellar population is produced by the scattering of the stars by the giant molecular clouds (Spitzer & Schwarzschild 1951, 1953) and/or transient spiral density waves in the Galaxy (Barbanis & Woltjer 1967; Carlberg & Sellwood 1985).
Gas in its turn intrinsically has a multicomponent nature caused by the constant energy input from supernovae explosions and various cooling and heating processes determining its thermal equilibrium (McKee & Ostriker 1977; Kulkarni & Heiles 1987). Recent studies have enabled us to distinguish 5 phases of the ISM: molecular gas in the form of the clouds (Scoville & Sanders 1987), cold neutral medium (CM) also in the form of clouds, and 3 more or less uniformly distributed gaseous components: warm neutral medium (WNM), warm ionized medium (WIM) (Kulkarni & Heiles 1987), and hot component (Savage 1987). They have different sound speeds and surface densities which makes one treat them separately.
Unperturbed gravitational field of our Galaxy is not axisymmetric and it limits our analysis to some extent. It was also shown that the stellar distribution function can not always be represented by formula (2) because nonaxisymmetries of the Galactic gravitational field cause vertex deviations of the velocity ellipsoid (Dehnen & Binney 1998). Young stars of O and B types are likely not to have Schwarzschild distribution because they have had no time to be sufficiently scattered by the giant molecular clouds or transient spiral arms, but they probably are not important mass contributors in the local part of the Galaxy. The local approximation itself maybe questionable because most unstable waves have $`\lambda `$ of the order of several kpc.
Other possible complications in real galactic disks involve the presence of cosmic rays and magnetic fields which could influence the gas dynamics (Elmegreen 1987). In this paper we are primarily interested in the purely gravitational aspects of the disk instability and for this reason we neglect them at all, though it makes our consideration less realistic when applied to the galaxies. For these reasons our analysis of the stability of the Solar neighborhood should be considered only as mostly illustrative though bearing sufficient resemblance to reality.
Following Holmberg & Flynn (2000) we split the galactic disk mass between $`13`$ major parts: $`4`$ gaseous and $`9`$ stellar. We neglected hot component of the gas because of its low number density, $`n0.003`$ cm<sup>-3</sup>, high temperature, $`T10^6`$ K, (Savage 1987) and, consequently, large thickness which makes the reduction effects very important (Romeo 1992). For the same reason we neglect the stellar halo component. Parameters of the gaseous components are taken from Kulkarni & Heiles (1987) and Scoville & Sanders (1987).
Surface densities of stellar components are taken from Holmberg & Flynn (2000). We got the radial velocity dispersions based on the recent data from $`Hipparcos`$ from Mignard (2000). Velocity dispersions of white and brown dwarfs have been chosen quite arbitrarily. All the parameters assumed for the mass constituents are listed in Table 1.
| i | Component | $`\mathrm{\Sigma }_i`$ | $`\sigma _{ri}`$ or $`c_s`$ |
| --- | --- | --- | --- |
| | | ($`M_{}\mathrm{pc}^2`$) | ($`\mathrm{km}\mathrm{s}^1`$) |
| $`1`$ | $`\mathrm{H}_2`$ | $`3.0`$ | $`4.0`$ |
| $`2`$ | CM | $`4.0`$ | $`6.9`$ |
| $`3`$ | WNM | $`4.0`$ | $`9.0`$ |
| $`4`$ | WIM | $`2.0`$ | $`9.0`$ |
| $`5`$ | giants | $`0.4`$ | $`26.0`$ |
| $`6`$ | $`M_V<2.5`$ | $`0.9`$ | $`17.0`$ |
| $`7`$ | $`2.5<M_V<3.0`$ | $`0.6`$ | $`20.0`$ |
| $`8`$ | $`3.0<M_V<4.0`$ | $`1.1`$ | $`22.5`$ |
| $`9`$ | $`4.0<M_V<5.0`$ | $`2.0`$ | $`26.0`$ |
| $`10`$ | $`5.0<M_V<8.0`$ | $`6.5`$ | $`30.5`$ |
| $`11`$ | $`M_V>8.0`$ | $`12.3`$ | $`32.5`$ |
| $`12`$ | white dwarfs | $`4.4`$ | $`32.5`$ |
| $`13`$ | brown dwarfs | $`6.2`$ | $`32.5`$ |
In Figure 4 we show the dependence of $`\omega ^2`$ upon the inverse radial wavelength $`k/2\pi `$. Epicyclic frequency $`\kappa =36`$ km s<sup>-1</sup> kpc<sup>-1</sup> is assumed throughout the calculation. Curve labeled $`1`$ corresponds to the choice of parameters described in the Table 1. One can see that for all radial wavelength $`\omega ^2`$ is positive that is the whole disk system is stable against local axisymmetric perturbations. To check how rigorous this conclusion is we varied some of the model parameters until the disk became unstable.
Data about the surface densities and velocity dispersions of the brown dwarfs (BDs) and white dwarfs (WDs) seem to be the most uncertain among all the model parameters, so we tried to vary them first. Curve labeled $`2`$ has all the surface densities as listed in Table 1 but the velocity dispersion of WDs and BDs was lowered to $`\sigma _r=19.2`$ km s<sup>-1</sup>. Only if this population is so cold can it make the system unstable with all other parameters being kept unchanged. It seems inevitable that real WD and BD populations must be sufficiently hotter because only young stars can have such a low velocity dispersion (Mignard 2000).
Curve labeled $`3`$ shows $`\omega ^2\lambda ^1`$ dependence for the case when total surface density of WDs and BDs was raised from $`10.6M_{}`$ pc<sup>-2</sup> to $`26M_{}`$ pc<sup>-2</sup> keeping the rest of model parameters unchanged. This leads to the neutral stability of the system but such a surface density seems to be too large despite large uncertainties and claims of some authors that such extreme values of surface density could be common. For example, Festin (1998) found high mass density of the WD, $`2.6`$ times larger than we assume here, but his conclusions were based on a small sample of $`7`$ sources only. Other authors (Ruiz & Takamiya 1995; Oswalt et al. 1996) claim values for WD mass density which are in agreement with what we take. Even more uncertainty is involved in determining the density of the BD. Different surveys quote values from $`0.6`$ to $`4`$ times what we assume in this research (Fuchs, Jahreiss, & Flynn 1998). Recent data (Reid et al. 1999) based on a large enough sample imply mass density of BD $`0.005M_{}`$ pc<sup>-3</sup> which is about $`60\%`$ of what we assume in our calculations.
Finally, the fourth curve shows the dispersion relation for the disk with the lowered sound speed in some of the gas components: in CM we set $`c_s=5.0`$ km s<sup>-1</sup>, in WNM $`c_s=7.5`$ km s<sup>-1</sup>, and in WIM $`c_s=8.0`$ km s<sup>-1</sup>, with all other parameters unchanged. In this case disk becomes marginally unstable. Important thing to notice here is that small variations in the gas sound speed can have stronger influence on the disk stability than large changes in the stellar velocity dispersion, even though the surface density of the gas is smaller than that of stars. This is a manifestation of the crucial importance of cold material for the stability of the whole disk, which was first noted by Jog & Solomon (1984a). It is also easy to see that variations of the gas parameters produce significant change of the most unstable wavelength compared to the variations of the stellar parameters; it is reduced from $`2`$ kpc to $`1`$ kpc by that small change in gas sound speeds.
The bottom line is that local Galactic disk seems to be stable against local gravitational axisymmetric perturbation even when allowance for a scant knowledge of some of the Galactic parameters is made, which confirms results of Jog & Solomon (1984b) and Elmegreen (1995) for two-fluid disks.
## 6 Conclusions
Due to the continuous interaction with the transient spiral structure and giant molecular clouds stars diffuse in the velocity space towards higher random velocities as their age increases and it was confirmed observationally (Dehnen & Binney 1998; Mignard 2000). It raises a necessity of considering the dynamics of the Solar neighborhood taking into account its complex multicomponent structure. In this paper we studied the stability of such a system against gravitational axisymmetric perturbations in the tight-winding limit. It is possible to derive an analytic dispersion relation characterizing multicomponent thin differentially rotating disk and study its stability.
In doing so we distinguished between two types of the disk constituents: stellar and gaseous. Stellar population is dynamically different from fluid because stars form collisionless system (Binney & Tremaine 1987) while gas must be treated as a fluid. We demonstrated that the difference in the results for stability produced by two approaches is small for multicomponent disks in many astrophysically interesting cases. Some disk models though could be sensitive to the choice of the stability condition and in that case one should use correct criterion given by the equation (23).
We apply our results to the stability of the Solar neighborhood and confirm the conclusions of the previous two-fluid studies that local Galactic disk is stable against axisymmetric perturbations in the WKB limit, even taking into account uncertainties associated with determining some of the disk parameters.
Keeping in mind previous two-fluid results (JS) it is not surprising that relatively small variations of the gaseous component parameters are very important for the overall disk stability. Indeed, Figure 4 shows that decrease of the sound speed of the gas by about $`1`$ km s<sup>-1</sup> drives instability of the disk, while in the case of stellar component one needs to reduce its velocity dispersion by $`10`$ km s<sup>-1</sup> to produce the same outcome. Even though the mass of the gas in the disk is smaller that the stellar mass its small random motion makes it much more susceptible to its own self-gravity than the hot stellar component.
Our study of the stability of the Solar neighborhood neglects a lot of physics such as magnetic fields or nonaxisymmetry of the Galactic gravitational field and their importance remains an open question. Measurement errors associated with determining some disk parameters also limit the applicability of multicomponent criterion because the larger the number of constituents the larger errors get accumulated. Future interferometric missions such as $`GAIA`$ and $`SIM`$ will probably solve this problem because of their anticipated accuracy. Nevertheless, even with all the simplifications and the observational uncertainties involved it seems that the Solar neighborhood is stable against purely gravitational axisymmetric perturbations.
## 7 Acknowledgements
It is my great pleasure to thank S.Tremaine for reading the manuscript and making valuable suggestions and B.T.Draine for useful discussions. Author would also like to acknowledge the financial support of this work by the Princeton University Science Fellowship.
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# 1 Introduction
## 1 Introduction
One of the phenomenologically most relevant corners of the string/M-theory moduli space is the Horava-Witten theory . It corresponds to compactifying M-theory on $`R^{10}\times S^1/Z_2`$, which results in an eleven dimensional “brane world”, with two ten dimensional branes sitting at the fixed points of the orbifold, each of them supporting an $`E_8`$ gauge theory. Compactifying six of the dimensions of $`R^{10}`$ on a Calabi-Yau three-fold ($`\mathrm{CY}_3`$), results in a five dimensional brane world with an $`𝒩=2`$ gauged supergravity in the bulk and an $`𝒩=1`$ supersymmetric $`E_8`$ gauge theory on each three-brane at the fixed points. A detailed derivation of the five dimensional low energy effective action for these compactifications was given in .
Such five dimensional brane world compactifications, besides their obvious phenomenological relevance due to the four dimensional $`𝒩=1`$ gauge theories coupled to chiral matter, offer a means for setting the four dimensional cosmological constant (CC) to zero, provided there are adjustable integration constants/parameters. Indeed, in another class of brane world scenarios, namely compactifications of type II string theory on curved spaces (such as $`S^5`$ or deformations of it) with the gauge theory coming from D-branes sitting at orientifold fixed planes, it is indeed possible to argue that one can get flat space after supersymmetry breaking by adjusting integration constants . In this paper, we discuss the same issue in the context of the Horava-Witten theory.
It was shown in the third paper of reference that a consistent compactification of M-theory ala Horava-Witten on a Calabi-Yau three-fold requires a non-trivial non-zero mode flux for the four form field strength $`𝒢`$. This corresponds to wrapped M5-branes around non-trivial $`(1,1)`$ cycles of the Calabi-Yau and one of the consequences of this background flux is a non-zero potential for the volume and shape moduli .
One could also consider the existence, for example of instanton corrections coming from (Euclidean) M5-branes wrapped around the Calabi-Yau manifold. Indeed, this seems to be the case, as discussed in , even though the exact form of this correction to the potential is not known. What we do know, however, is that the volume modulus dependence is the same as for the potential from the compactification, a fact that can be easily justfied once we realize that both contributions can be seen as arising from gauging Abelian isometries of the universal hypermultiplet coset space , . The total potential due to M5-branes can be therefore written as
$$V_{M5}=e^{2\phi }(V_{(1,1)}+V_{(3,3)}),$$
(1)
where $`\phi `$ is the volume modulus, $`V_{(1,1)}`$ and $`V_{(3,3)}`$ correspond to M5-branes wrapped around $`(1,1)`$ and the $`(3,3)`$ cycle respectively.
In section 2, we rederive $`V_{(1,1)}`$ following ignoring for simplicity instanton contributions. In section 3 we discuss the cosmological constant problem from the point of view of the five dimensional brane world after supersymmetry breaking. We will conclude that after supersymmetry breaking there do not exist flat (or nearly flat) brane solutions to these equations. In section 4, we summarize our conclusions.
## 2 Horava-Witten on $`\mathrm{CY}_3`$
The bosonic part of the low energy effective action of $`M`$-theory is the bosonic part of the eleven dimensional supergravity of :
$$S=\frac{1}{2\kappa _{}^{(11)}{}_{}{}^{2}}\{d^{11}x\sqrt{g^{(11)}}^{(11)}\frac{1}{2}(𝒢𝒢+\frac{1}{3}𝒞𝒢𝒢)\},$$
(2)
with $`𝒞`$ the 3-form of 11 dimensional supergravity and $`𝒢=d𝒞`$. This action, when supplemented by appropriate boudary terms, describes the (low energy) strong coupling limit of the heterotic string theory, of Horava and Witten.
We first consider the effective action obtained by reducing (2) on a Calabi-Yau 3-fold . The appropriate metric ansatz is
$$ds_{}^{(11)}{}_{}{}^{2}=ds_{str}^{(5)}{}_{}{}^{2}+ds_{CY}^{(6)}{}_{}{}^{2}.$$
(3)
The subscript $`{}_{}{}^{\prime \prime }str^{\prime \prime }`$ indicates string frame and superscripts indicate the dimensionality. We assume, following , that the “standard embedding” requires in the Horava-Witten picture a non-trivial $`𝒢`$-flux of the form
$$𝒢=\frac{i}{2𝒱}n_iG^{ij}(\omega _j),$$
(4)
where the $`n_i`$ are integers. The form of the $`𝒢`$-flux for non-standard embeddings and additional five branes in the bulk is similar; the only modification is that the $`n_i`$ in the above is replaced by a sum over fluxes . The low energy spectrum includes the five dimensional graviton multiplet, $`h^{1,1}1`$ vector multiplets containing one shape modulus scalar ($`b^i`$) each ($`i=1\mathrm{}h^{1,1}`$ but there is one constraint, see below), and the universal hypermultiplet that includes the volume modulus ($`\phi `$), the scalar dual to the five dimensional three-form ($`\sigma `$) and a pair of complex scalars ($`\xi `$ and $`\overline{\xi }`$). The metric on the vector moduli space is
$$G_{ij}(a)=\frac{i}{2𝒱}_{CY}\omega _i\omega _j=\frac{1}{2}_i_j\mathrm{ln}𝒱(a),$$
(5)
The K$`\ddot{\mathrm{a}}`$hler form is $`J=a^i\omega _i`$, with the $`\omega _i`$ a basis of $`h^{1,1}`$ 2-forms and $`a^i`$ are the $`h^{1,1}`$ K$`\ddot{\mathrm{a}}`$hler moduli of the Calabi-Yau. The Calabi-Yau volume $`𝒱`$ is
$$𝒱(a)=\frac{1}{3!}_{CY}JJJ=\frac{1}{6}c_{ijk}a^ia^ja^k,$$
(6)
with $`c_{ijk}`$ the intersection numbers of the Calabi-Yau.
For our purposes, it is sufficient to consider the part of the effective action which is 5D gravity coupled to the scalars of the vector multiplets and the volume modulus (breathing mode) of the universal hypermultiplet. The corresponding five dimensional string frame bulk effective action resulting from the first two terms of (2), is
$$S=\frac{1}{2\kappa _{}^{(5)}{}_{}{}^{2}}d^5x\sqrt{g_{str}^{(5)}}\left[𝒱^{(5)}+\left(𝒱G_{ij}(a)+_i_j𝒱\right)_Ia^i^Ia^j\frac{1}{4𝒱}G^{ij}(a)n_in_j\right].$$
(7)
The five dimensional index is $`I=\{\mu ,r\}`$. We have neglected the terms coming from the Chern-Simons term in (2) since they are irrelevant to our dicussion. After the Weyl rescaling $`ds_{E}^{(5)}{}_{}{}^{2}=𝒱^{\frac{2}{3}}ds_{str}^{(5)}{}_{}{}^{2}`$ and separation of the volume modulus from the shape moduli which can be done by defining $`b^i=a^i𝒱^{1/3}`$, we arrive at a bulk action of the form
$$S_{bulk}=\frac{1}{2\kappa _{}^{(5)}{}_{}{}^{2}}d^5x\sqrt{g_E^{(5)}}\left[^{(5)}G_{ij}(b)_Ib^i^Ib^j\frac{1}{2}_I\phi ^I\phi \frac{1}{4}e^{2\phi }G^{ij}(b)n_in_j+\lambda (c_{ijk}b^ib^jb^k6)\right],$$
(8)
where we have defined $`𝒱=e^\phi `$ and $`\lambda `$ is a Lagrange multiplier. This is not the whole relevant Horava-Witten 5D action, because we have not taken into account yet the Horava-Witten Wall/M5-brane action sitting at the fixed points of the orbifold. The additional brane terms are
$$S_{branes}=\frac{1}{2\kappa _{}^{(5)}{}_{}{}^{2}}d^4x\sqrt{g_E^{(4)}}T_1(\phi )\delta (r)+\frac{1}{2\kappa _{}^{(5)}{}_{}{}^{2}}d^4x\sqrt{g_E^{(4)}}T_2(\phi )\delta (r\pi R),$$
(9)
with $`g_E^{(4)}`$ the induced metric on the brane (we will assume static gauge and ignore fluctuations) and $`T_{1,2}`$ are the tensions of the branes. As we mentioned, these 3-branes arise from the M5-branes of the original 11 dimensional theory that are wrapped around non-trivial 2-cycles of the Calabi-Yau and the brane tensions at the string scale, i.e. before supersymmetry breaking, are simply proportional to the volume of the 2-cycle , :
$$T_1=T_2(\phi )=\frac{2}{1!}J=2e^\phi \zeta .$$
(10)
Defining $`n^2G^{ij}n_in_j=n_in^i`$ and eliminating the Lagrange multiplier, we can write the equations of motion for the $`b^i`$ as
$`{\displaystyle \frac{1}{4}}e^{2\phi }(_kn^2)+{\displaystyle \frac{1}{6}}e^{2\phi }b_kb^l(_ln^2)`$
$`{\displaystyle \frac{3}{2}}c_{ijk}(_\alpha b^i)(^\alpha b^j)+{\displaystyle \frac{2}{3}}c_{ijl}b^lb_k(_\alpha b^i)(^\alpha b^j)+{\displaystyle \frac{4}{3}}b_kb_i(_\alpha ^\alpha b^i)`$
$`=e^\phi (\delta (r)\delta (r\pi R))\left({\displaystyle \frac{4}{3}}n_lb^lb_k2n_k\right).`$ (11)
To simplify these equations, we will assume that the $`b^i`$ take constant values in the vacuum. All the terms with the covariant derivatives drop out, the equations of motion for $`b^i`$ decouple from the equation of motion for $`\phi `$ and then the unique ansatz that solves the remaining of (11) is ,
$$\zeta b_k=\frac{3}{2}n_k,$$
(12)
where $`\zeta n_lb^l`$. By solving the above system of equations, we obtain the values that the squashing modes take in the vacuum. Notice that for the ansatz (12), the term multiplying the $`\delta `$-functions in (11) vanishes, which is a necessary condition for a consistent solution with constant $`b^i`$. Now we turn to the breathing mode $`\phi `$.
The bulk action, in the Einstein frame, with all the moduli besides the breathing mode $`\phi `$ stabilized, is
$$S_{bulk}=\frac{1}{2\kappa _{}^{(5)}{}_{}{}^{2}}d^5x\sqrt{g_E^{(5)}}\left[^{(5)}\frac{1}{2}(\phi )^2\frac{1}{6}\zeta ^2e^{2\phi }\right].$$
(13)
In the above, the constant $`\zeta ^2`$ in front of the potential is determined in terms of the integer $`n_i`$ (therefore it is not a continuous quantity) and it corresponds to contributions from fluxes associated with wrapping M5 branes around the non trivial $`(1,1)`$ cycles of the Calabi Yau. ¿From the 5D gauged supergravity point of view, the potential comes from gauging the $`U(1)`$ isometry of the universal hypermultiplet moduli space $`SU(2,1)/U(2)`$ corresponding to a shift symmetry of the scalar $`\sigma `$. It is interesting to note here that the universal hypermultiplet moduli space has more $`U(1)`$ isometries that could be gauged , associated with shift symmetries of $`\xi `$ and $`\overline{\xi }`$. Gauging all of the $`U(1)`$ isometries, in principle, can produce additional terms in the potential. ¿From the M-theory point of view, the additional pieces correspond to M5-branes wrapped around the $`(3,3)`$ cycle (i.e. the whole Calabi-Yau) and/or to M2-branes wrapped around $`(3,0)`$ and $`(0,3)`$ cycles of the Calabi-Yau . The potential in (13) is therefore the potential for $`\mathrm{CY}_3`$ compactifications of the Horava-Witten theory for constant shape moduli and without an M5/M2 instanton gas. In the following, we will neglect instanton contributions, since our subsequent arguments about the cosmological constant are not affected by their presence.
## 3 Supersymmetry Breaking and the Cosmological Constant
The equations of motion from action (13), for $`\phi =\phi (r)`$ and for
$$ds_{E}^{(5)}{}_{}{}^{2}=e^{2A(r)}\eta _{\mu \nu }dx^\mu dx^\nu +(dr)^2,$$
(14)
are
$$\phi ^{\prime \prime }+4A^{}\phi ^{}=\frac{V(\phi )}{\phi }$$
(15)
$$A^{\prime \prime }=\frac{1}{6}\phi ^2$$
(16)
$$A^2=\frac{1}{12}V(\phi )+\frac{1}{24}\phi ^2,$$
(17)
where $`V(\phi )e^{2\phi }V_{(1,1)}=e^{2\phi }\frac{1}{6}\zeta ^2`$. The prime denotes differentiation with respect to $`r`$. The first order set of equations corresponding to the above second order set, in terms of the superpotential $`W`$, is
$$\phi ^{}=\frac{W}{\phi }$$
(18)
$$A^{}=\frac{1}{6}W$$
(19)
$$V(\phi )=\frac{1}{2}\left(\frac{W}{\phi }\right)^2\frac{1}{3}W^2.$$
(20)
We note here that the proper ansatz for the five dimensional metric would be to replace the flat Minkowski metric $`\eta _{\mu \nu }`$ by $`g_{\mu \nu }^{(A)dS}`$ (i.e. a metric for deSitter or Anti-deSitter four-space) but since the (possibly) measured value of the CC is many tens of orders of magnitude smaller than all relevant scales in this analysis, this is not a meaningful difference. We can rewrite (20) as
$$\left(\frac{\dot{W}}{\sqrt{2V}}\right)^2\left(\frac{W}{\sqrt{3V}}\right)^2=1,$$
(21)
where the dot stands for differentiation with respect to $`\phi `$. The ansatz
$$\dot{W}=\sqrt{2V}\mathrm{cosh}f(\phi )\mathrm{and}W=\sqrt{3V}\mathrm{sinh}f(\phi )$$
(22)
allows us to separate variables, so that combining the above, we can integrate over $`\phi `$ and $`f`$:
$$𝑑\phi =\frac{df}{\sqrt{\frac{2}{3}}\mathrm{tanh}f}.$$
(23)
The integral yields
$$c_1e^\phi =(\mathrm{tanh}f\sqrt{\frac{2}{3}})^3\frac{(\mathrm{tanh}f+1)^{\sqrt{\frac{3}{2}}(1\sqrt{\frac{2}{3}})}}{(\mathrm{tanh}f1)^{\sqrt{\frac{3}{2}}(1+\sqrt{\frac{2}{3}})}},$$
(24)
with $`c_1`$ an integration constant. Now equations (18) and (19) can be integrated as follows:
$$r+c_2=\frac{\sqrt{3}}{\zeta }𝑑\phi e^\phi \sqrt{1\mathrm{tanh}^2f(\phi )}$$
(25)
and
$$A=\frac{1}{6}\sqrt{\frac{3}{2}}𝑑\phi \mathrm{tanh}f(\phi ).$$
(26)
The method therefore to find a solution is to invert (24) for $`\mathrm{tanh}f`$ in terms of $`\phi `$, substitute into (25) and (26) and integrate. Then, to find the expression for $`\phi (r)`$, invert (25) and finally use this to obtain the expression for $`A(r)`$.
The construction of the solution suited for an orbifold is described in detail for example in . Consistency of our flat domain wall ansatz with the vanishing of the total 4 dimensional CC, amounts to satisfying the following jump conditions:
$$2\phi ^{}(r)=+\frac{T_1}{\phi }(\phi (r))_{r=0}$$
(27)
$$2A^{}(r)=\frac{1}{6}T_1(\phi (r))_{r=0}$$
(28)
$$2\phi ^{}(r)=\frac{T_2}{\phi }(\phi (r))_{r=\pi R}$$
(29)
$$2A^{}(r)=+\frac{1}{6}T_2(\phi (r))_{r=\pi R}.$$
(30)
Equations (27), (29), (28) and (30) is the system of equations that has to be satisfied in a model with vanishing CC. To satisfy these, we have the two integration constants $`c_1`$ and $`c_2`$, the size of the orbifold $`R`$ and the discrete values of the $`n_i`$ that determine $`\zeta ^2`$. At the supersymmetric point, we can rewrite (27)-(30) in a more convenient form using (22):
$$\mathrm{cosh}f\left(\phi (0)\right)=\sqrt{3}$$
(31)
$$\mathrm{sinh}f\left(\phi (0)\right)=\sqrt{2}$$
(32)
$$\mathrm{cosh}f\left(\phi (\pi R)\right)=\sqrt{3}$$
(33)
$$\mathrm{sinh}f\left(\phi (\pi R)\right)=\sqrt{2}.$$
(34)
A simple solution to the above can be found if we take $`W=\zeta e^\phi `$, which is a solution to (20). Then,
$$\mathrm{cosh}f\left(\phi (r)\right)=\pm \sqrt{3}\mathrm{and}\mathrm{sinh}f\left(\phi (r)\right)=\pm \sqrt{2},$$
(35)
and the equations of motion (25) and (26) can be solved easily, yielding
$$\phi (r)=\mathrm{ln}(\zeta |r|+c)\mathrm{and}A(r)=\frac{1}{6}\mathrm{ln}(\zeta |r|+c),$$
(36)
with $`c\zeta c_2`$ the (only) integration constant. This is a trivial solution, trivial in the sense that the jump conditions are satisfied identically, without restriction on $`\zeta `$ and the integration constant $`c`$. In fact, since $`\mathrm{tanh}f\left(\phi (r=0,\pi R)\right)=\sqrt{\frac{2}{3}}`$ at the supersymmetric point, from (23) we see that at the supersymmetric point, this is actually the only possible solution. <sup>1</sup><sup>1</sup>1In general, even for compactifications on curved spaces such as type IIB on (squashed) $`S^5`$ for example, such a trivial solution is always possible if the brane tension is taken to be $`W`$ . Turning this around, we conclude that after supersymmetry breaking, the simple ansatz $`W=\zeta e^\phi `$ can not be used anymore to satisfy (27)-(30). But this is expected, since in general it is not possible to satisfy the four jump equations with the three parameters $`\zeta `$, $`c`$ and $`R`$ (even if $`\zeta `$ were continuous). Thus, for a consistent model with zero cosmological constant away from the supersymmetric point, we have to look for more general solutions to (20).
After supersymmetry breaking (on the brane), the brane tensions get renormalized as
$$T_1(\phi )2\zeta e^\phi (1+ϵ\psi _1(\zeta ϵ^\phi )),T_2(\phi )2\zeta e^\phi (1+ϵ\psi _2(\zeta ϵ^\phi ))$$
(37)
with $`ϵ`$ being a small parameter characterizing the size of supersymmetry breaking. The scalar potential on the other hand, at least at the 5D level, remains the same .
The jump conditions therefore, after supersymmetry breaking become
$$\mathrm{cosh}f\left(\phi (0)\right)=\sqrt{3}\left(1+ϵ\frac{d}{dx}(x\psi _1(x))|_{x=x(0)}\right)$$
(38)
$$\mathrm{sinh}f\left(\phi (0)\right)=\sqrt{2}\left(1+ϵ\psi _1(x)|_{x=x(0)}\right)$$
(39)
$$\mathrm{cosh}f\left(\phi (\pi R)\right)=\sqrt{3}\left(1+ϵ\frac{d}{dx}(x\psi _2(x))|_{x=x(\pi R)}\right)$$
(40)
$$\mathrm{sinh}f\left(\phi (\pi R)\right)=\sqrt{2}\left(1+ϵ\psi _2(x)|_{x=x(\pi R)}\right),$$
(41)
where $`x=\zeta e^\phi `$. Now let us recall that the difference of $`\mathrm{cosh}f(\mathrm{sinh}f)`$ from $`\sqrt{3}(\sqrt{2})`$ vanishes like $`c_1`$ which is of $`O(ϵ)`$. Thus we may write the above equations as relations between $`O(1)`$ functions. Now in accordance with the general argument of , the integration constants $`c_1,c_2`$ may be adjusted to satisfy say the first two equations above. In addition the distance $`\pi R`$ may be freely adjusted to satisfy one more matching condition leaving us with one more condition to satisfy. To satisfy the fourth condition then requires a fine tuning of the parameters in the bulk potential. In our case this means a fine tuning of the quantity $`\zeta `$. However, as noted earlier, the latter is completely determined in terms of integers $`n_i`$ governing the fluxes and the integers $`d_{ijk}`$ and $`h_{11}`$ of the Calabi-Yau manifold. There is no way these can be chosen to satisfy this last equation since any change of these integers results in a O(1) change of the functions $`\psi _{1,2}`$ and it would be a miracle if for any choice of the integers the condition could be satisfied <sup>2</sup><sup>2</sup>2Of course one does not require exact satisfaction of the matching conditions. Strictly speaking all we need is that the cosmological constant on the brane be of $`O(10^{120}M_P^4)`$. This would still imply an adjustment of parameters to this accuracy as explained in and is clearly ruled out in the present situation.. In other words there is no adjustable continuous parameter that could be tuned to satisfy the equation.
Let us now discuss possible generalizations. One possible non-vanishing supersymmetric correction to the scalar potential comes from the instanton gas obtained from wrapping M5-branes around the Calabi-Yau manifold or M2-branes wrapping around 3-cycles<sup>3</sup><sup>3</sup>3It is not clear whether there are actually non-zero contributions due to transverse M2 branes parallel to the walls since there is no S-dual analog in the weak coupling i.e. the heterotic string limit.. The potential would then be as in (1) plus an analogous contribution coming from M2-branes. Clearly, the original problem associated with the discreteness of the fourth adjustable parameter still remains since we are still talking about contributions that are completely determined by a set of integers.
Another possible generalization would be to look for solutions to the 5D equations of motion (11) with $`r`$-dependent $`b^i`$. Such a solution, for the supersymmetric case, was found in . After supersymmetry breaking, we would have to satisfy two additional jump conditions for each $`b^i`$, which is possible, since we get two additional integration constants for each $`b^i`$ from their equations of motion. However, as before, we still have one fine tuning to do and the previous argument, that since there is no continuous parameter there is no possibility of doing this, still applies.
## 4 Conclusions
We have discussed in this paper the two brane scenario coming from the Horava-Witten theory compactified on a Calabi-Yau manifold obtained by Lukas et al . Those authors discussed the supersymmetry preserving case when there is sufficient degeneracy in the system of equations to satisfy the matching conditions without any adjustment of integration constants. When supersymmetry is broken however, as has been discussed in earlier work, we need four adjustable parameters (integration constants). We have shown explicitly that how these four parameters arise in the case at hand but that since one of the parameters necessarily takes discrete values, it does not seem possible to tune the cosmological constant to zero (or a small value).
The situation discussed in this paper is to be contrasted with the one obtained for compactifications on Ricci non-flat manifolds, as for example in IIB string theory on a (squashed) $`S^5`$, where one has two exponentials in the potential for the breathing mode. In these compactifications, besides a discrete parameter like $`\zeta `$ which is also present, (coming from turning on five-form flux) there is an additional continuous parameter not present in Calabi-Yau compactifications. It is associated with the Ricci curvature of the compact space. Thus, in such cases, it is in principle always possible to satisfy the jump conditions and readjust integration constants by an arbitrarily small quantity after supersymmetry breaking to get flat brane solutions i.e. a cosmological constant that is zero<sup>4</sup><sup>4</sup>4Or indeed of order $`O((10^3eV)^4`$ as seems to be fashionable these days! after supersymmetry breaking. Of course in the absence of a theory of integration constants one cannot claim to have solved the cosmological constant problem even in the case of compactifications on Ricci-non-flat manifolds. All one could claim there is that one has sufficient freedom to get a flat space solution. In the Ricci flat case what we have argued is that this freedom does not seem to exist.
Acknowledgements
This work is partially supported by the Department of Energy contract No. DE-FG02-91-ER-40672.
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# The Scalar-flat Kähler Metric and Painlevé III
## 1 introduction
The aim of this paper is to analyse scalar-flat Kähler metrics $`g`$ in real dimmension four admitting an isometric action of $`SU(2)`$ with generically three-dimensional orbits. A scalar-flat Kähler metric is a metric with zero scalar curvature which is Kählerian with respect to a complex structure on $`M`$. It is automatically anti-self-dual with respect to the canonical orientation.
Hitchin shows that the $`SU(2)`$-invariant anti-self-dual metric is generically specified by a solution of Painlevé VI type equation, and if the metric is scalar-flat Kähler it is specified by a solution of Painlevé III type equation. Hitchin used the twistor correspondence to associate the anti-self-dual equation and Painlevé equation. The lifted action of $`SU(2)`$ determines a pre-homogenious action of $`SU(2)`$ on the twistor space $`Z`$, and it determines a isomonodromic family of connections on $`^1`$, and then we have Painlevé equations. In this way, Dancer analyse the scalar-flat Kähler metric with $`SU(2)`$-symmetry. In , Hichin obtains a complete classification of anti-self-dual Einstein metrics admitting an isometric action of $`SU(2)`$ with three-dimensional orbits. For the completeness analysis, it is impotant to have the explicit form of anti-self-dual equation.
If the metric is diagonal, the explicit form of anti-self-dual equation is known, but if the metric is non-diagonal, it is known very little. For scalar-flat Kähler metric, complex structure is not known for non-diagonal metric.
In Section 2 we show how non-diagonal metric is represented, and by use of the block form of curvature tensor given by Besse , we have the ninth-order system equivalent to the anti-self-dual equation.
In Section 3 we establishes the relationship between $`SU(2)`$-invariant anti-self-dual manifold and the isomonodromic deformation. It is essentially equivalent to Hitchin’s ansatz. Still in our way, we have the explicit form of the isomonodromic deformations, and we have the condition that the corresponding Painlevé equation is type III.
In Section 4, we show that the anti-self-dual equation reduce to Painlevé III if and only if the metric admits an Hermitian structure. In this case, the anti-self-dual equation is equivalent to a seventh-order system, and it also admits Kähler structure, the seventh-order system reduce to a sixth-order system.
Acknowledgements
The author would express his sincere gratitude to Professor Yousuke Ohyama, who intoroduced him to the subject, for enlightening discussions.
## 2 The non-diagonal anti-self-dual equations
We can write the $`SU(2)`$-invariant metric in the form
$`g=f(\tau )d\tau ^2+{\displaystyle \underset{l,m=1}{\overset{3}{}}}h_{lm}(\tau )\sigma _l\sigma _m,`$
where $`\{\sigma _1,\sigma _2,\sigma _3\}`$ is a basis of left invariant one-forms on each $`SU(2)`$-orbit satisfying
$`d\sigma _1`$ $`=\sigma _2\sigma _3,`$ $`d\sigma _2`$ $`=\sigma _3\sigma _1,`$ $`d\sigma _3=\sigma _1\sigma _2.`$
Using the Killing form, we can diagonalize the metric $`g`$ on each $`SU(2)`$-orbit. Then we can express the metric as follows:
$`g=(abc)^2dt^2+a^2d\stackrel{~}{\sigma }_1^2+b^2\stackrel{~}{\sigma }_2^2+c^2\stackrel{~}{\sigma }_3^2,`$
where $`t=t(\tau ),a=a(t),b=b(t),c=c(t)`$ and
$`\left(\begin{array}{c}\stackrel{~}{\sigma }_1\\ \stackrel{~}{\sigma }_2\\ \stackrel{~}{\sigma }_3\end{array}\right)=R(t)\left(\begin{array}{c}\sigma _1\\ \sigma _2\\ \sigma _3\end{array}\right),`$
$`R(t)`$ is $`SO(3)`$-valued function.
Since $`\dot{R}R^1`$ (where $`\dot{}=\frac{d}{dt}`$) is $`𝔰𝔬(3)`$-valued, we can write
$`d\left(\begin{array}{c}\stackrel{~}{\sigma }_1\\ \stackrel{~}{\sigma }_2\\ \stackrel{~}{\sigma }_3\end{array}\right)`$ $`=R(t)\left(\begin{array}{c}\sigma _2\sigma _3\\ \sigma _3\sigma _1\\ \sigma _2\sigma _2\end{array}\right)+\dot{R}dt\left(\begin{array}{c}\sigma _1\\ \sigma _2\\ \sigma _3\end{array}\right)`$
$`=\left(\begin{array}{c}\stackrel{~}{\sigma }_2\stackrel{~}{\sigma }_3\\ \stackrel{~}{\sigma }_3\stackrel{~}{\sigma }_1\\ \stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2\end{array}\right)+\left(\begin{array}{ccc}\hfill 0& \hfill \xi _3& \hfill \xi _2\\ \hfill \xi _3& \hfill 0& \hfill \xi _1\\ \hfill \xi _2& \hfill \xi _1& \hfill 0\end{array}\right)dt\left(\begin{array}{c}\stackrel{~}{\sigma }_1\\ \stackrel{~}{\sigma }_2\\ \stackrel{~}{\sigma }_3\end{array}\right),`$
for some $`\xi _1=\xi _1(t)`$, $`\xi _2=\xi _2(t)`$, $`\xi _3=\xi _3(t)`$.
If $`\xi _1=0,\xi _2=0,\xi _3=0`$, then the matrix $`(h_{lm})`$ can be chosen to be diagonal for all $`\tau `$, and then we say that $`g`$ has diagonal form.
In this paper we mainly study the non-diagonal case.
To compute the curvature tensor we choose a basis for $`^2`$
$`\{\mathrm{\Omega }_1^+,\mathrm{\Omega }_2^+,\mathrm{\Omega }_3^+,\mathrm{\Omega }_1^{}\mathrm{\Omega }_2^{},\mathrm{\Omega }_3^{}\},`$
where
$`\mathrm{\Omega }_1^+`$ $`=a^2bcdt\stackrel{~}{\sigma }_1+bc\stackrel{~}{\sigma }_2\stackrel{~}{\sigma }_3,`$
$`\mathrm{\Omega }_2^+`$ $`=ab^2cdt\stackrel{~}{\sigma }_2+ca\stackrel{~}{\sigma }_3\stackrel{~}{\sigma }_1,`$
$`\mathrm{\Omega }_3^+`$ $`=abc^2dt\stackrel{~}{\sigma }_3+ab\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2,`$
$`\mathrm{\Omega }_1^{}`$ $`=a^2bcdt\stackrel{~}{\sigma }_1bc\stackrel{~}{\sigma }_2\stackrel{~}{\sigma }_3,`$
$`\mathrm{\Omega }_2^{}`$ $`=ab^2cdt\stackrel{~}{\sigma }_2ca\stackrel{~}{\sigma }_3\stackrel{~}{\sigma }_1,`$
$`\mathrm{\Omega }_3^{}`$ $`=abc^2dt\stackrel{~}{\sigma }_3ab\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2.`$
With respect to this frame, the curvature tensor has the following block form
$`\left(\begin{array}{cc}\hfill A& \hfill B\\ \hfill {}_{}{}^{t}B& \hfill D\end{array}\right),`$
where $`s=4\text{trace}D`$ is the scalar curvature, $`W^+=A\frac{1}{12}s`$ and $`W^{}=D\frac{1}{12}s`$ are the self-dual and anti-self-dual parts of the Weyl tensor and $`B`$ is the trace free parts of Ricci tensor.
We set $`w_1=bc,w_2=ca,w_3=ab`$ and determine $`\alpha _1,\alpha _2,\alpha _3`$ by
$`\begin{array}{cc}\hfill \dot{w}_1& =w_2w_3+w_1(\alpha _2+\alpha _3),\hfill \\ \hfill \dot{w}_2& =w_3w_1+w_2(\alpha _3+\alpha _1),\hfill \\ \hfill \dot{w}_3& =w_1w_2+w_3(\alpha _1+\alpha _2).\hfill \end{array}`$ (1)
Calculating the condition $`A=0`$, we have the following theorem.
###### Theorem 2.1
The metric is anti-self-dual with vanishing scalar curvature if and only if $`\alpha _1,\alpha _2,\alpha _3`$ and $`\xi _1,\xi _2,\xi _3`$ satisfies the following equations:
$`\begin{array}{cc}\hfill \dot{\alpha }_1=& \alpha _2\alpha _3+\alpha _1(\alpha _2+\alpha _3)+{\displaystyle \frac{1}{4}}(w_2^2w_3^2)^2\left({\displaystyle \frac{\xi _1}{w_2w_3}}\right)^2\hfill \\ & +{\displaystyle \frac{1}{4}}(w_3^2w_1^2)(3w_1^2+w_3^2)\left({\displaystyle \frac{\xi _2}{w_3w_1}}\right)^2\hfill \\ & +{\displaystyle \frac{1}{4}}(w_2^2w_1^2)(3w_1^2+w_2^2)\left({\displaystyle \frac{\xi _3}{w_1w_2}}\right)^2,\hfill \\ \hfill \dot{\alpha }_2=& \alpha _3\alpha _1+\alpha _2(\alpha _3+\alpha _1)+{\displaystyle \frac{1}{4}}(w_3^2w_1^2)^2\left({\displaystyle \frac{\xi _2}{w_3w_1}}\right)^2\hfill \\ & +{\displaystyle \frac{1}{4}}(w_1^2w_2^2)(3w_2^2+w_1^2)\left({\displaystyle \frac{\xi _3}{w_1w_2}}\right)^2\hfill \\ & +{\displaystyle \frac{1}{4}}(w_3^2w_2^2)(3w_2^2+w_3^2)\left({\displaystyle \frac{\xi _1}{w_2w_3}}\right)^2,\hfill \\ \hfill \dot{\alpha }_3=& \alpha _1\alpha _2+\alpha _3(\alpha _1+\alpha _2)+{\displaystyle \frac{1}{4}}(w_1^2w_2^2)^2\left({\displaystyle \frac{\xi _3}{w_1w_2}}\right)^2\hfill \\ & +{\displaystyle \frac{1}{4}}(w_2^2w_3^2)(3w_3^2+w_2^2)\left({\displaystyle \frac{\xi _1}{w_2w_3}}\right)^2\hfill \\ & +{\displaystyle \frac{1}{4}}(w_1^2w_3^2)(3w_3^2+w_1^2)\left({\displaystyle \frac{\xi _2}{w_3w_1}}\right)^2,\hfill \end{array}`$ (2)
and
$`\begin{array}{cc}\hfill (w_2^2w_3^2){\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{\xi _1}{w_2w_3}}\right)=& {\displaystyle \frac{\xi _2}{w_3w_1}}{\displaystyle \frac{\xi _3}{w_1w_2}}(2w_2^2w_3^2+w_3^2w_1^2+w_1^2w_2^2)\hfill \\ & +{\displaystyle \frac{\xi _1}{w_2w_3}}(\alpha _2w_2^2\alpha _3w_3^2+3\alpha _2w_3^2+3\alpha _3w_2^2),\hfill \\ \hfill (w_3^2w_1^2){\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{\xi _2}{w_3w_1}}\right)=& {\displaystyle \frac{\xi _3}{w_1w_2}}{\displaystyle \frac{\xi _1}{w_2w_3}}(2w_3^2w_1^2+w_1^2w_2^2+w_2^2w_3^2)\hfill \\ & +{\displaystyle \frac{\xi _2}{w_3w_1}}(\alpha _3w_3^2\alpha _1w_1^2+3\alpha _3w_1^2+3\alpha _1w_2^2),\hfill \\ \hfill (w_1^2w_2^2){\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{\xi _3}{w_1w_2}}\right)=& {\displaystyle \frac{\xi _1}{w_1w_3}}{\displaystyle \frac{\xi _2}{w_3w_1}}(2w_1^2w_2^2+w_2^2w_3^2+w_3^2w_1^2)\hfill \\ & +{\displaystyle \frac{\xi _3}{w_1w_2}}(\alpha _1w_1^2\alpha _2w_2^2+3\alpha _1w_2^2+3\alpha _2w_1^2).\hfill \end{array}`$ (3)
###### Remark 2.2
If $`\xi _1=0`$, $`\xi _2=0`$ and $`\xi _3=0`$ then the system (1), (2), (3) reduce to a sixth-order system given by Tod. Furthermore, if $`\alpha _1=w_1,\alpha _2=w_2,\alpha _3=w_3`$ then (1),(2),(3) reduce to a third-order system which determines Atiyha-Hitchin family , and if $`\alpha _1=0,\alpha _2=0,\alpha _3=0`$ then the system reduce to a third-order system which determines BGPP family .
###### Remark 2.3
If $`w_2=w_3`$, then we can set $`\xi _1=0`$, $`\xi _2=0`$ and $`\xi _3=0`$ by taking another flame. This is also a diagonal case. Therefore we assume $`(w_2w_3)(w_3w_1)(w_1w_2)0`$.
## 3 The Isomonodromic Deformations and Painlevé equation
Let $`(M,g)`$ be an oriented Riemannian four manifold. We define $`Z`$ to be the unit sphere bundle in the bundle of self-dual two-forms, and let $`\pi :ZM`$ denote the projection. Each point $`z`$ in the fiber over $`\pi (z)`$ defines a complex structure on the tangent space $`T_{\pi (z)}M`$, compatible with the metric and its orientation.
Using the Levi-Civita connection, we can split the tangent space $`T_zZ`$ into horizontal and vertical spaces, and the projection $`\pi `$ identifies the horizontal space with $`T_{\pi (z)}M`$. This space has a complex structure defined by $`z`$ and the vertical space is the tangent space of the fiber $`S^3^1`$ which has its natural complex structure. The almost complex structure on $`Z`$ is integrable if and only if the metric is anti-self-dual . In this situation $`Z`$ is called the twistor space of $`(M,g)`$ and The fibers are called the real twistor lines.
The almost complex structure on $`Z`$ can be determined by the following $`(1,0)`$-forms:
$`\begin{array}{cc}\hfill \mathrm{\Theta }_1=& z(e^2+\sqrt{1}e^3)(e^0+\sqrt{1}e^1),\hfill \\ \hfill \mathrm{\Theta }_2=& z(e^0\sqrt{1}e^1)+(e^2\sqrt{1}e^3),\hfill \\ \hfill \mathrm{\Theta }_3=& dz+{\displaystyle \frac{1}{2}}z^2(\omega _2^0\omega _1^3+\sqrt{1}(\omega _3^0\omega _2^1))\hfill \\ & \sqrt{1}z(\omega _1^0\omega _3^2)+{\displaystyle \frac{1}{2}}(\omega _2^0\omega _1^3\sqrt{1}(\omega _3^0\omega _2^1)),\hfill \end{array}`$ (4)
where $`\{e^0,e^1,e^2,e^3\}`$ is an orthonormal flame, and $`\omega _j^i`$ are the connection forms determined by $`de^i+\omega _j^ie^j=0`$ and $`\omega _j^i+\omega _i^j=0`$. Then the anti-self-dual condition is
$`d\mathrm{\Theta }_1`$ $`0,`$ $`d\mathrm{\Theta }_2`$ $`0,`$ $`d\mathrm{\Theta }_3`$ $`0`$ $`(\text{mod}\mathrm{\Theta }_1,\mathrm{\Theta }_2,\mathrm{\Theta }_3).`$ (5)
If the metric is $`SU(2)`$ invariant, we can write
$`\left(\begin{array}{c}\mathrm{\Theta }_1\\ \mathrm{\Theta }_2\\ \mathrm{\Theta }_3\end{array}\right)=\left(\begin{array}{c}0\\ 0\\ 1\end{array}\right)dz+\left(\begin{array}{c}v_1\\ v_2\\ v_3\end{array}\right)dt+A\left(\begin{array}{c}\sigma _1\\ \sigma _2\\ \sigma _3\end{array}\right),`$ (18)
where $`v_1=v_1(z,t)`$, $`v_2=v_2(z,t)`$, $`v_3=v_3(z,t)`$; $`A=\left(a_{ij}(z,t)\right)_{i,j=1,2,3}`$.
If $`\text{det}A0`$, then metric is in the BGPP family .
If $`\text{det}A0`$, then we can write
$`\left(\begin{array}{c}\sigma _1\\ \sigma _2\\ \sigma _3\end{array}\right)A^1(\left(\begin{array}{c}0\\ 0\\ 1\end{array}\right)dz+\left(\begin{array}{c}v_1\\ v_2\\ v_3\end{array}\right)dt)=:\left(\begin{array}{c}\varsigma _1\\ \varsigma _2\\ \varsigma _3\end{array}\right),`$ (31)
and then
$`d\left(\begin{array}{c}\varsigma _1\\ \varsigma _2\\ \varsigma _3\end{array}\right)\left(\begin{array}{c}\varsigma _2\varsigma _3\\ \varsigma _3\varsigma _1\\ \varsigma _1\varsigma _2\end{array}\right).`$ (38)
Since $`\varsigma _1,\varsigma _2,\varsigma _3`$ are one-forms on $`(z,t)`$plane,
$`d\left(\begin{array}{c}\varsigma _1\\ \varsigma _2\\ \varsigma _3\end{array}\right)=\left(\begin{array}{c}\varsigma _2\varsigma _3\\ \varsigma _3\varsigma _1\\ \varsigma _1\varsigma _2\end{array}\right).`$ (45)
If we set
$`\mathrm{\Sigma }`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}\sqrt{1}\varsigma _1& \varsigma _3+\sqrt{1}\varsigma _2\\ \varsigma _3+\sqrt{1}\varsigma _2& \sqrt{1}\varsigma _1\end{array}\right)`$ (48)
$`=:B_1dzB_2dt,`$ (49)
then
$`d\mathrm{\Sigma }+\mathrm{\Sigma }\mathrm{\Sigma }=0.`$ (50)
This is the isomonodromic condition of the equation
$`\left({\displaystyle \frac{d}{dz}}B_1\right)\left(\begin{array}{c}y_1\\ y_2\end{array}\right)=0.`$ (53)
$`B_1`$ has poles on $`\{z|\text{det}A=0\}`$.
###### Lemma 3.1
$`\text{det}A=0`$ is equivalent to the following equation
$$\begin{array}{c}z^4\left(\left(\alpha _2+\alpha _3\right)\sqrt{1}X_1\right)2z^3\left(X_2\sqrt{1}X_3\right)+2z^2\left(2\alpha _1+\alpha _2+\alpha _3\right)\hfill \\ \hfill +2z\left(X_2+\sqrt{1}X_3\right)+\left(\left(\alpha _2+\alpha _3\right)+\sqrt{1}X_1\right)=0,\end{array}$$
(54)
where
$`X_1`$ $`={\displaystyle \frac{w_2^2w_3^2}{w_2w_3}}\xi _1,`$ $`X_2`$ $`={\displaystyle \frac{w_3^2w_1^2}{w_3w_1}}\xi _2,`$ $`X_3`$ $`={\displaystyle \frac{w_1^2w_2^2}{w_1w_2}}\xi _3.`$
For this lemma, generically $`B_1`$ has four simple poles.
###### Theorem 3.2
The anti-self-dual equation on $`SU(2)`$invariant metrics generically reduce to a Painlevé VI type equation.
###### Remark 3.3
If $`z=\zeta `$ is a solution of the equation then $`z=1/\overline{\zeta }`$ is also a solution. Therefore the equation is compatible with the real structure of twistor space.
###### Remark 3.4
The idea of Hitchin is that the lifted action of $`SU(2)`$ on the twistor space $`Z`$ gives a homomorphism of vector bundles $`\alpha :Z\times 𝔰𝔲(2)^{}TZ`$, and the inverse of $`\alpha `$ gives a flat meromorphic $`SL(2,)`$-connection, which determine isomonodromic deformations. we can think that one-forms $`\mathrm{\Theta }_1,\mathrm{\Theta }_2,\mathrm{\Theta }_3`$ on $`Z`$ are infinitesimal variations, therefore we can identify $`\mathrm{\Sigma }`$ with $`\alpha ^1`$.
###### Lemma 3.5
Let $`g`$ be a non-diagonal $`SU(2)`$-invariant metric. Then (54) has two solutions of order two if and only if there exists a function $`f(t)`$ satisfying
$`X_1^2=`$ $`4(f\alpha _2)(f\alpha _3),`$
$`X_2^2=`$ $`4(f\alpha _3)(f\alpha _1),`$
$`X_3^2=`$ $`4(f\alpha _1)(f\alpha _2).`$
And then the anti-self-dual equation reduce to (1), (2) and $`\dot{f}=f^2`$.
Proof.
We can write (54) as
$`\overline{a}z^4\overline{b}z^3+cz^2+bz+a=0,`$ (55)
where $`a,b`$ are complex coefficient and $`c`$ is a real coefficient. By an linear fractional transformation
$`z{\displaystyle \frac{\left(b|b|\right)\zeta b+|b|}{\left(\overline{b}+|b|\right)\zeta \overline{b}+|b|}}`$ (56)
preserving the real structure, we can write (54) as
$`\zeta ^4\overline{b}_0\zeta ^3+c_0\zeta ^2+b_0\zeta +1=0,`$ (57)
where $`b_0`$ is a complex coefficient and $`c_0`$ is a real coefficient. Since this equation is also compatible with the real structure, if $`\zeta =\zeta _0`$ is a solution of order two then $`\zeta =1/\overline{\zeta }_0`$ is also a solution of order two. Therefore
$`\zeta ^4\overline{b}_0\zeta ^3+c_0\zeta ^2+b_0\zeta +1=(\zeta \zeta _0)^2(\zeta +1/\overline{\zeta }_0)^2,`$ (58)
then we have $`\zeta _0^2\left(1/\overline{\zeta }_0\right)^2=1`$ and then $`\zeta _0=\pm \overline{\zeta }_0`$, which implies $`\zeta _0`$ is real or pure-imaginary. Therefore $`b_0`$ must be real or pure-imaginary. Calculating this condition, we have the Lemma.
## 4 The Hermitian Structure
Hitchin shows that if a metric is scalar-flat Kähler but not Hyper-Kähler, then the anti-self-dual equation reduce to a Painleve III type equation. We can interprets this result as the following result.
###### Corollary 4.1
If a metric is scalar-flat-Kähler but not hyper-Kähler then the equation (54) has two double zeros.
Therefore we analyze the case (54) has two double zeros.
Let $`z=z(t)`$ is a solution of (54). If we restrict $`(1,0)`$-forms $`\mathrm{\Theta }_1,\mathrm{\Theta }_2`$ on $`Z`$ to $`z=z(t)`$, we have $`(1,0)`$-forms on $`M`$, which determine an almost complex structure on $`M`$. Analyzing this almost complex structure, we have the following theorem.
###### Theorem 4.2
Let $`g`$ be an $`SU(2)`$-invariant anti-self-dual scalar-flat metric. There exists a $`SU(2)`$-invariant hermitian structure $`(g,I)`$ if and only if (54) has solutions of order two.
proof.
Let $`(g,i)`$ be a $`SU(2)`$-invariant hermitian structure. The complex structure $`I`$ is determined by $`(1,0)`$-forms $`\mathrm{\Theta }_1|_{z=z(t)}`$ and $`\mathrm{\Theta }_2|_{z=z(t)}`$, where $`z=z(t)`$ is a function on $`M`$ depending on $`t`$ only. Since the complex structure is integrable, $`\mathrm{\Theta }_3|_{z=z(t)}0`$ $`(\text{mod}\mathrm{\Theta }_1|_{z=z(t)},\mathrm{\Theta }_2|_{z=z(t)})`$. Therefore we have
$$\begin{array}{c}dz|_{z=z(t)}=\{\frac{1}{4}(\alpha _2+\alpha _3+\sqrt{1}X_1)z^3\hfill \\ \hfill \frac{1}{2}\left(\frac{w_1^2}{w_3^2w_1^2}X_2+\sqrt{1}\frac{w_1^2}{w_1^2w_2^2}X_3\right)z^2+\frac{\sqrt{1}}{2}X_1z\\ \hfill \frac{1}{2}\left(\frac{w_1^2}{w_3^2w_1^2}X_2\sqrt{1}\frac{w_1^2}{w_1^2w_2^2}X_3\right)\\ \hfill +\frac{1}{4}(\alpha _2\alpha _3+\sqrt{1}X_1)z^3\}dt\end{array}$$
(59)
On the other hand, since $`\mathrm{\Theta }_3|_{z=z(t)}0`$, $`z=z(t)`$ is a solution of (54). Moreover if we substitute $`z=z(t)`$ and (59) into the derivative of left hand side of (54), it also becomes zero. Therefore (54) has solutions of order two.
Conversely, let $`z=z_0`$ be a solution of order two, then from lemma 3.1 we have
$`z_0={\displaystyle \frac{X_2X_3\pm \sqrt{X_2^2X_3^2+X_3^2X_2^2+X_1^2X_2^2}}{X_1(X_2\sqrt{1}X_3)}},`$ (60)
if $`X_1X_2X_30`$. And then we have $`\mathrm{\Theta }_3|_{z=z(t)}0`$. Therefore the almost complex structure determined by the $`(1,0)`$-forms $`\mathrm{\Theta }_1|_{z=z(t)}`$ and $`\mathrm{\Theta }_2|_{z=z(t)}`$ is integrable.
If $`X_1X_2X_3=0`$, $`f`$ must be $`\alpha _1,\alpha _2`$ or $`\alpha _3`$. Let $`f=\alpha _1`$, then we have $`X_2=0`$ and $`X_3=0`$, and then
$`z_0={\displaystyle \frac{\sqrt{\alpha _3\alpha _1}+\sqrt{1}\sqrt{\alpha _2\alpha _1}}{\sqrt{\alpha _2+\alpha _3+2\alpha _1}}},`$ (61)
and then $`\mathrm{\Theta }_3|_{z=z(t)}0`$. In this case the almost complex structure is also integrable.
###### Theorem 4.3
The hermitian structure $`(g,I)`$ determined by theorem 4.2 is Kähler if and only if
$`X_1^2`$ $`=4\alpha _2\alpha _3,`$ $`X_2^2`$ $`=4\alpha _3\alpha _1,`$ $`X_3^2`$ $`=4\alpha _1\alpha _2.`$ (62)
proof.
If $`X_1X_2X_30`$, the Kähler form is determined by (60) as
$`\mathrm{\Omega }=`$ $`{\displaystyle \frac{X_2X_3}{\sqrt{X_2^2X_3^2+X_3^2X_2^2+X_1^2X_2^2}}}\mathrm{\Omega }_1^+`$
$`+{\displaystyle \frac{X_3X_1}{\sqrt{X_2^2X_3^2+X_3^2X_2^2+X_1^2X_2^2}}}\mathrm{\Omega }_2^+`$
$`+{\displaystyle \frac{X_1X_2}{\sqrt{X_2^2X_3^2+X_3^2X_2^2+X_1^2X_2^2}}}\mathrm{\Omega }_3^+.`$
By the anti-self-dual equations (1),(2),(3), we have
$`d\mathrm{\Omega }=`$ $`{\displaystyle \frac{2fw_1X_2X_3}{\sqrt{X_2^2X_3^2+X_3^2X_2^2+X_1^2X_2^2}}}dt\stackrel{~}{\sigma }_2\stackrel{~}{\sigma }_3`$
$`+{\displaystyle \frac{2fw_2X_3X_1}{\sqrt{X_2^2X_3^2+X_3^2X_2^2+X_1^2X_2^2}}}dt\stackrel{~}{\sigma }_3\stackrel{~}{\sigma }_1`$
$`+{\displaystyle \frac{2fw_3X_1X_2}{\sqrt{X_2^2X_3^2+X_3^2X_2^2+X_1^2X_2^2}}}dt\stackrel{~}{\sigma }_1\stackrel{~}{\sigma }_2.`$
Since $`w_1w_2w_30`$ and $`X_1X_2X_30`$, we have $`d\mathrm{\Omega }=0`$ if and only if $`f=0`$.
If $`X_1X_2X_3=0`$, then $`f`$ must be $`\alpha _1,\alpha _2`$ or $`\alpha _3`$. Let $`f=\alpha _1`$, then $`X_1^2=4(\alpha _2\alpha _1)(\alpha _3\alpha _1)`$, $`X_2=0`$, $`X_3=0`$. The Kähler form is determined by (61) as
$`\mathrm{\Omega }={\displaystyle \frac{\sqrt{\alpha _2\alpha _1}}{\sqrt{\alpha _2+\alpha _32\alpha _1}}}\mathrm{\Omega }_2^++{\displaystyle \frac{\sqrt{\alpha _3\alpha _1}}{\sqrt{\alpha _2+\alpha _32\alpha _1}}}\mathrm{\Omega }_3^+.`$ (63)
Then
$`d\mathrm{\Omega }={\displaystyle \frac{2w_2\alpha _1\sqrt{\alpha _2\alpha _1}}{\sqrt{\alpha _2+\alpha _32\alpha _1}}}dt\stackrel{~}{\alpha }_3\stackrel{~}{\alpha }_1+{\displaystyle \frac{2w_3\alpha _1\sqrt{\alpha _3\alpha _1}}{\sqrt{\alpha _2+\alpha _32\alpha _1}}}dt\stackrel{~}{\alpha }_1\alpha _2.`$ (64)
Since the metric is non-diagonal, $`X_1^2=4(\alpha _2\alpha _1)(\alpha _3\alpha _1)0`$ and then $`d\mathrm{\Omega }=0`$ if and only if $`\alpha _1=0`$.
###### Remark 4.4
If the metric is scalar-flat Kähler, the anti-self-dual equation reduce to a sixth-order equation.
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# Quantum gauge symmetry from finite field dependent BRST transformations
## Abstract
Using the technique of finite field dependent BRST transformations we show that the classical massive Yang-Mills theory and the pure Yang-Mills theory whose gauge symmetry is broken by a gauge fixing term are identical from the view point of quantum gauge symmetry. The explicit infinitesimal transformations which leave the massive Yang-Mills theory BRST invariant are given.
In a recent paper it was shown that a classical massive gauge theory does not have an essential difference, at the quantum level, from a gauge invariant theory whose gauge symmetry is broken by a gauge fixing term. Specifically, the classical lagrangians,
$$=_{YM}\frac{m^2}{2}A_\mu ^aA_\mu ^a$$
(1)
and
$$=_{YM}\frac{1}{2}(_\mu A_\mu ^a)^2$$
(2)
where $`_{YM}`$ is the Yang-Mills lagrangian,
$$_{YM}=\frac{1}{4}(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)^2$$
(3)
could be given an identical physical meaning, both representing an effective gauge fixed lagrangian associated with the quantum theory defined by
$$𝒟A_\mu ^a𝒟B^a𝒟\overline{c}^a𝒟c^a\mathrm{exp}\left\{S_{YM}+d^4x\left[iB^a(^\mu A_\mu ^a)+\overline{c}^a(_\mu (D^\mu c)^a)\right]\right\}$$
(4)
that is invariant under the BRST transformations,
$`\delta A_\mu ^a`$ $`=`$ $`i(D_\mu c)^a\delta \lambda `$ (5)
$`\delta c^a`$ $`=`$ $`i{\displaystyle \frac{g}{2}}f^{abc}c^bc^c\delta \lambda `$ (6)
$`\delta \overline{c}^a`$ $`=`$ $`B^a\delta \lambda `$ (7)
$`\delta B^a`$ $`=`$ $`0`$ (8)
where $`\delta \lambda `$ is an infinitesimal Grassmann parameter. This is to be contrasted with the conventional interpretation of regarding (1) as a massive vector theory and (2) as an effective Yang-Mills theory with a covariant gauge fixing term.
In this paper we shall show the equivalence of the quantum theories defined by (1) and (2) by following the method of finite field dependent BRST (FFBRST) transformations developed by one of us . In particular this method, which will be briefly reviewed below, connects quantum gauge theories in different gauges. Here we start from the conventional gauge fixed Yang-Mills lagrangian defined by (2). The explicit FFBRST transformations are then stated which maps this theory to one whose lagrangian is defined by (1), thereby showing the connection between them. We also get the form of the transformations that preserve the BRST invariance of the quantum theory defined by (1). Finally we suggest a possible connection between our approach and that adopted in , which was based on a modified quantization scheme , where the variation of the gauge field in the path integral is taken over the entire gauge orbit.
Let us now briefly review the FFBRST approach . FFBRST transformations are obtained by an integration of infinitesimal (field dependent ) BRST transformations . In this method all the fields are function of some parameter, $`\kappa :0\kappa 1`$. For a generic field $`\varphi (x,\kappa ),\varphi (x,\kappa =0)=\varphi (x)`$ and $`\varphi (x,\kappa =1)=\varphi ^{}(x).`$ Then the infinitesimal field dependent BRST transformations are defined as,
$$\frac{d}{d\kappa }\varphi (x,\kappa )=\delta _{BRST}\varphi (x,\kappa )\mathrm{\Theta }^{}[\varphi (x,\kappa )]$$
(9)
where $`\mathrm{\Theta }^{}d\kappa `$ is an infinitesimal field dependent parameter. It has been shown by integrating these equations from $`\kappa =0`$ to $`\kappa =1`$ that $`\varphi ^{}(x)`$ are related to $`\varphi (x)`$ by FFBRST,
$$\varphi ^{}(x)=\varphi (x)+\delta _{BRST}\varphi (x)\mathrm{\Theta }[\varphi (x)]$$
(10)
where $`\mathrm{\Theta }[\varphi (x)]`$ is obtained from $`\mathrm{\Theta }^{}[\varphi (x)]`$ through the relation,
$$\mathrm{\Theta }[\varphi (x)]=\mathrm{\Theta }^{}[\varphi (x)]\frac{\mathrm{exp}f[\varphi (x)]1}{f[\varphi (x)]}$$
(11)
and $`f`$ is given by $`f=_i\frac{\delta \mathrm{\Theta }^{}(x)}{\delta \varphi _i(x)}\delta _{BRST}\varphi _i(x)`$
The choice of the parameter $`\mathrm{\Theta }^{}`$ is crucial in connecting different effective gauge theories by means of the FFBRST. In particular the FFBRST of Eq. (10) with $`\mathrm{\Theta }^{}[\varphi (x,\kappa )]=i\overline{c}^a(y)\left[F^a[A(\kappa )]F^a[A(\kappa )]\right]`$ relates the Yang-Mills theory with an arbitrary gauge fixing $`F[A]`$ to the Yang-Mills theory with another arbitrary gauge fixing $`F^{}[A]`$ .
The meaning of these field transformations is as follows. We consider the vacuum expectation value of a gauge invariant functional $`G[\varphi ]`$ in some arbitrary gauge $`F[A]`$,
$$<<G[\varphi ]>>𝒟\varphi G[\varphi ]\mathrm{exp}(iS_{eff}^F[\varphi ])$$
(12)
where,
$$S_{eff}^F=S_0\frac{1}{2}d^4xF^2[A]d^4x\overline{c}^aW^{ab}c^b$$
(13)
with
$$W^{ab}=\frac{\delta F^a}{\delta A_\mu ^c}D_\mu ^{cb}[A]$$
(14)
Here $`S_0`$ is the pure Yang-Mills action obtained from (3) and the covariant derivative, $`D_\mu ^{ab}[A]\delta ^{ab}_\mu +gf^{abc}A_\mu ^c`$. For simplicity we have set the gauge parameter $`\lambda =1`$ in the gauge fixing term $`\frac{1}{2\lambda }d^4xF^2[A]`$.
Now we perform the FFBRST transformations $`\varphi \varphi ^{}`$ given by (10). We have then
$$<<G[\varphi ]>>=<<G[\varphi ^{}]>>=𝒟\varphi ^{}J[\varphi ^{}]G[\varphi ^{}]\mathrm{exp}(iS_{eff}^F[\varphi ^{}])$$
(15)
on account of BRST invariance of $`S_{eff}^F`$ and gauge invariance of $`G[\varphi ]`$. Here $`J[\varphi ^{}]`$ is the Jacobian associated with FFBRST and defined as,
$$𝒟\varphi =𝒟\varphi ^{}J[\varphi ^{}]$$
(16)
As shown in for the special case $`G[\varphi ]=1`$ the Jacobian $`J[\varphi ^{}]`$ in Eq (15) can always be replaced by $`\mathrm{exp}(iS_1[\varphi ^{}]`$) with,
$$S_{eff}^F[\varphi ^{}]+S_1[\varphi ^{}]=S_{eff}^F^{}[\varphi ^{}]$$
(17)
where
$$S_{eff}^F^{}=S_0\frac{1}{2}d^4xF^2[A]d^4x\overline{c}^aW^{ab}c^b$$
(18)
with
$$W^{ab}=\frac{\delta F^a}{\delta A_\mu ^c}D_\mu ^{cb}[A]$$
(19)
The extra piece in the action which arises from the Jacobian of such FFBRST is given by,
$$S_1[\varphi ]=d^4x\left[\frac{1}{2}F^2[A]+\frac{1}{2}F^2[A]+\overline{c}[WW^{}]c\right]$$
(20)
Thus the FFBRST in Eq. (10) takes the theory with gauge $`F`$ to the corresponding theory with gauge $`F^{}`$.
We are now ready to apply this machinery to the present problem. We start with the generating functional for the Yang-Mills theory in the Lorentz gauge,
$$Z=𝒟A_\mu 𝒟c𝒟\overline{c}\mathrm{exp}(iS_{eff}^L)$$
(21)
where
$$S_{eff}^L=S_0\frac{1}{2}d^4x(^\mu A_\mu )^2d^4x\overline{c}Wc$$
(22)
with $`W=^\mu D_\mu `$ is the Faddeev-Popov determinant. We now apply FFBRST \[Eq. (10) \] with
$$\mathrm{\Theta }^{}=id^4y\overline{c}^a(y)\left[^\mu A_\mu ^am\frac{\omega ^\mu }{|\omega |}A_\mu ^a\right](y)$$
(23)
where $`\omega ^\mu `$ is an arbitrary 4-vector, to the expression for the generating functional to obtain,
$$Z=𝒟A_\mu ^{}𝒟c^{}𝒟\overline{c}^{}\mathrm{exp}i(S_{eff}^L+S_1)$$
(24)
The additional piece in the action comes from the non-trivial Jacobian of the FFBRST and can be written using Eq (20)
$$S_1=d^4x\left[\frac{1}{2|\omega |^2}m^2(\omega ^\mu A_\mu )^2+\frac{1}{2}(^\mu A_\mu )^2\overline{c}(W^{}W)c\right]$$
(25)
with $`W^{}=m\frac{\omega ^\mu }{|\omega |}D_\mu `$. Hence we obtain the generating functional for a new effective action given by,
$$S_{eff}=S_0d^4x\left[\frac{1}{2}A_\mu M^{\mu \nu }A_\nu +\overline{c}m\frac{\omega ^\mu }{|\omega |}D_\mu c\right]$$
(26)
where $`M^{\mu \nu }`$ is a generalized mass matrix,
$$M^{\mu \nu }=m^2\frac{\omega ^\mu \omega ^\nu }{|\omega |^2}$$
(27)
This effective action (26) corresponds to the Yang-Mills lagrangian with a generalized mass term. It shows the connection between the Lorentz gauge and a generalized ‘mass’ gauge $`\frac{1}{2}A_\mu M^{\mu \nu }A_\nu `$ in the context of Yang-Mills theory. To exactly reproduce the familiar mass term, we restrict the arbitrary vector $`\omega ^\mu `$ to be of infinitesimal form satisfying the symmetric multiplication rule,
$$\frac{\omega ^\mu \omega _\nu }{|\omega |^2}=\frac{g_\nu ^\mu }{4}$$
(28)
In that case the gauge fixing term is $`\frac{1}{8}m^2A_\mu A^\nu `$ which coincides with the standard mass term, after a proper normalization of $`m`$.
The infinitesimal BRST transformations which leave the Yang-Mills theory with a mass term (26) invariant are given by
$`\delta A_\mu ^a`$ $`=`$ $`D_\mu ^{ab}c^b\delta \lambda `$ (29)
$`\delta c^a`$ $`=`$ $`{\displaystyle \frac{g}{2}}f^{abc}c^bc^c\delta \lambda `$ (30)
$`\delta \overline{c}^a`$ $`=`$ $`m{\displaystyle \frac{\omega ^\mu }{|\omega |}}A_\mu ^a\delta \lambda `$ (31)
We have shown how, by means of finite field dependent BRST transformations, it was possible to interpolate between the Yang-Mills theory in the covariant gauge to the Yang-Mills theory in a mass like gauge. Since FFBRST also connects the Yang-Mills theory in the axial and covariant gauges it is clear that the Yang-Mills theory with a mass like gauge fixing term can also be obtained from other starting points. In this paper we took the covariant gauge as the starting point for reasons of convenience and also comparing our analysis with . It may be pointed out that the latter approach is based on the variation of the gauge variable along the entire gauge orbit, without taking any specific limit of the gauge fixing parameter. Consequently there seems to be a connection between this approach and the FFBRST method, which is not altogether surprising. Carrying out the integration over the complete gauge orbit would be simulated by finite BRST transformations instead of the conventional infinitesimal one. We feel it might be useful to pursue this connection in a later work.
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# DEFINING ENTROPY BOUNDS Alberta-Thy-09-00, hep-th/0007238
## 1 Introduction
Bekenstein has conjectured that the entropy $`S`$ of a system confined to radius $`R`$ or less and energy $`E`$ or less would obey the inequality (using units $`\mathrm{}=c=k_{\mathrm{B}oltzmann}=1`$)
$$S2\pi ER.$$
(1)
He and colleagues have supported this conjecture with many arguments and examples \[1-18\]. However, many counterarguments and counterexamples have also been noted \[19-34\]. Whether the conjectured bound (1) holds or not depends on what systems are considered and how $`R`$, $`E`$, and $`S`$ are defined.
Perhaps the simplest procedure would be to just consider quantum fields inside some bounded region within a sphere of radius $`R`$ and put boundary conditions on the fields at the boundary of the region. However, this procedure leads to a large number of counterexamples to Bekenstein’s conjectured bound. For example , the Casimir effect can make $`E<0`$ for certain states of quantum fields confined within a certain regions of radius $`R`$, violating the bound. If states with $`E<0`$ are excluded by definition, one can still consider a mixed state with arbitrarily small positive $`E`$ that violates the bound. Even if $`E`$ is redefined to be the nonnegative energy excess over that of the ground state , one can violate the bound by a mixed state that is almost entirely the ground state and a tiny incoherent mixture of excited states, at least if the entropy is defined to be $`S=tr\rho \mathrm{ln}\rho `$ . If $`S`$ is instead defined to be $`S=\mathrm{ln}n`$ for a mixture of $`n`$ orthogonal pure states (which would agree with $`S=tr\rho \mathrm{ln}\rho `$ if the mixture had equal probabilities $`1/n`$ for each of those $`n`$ pure states), then one can violate the bound by an equal mixture of the ground state and the first excited state of certain interacting fields with certain boundary conditions that have the two lowest states nearly degenerate in energy (separated by exponentially small tunneling effects) . If interacting fields are excluded from the definition of allowable systems, one can get a violation by considering a sufficiently large number $`N`$ of identical free fields, giving $`n=N`$ degenerate first excited states of finite energy but sufficiently large entropy $`S=\mathrm{ln}n=\mathrm{ln}N`$ to violate (1) . And even for a single free electromagnetic field, $`S=\mathrm{ln}n`$ can exceed $`2\pi ER`$ by an arbitrarily large factor by using boundary conditions corresponding to an arbitrarily large number of parallel perfectly conducting plates within the region of radius $`R`$ , or by using boundary conditions corresponding to an arbitrarily long coaxial cable loop coiled up within the region .
However, other than in his papers with Schiffer , Bekenstein has generally advocating taking $`E`$ to be the total energy of a complete system . This would disallow using just the energy of fields within a bounded region with boundary conditions, since that would ignore the energy of whatever it is that is providing the boundary conditions. Therefore, all of the counterexamples mentioned above would be excluded by this restriction. However, then the problem is to define what one means by the radius $`R`$ of the system. In the weakly gravitating case (essentially quantum fields in flat Minkowski spacetime) that we shall focus on here, Bekenstein takes $`R`$ to mean the radius of a sphere which circumscribes the system, which leaves the problem of what it means for a sphere to circumscribe the complete system.
In quantum field theory in Minkowski spacetime, the complete system is the quantum state of the fields. Since the quantum fields extend all the way out to radial infinity, a sphere circumscribing the entire system would have to be at $`R=\mathrm{}`$, which makes the Bekenstein bound true (at least for states of positive energy and finite entropy) but trivial. To get a nontrivial bound, one needs to suppose that a sphere of finite $`R`$ can circumscribe the system. For example, one might try to say that the sphere encloses all of the excitations of the fields from the vacuum. However, it is also hard to get this to occur for a finite $`R`$. For example, the wavefunction for any single particle state that is a superposition of energy eigenstates of bounded energy will not vanish outside any finite radius $`R`$, since a one-particle wavefunction that does vanish outside a finite region must be a superposition of arbitrarily large momentum components, which will have unbounded energy. Even if one looks at a composite system, such as a hydrogen atom, and ignores the fact that its center of mass will have amplitudes to be outside any finite sphere if it is made of purely bounded energy components, the wavefunction for the relative position of the electron and proton does not drop identically to zero outside any finite separation distance for states that are superpositions of energy eigenstates of bounded energy. In particular, even if one fixed the center of mass of a hydrogen atom in its ground state and ignored the infinite energy from the resulting infinite uncertainty of the center of mass momentum, the density matrix for the electron position would decay only exponentially with distance from the center of mass and never go to zero outside any sphere of finite radius $`R`$.
Therefore, it is problematic to define the radius $`R`$ of a sphere circumscribing a complete system in any quantum field theory. This issue has not been addressed by Bekenstein and his collaborators, but without such a definition, there is no nontrivial formulation of the conjectured bound (1) for complete systems, but only its trivial truth for any complete system with positive energy and finite entropy that can only be circumscribed by the sphere enclosing all of space, $`R=\mathrm{}`$.
Here new ways are proposed to define systems and their radii $`R`$, energies $`E`$, and entropies $`S`$, so that for each, there is a bound on $`S`$ for a given system as a function of finite $`R`$ and $`E`$. These bounds will not have the form of Bekenstein’s conjectured inequality (1), though in some cases they may obey that inequality.
## 2 Vacuum-Outside-R States
The main new element of the present paper is a proposal is to define a system of radius $`R`$ (in flat spacetime for the present) not by imposing boundary conditions on the field itself, but by imposing conditions on the quantum state of the field so that outside a closed ball of radius $`R`$ the quantum state is indistinguishable from the vacuum at some time. Such a state will be called a vacuum-outside-$`R`$ state. (For simplicity, set this time to be $`t=0`$, and take the closed ball, say $`B`$, to be the region $`rR`$ on the $`t=0`$ hypersurface, where $`r`$ is the standard radial polar coordinate giving the proper distance from the coordinate origin on that hypersurface.) In other words, a vacuum-outside-$`R`$ state of the system, say as expressed by its density matrix $`\rho `$, is such that the expectation value of any operator $`O`$ which is completely confined to the region $`r>R`$ when written in terms of field and conjugate operators at $`t=0`$, is precisely the same as the expectation value of the same operator in the vacuum state $`|0><0|`$,
$$tr(O\rho )=<0|O|0>.$$
(2)
In particular, all the $`n`$-point functions for the field and for its conjugate momentum in the state $`\rho `$ are the same as in the vacuum state, if all of the $`n`$ points are outside the ball of radius $`R`$ and on the hypersurface $`t=0`$. Of course, the $`n`$-point functions need not be the same as their vacuum values if some or all of the points are inside the ball.
If operators confined to the three-dimensional region $`r>R`$ and $`t=0`$ (say $`C`$, to give a name to this achronal spacelike surface, the $`t=0`$ hypersurface with the central closed ball $`B`$, $`rR`$, excluded) have the same expectation value in the vacuum-outside-$`R`$ state as in the vacuum state, the same will be true in any quantum field theory that I shall call “strongly causal” for all operators confined to the Cauchy development or domain of dependence of $`C`$, the larger four-dimensional region $`r>R+|t|`$ (say $`D`$) that is the set of all points in the Minkowski spacetime such that every inextendible (endless) causal, or non-spacelike (everywhere timelike or lightlike), curve through such a point intersects the partial Cauchy surface $`C`$. Just as solutions of hyperbolic wave equations in $`D`$ are determined by the data on $`C`$, so the part of the quantum state of a strongly causal field in $`D`$, as represented by the expectation values of operators confined to $`D`$, is determined by the part of the quantum state in $`C`$, as represented by the expectation values of operators confined to $`C`$. (For some interacting quantum field theories, the expectation values of operators confined to the three-dimensional spacelike surface $`C`$ may be too ill-defined for these theories to be “strongly causal” in my sense, but a wider class of these theories may be “weakly causal” in the sense that sufficiently many operators smeared over, but confined to, an arbitrarily thin-in-time four-dimensional slab, say $`E`$, containing $`C`$ within $`D`$, have well-defined expectation values that determine the expectation values of all operators smeared over, but confined to, any part of $`D`$.)
Henceforth I shall restrict attention to strongly causal and weakly causal quantum field theories, calling them simply causal quantum field theories for short. I shall also assume, until discussing gravitational theories later, that any quantum field theory under consideration is a nongravitational Lorentz-invariant quantum field theory in Minkowski spacetime, and that it has a unique pure state of lowest Minkowski energy $`E=0`$ (the expectation value of the Hamiltonian $`H`$ that generates translations in the time coordinate $`t`$ in some Lorentz frame, with the arbitrary constant in the Hamiltonian being adjusted to give the lowest energy state zero energy).
Therefore, for such a causal nongravitational quantum field theory in Minkowski spacetime, I shall propose that the radius $`R`$ be defined so that all of the operators constructed from field and conjugate momentum operators smeared over regions confined to the region $`D`$, $`r>R+|t|`$ in some Lorentz frame, have in the particular quantum state being considered (a vacuum-outside-$`R`$ state) the same expectation values that they have in the vacuum state for that quantum field theory. The energy $`E`$ of the state $`\rho `$ can then be simply defined to be the expectation value,
$$Etr(H\rho ),$$
(3)
of the Hamiltonian $`H`$ that generates time translations in the same Lorentz frame. Because the energy $`E`$ has been defined to have the minimum value of zero for the unique pure vacuum state, there is no problem here with negative Casimir energies. In other words, the energy is that of the complete system over all of Minkowski spacetime.
Obviously we would also like a definition of the entropy $`S`$ that has a minimum value of zero, which it should attain for the pure vacuum state. One simple definition is the von Neumann entropy,
$$S=S_{\mathrm{vN}}tr\rho \mathrm{ln}\rho ,$$
(4)
using the density matrix $`\rho `$ for the full state of the quantum field, over the entire Minkowski spacetime.
## 3 Entropy Bounds for Vacuum-Outside-R States
Now we may conjecture that for any vacuum-outside-$`R`$ state of any particular causal nongravitational quantum field theory in Minkowski spacetime, one which has the vacuum expectation values in the region $`D`$, $`r>R+|t|`$ (the region causally disconnected from the ball $`rR`$ at $`t=0`$), the von Neumann entropy is bounded above by some function $`\sigma _{\mathrm{vN}}`$ (depending on the quantum field theory in question) of the radius $`R`$ and energy $`E`$:
$$S_{\mathrm{vN}}\sigma _{\mathrm{vN}}(R,E).$$
(5)
Define this function $`\sigma _{\mathrm{vN}}(R,E)`$ to be the least upper bound on the von Neumann entropy of any state which is vacuum outside the radius $`R`$ and which has energy $`E`$.
In the case of a scale-invariant quantum field, such as a free massless field, or say a massless scalar field $`\varphi `$ with a $`\lambda \varphi ^4`$ self-coupling potential, the least upper bound function $`\sigma _{\mathrm{vN}}(R,E)`$ will actually be a function of the single dimensionless variable
$$x2\pi RE,$$
(6)
say
$$\sigma _{\mathrm{vN}}(R,E)=\sigma _\mathrm{N}(x).$$
(7)
Bekenstein’s conjectured entropy bound (1), if $`R`$, $`E`$, and $`S`$ were defined as done herein, would be $`\sigma _{\mathrm{vN}}(R,E)x`$, whether or not the quantum field theory is scale invariant, or
$$B_{\mathrm{vN}}(R,E)\frac{\sigma _{\mathrm{vN}}(R,E)}{x}\frac{\sigma _{\mathrm{vN}}(R,E)}{2\pi RE}1.$$
(8)
If the quantum field theory is scale invariant, we can define
$$B_\mathrm{N}(x)\frac{\sigma _\mathrm{N}(x)}{x},$$
(9)
which should also be less than or equal to unity if Bekenstein’s bound applies.
For a set of one or more free massless fields and vacuum-outside-$`R`$ states with $`x1`$, one would expect that the highest entropy would be given by a mixed state that at $`t=0`$ is approximately a high-temperature ($`RT1`$) thermal radiation state for $`r<R`$, surrounded by vacuum for $`r>R`$. A high-temperature thermal radiation state has an energy density for massless fields of approximately $`a_rT^4`$, and hence an entropy density $`(4/3)a_rT^3`$, where
$$a_r=\frac{\pi ^2}{30}(n_b+\frac{7}{8}n_f)$$
(10)
is the radiation constant for $`n_b`$ independent bosonic degrees of freedom for each momentum (e.g., $`n_b`$ different spin or helicity states) and for $`n_f`$ fermionic degrees of freedom. Therefore, in this case with $`x1`$,
$$B_{\mathrm{vN}}(R,E)=\frac{\sigma _\mathrm{N}(x)}{x}\left(\frac{2^7a_r}{3^5\pi ^2x}\right)^{\frac{1}{4}}=\left[\frac{2^6}{3^65x}(n_b+\frac{7}{8}n_f)\right]^{\frac{1}{4}},$$
(11)
which is indeed less than 1, thus obeying Bekenstein’s conjectured bound, for
$$x\frac{2^7a_r}{3^5\pi ^2}=\frac{2^6}{3^65}(n_b+\frac{7}{8}n_f)=\frac{64n_b+56n_f}{3645},$$
(12)
if $`x`$ is also large enough that Eq. (11) is a good approximation. Thus one would expect that Bekenstein’s conjectured bound, using the definitions above for $`R`$, $`E`$, and $`S`$, holds for a fixed set of free massless quantum fields at sufficiently large $`x2\pi RE`$.
On the other hand, the definitions above for $`R`$, $`E`$, and $`S`$ still permit Bekenstein’s conjectured bound applied to them to be violated for sufficiently small $`x`$, as we can see by the following construction:
A way to construct vacuum-outside-$`R`$ states, quantum states of a free quantum field theory that have vacuum expectation values in the region $`D`$, $`r>R+|t|`$, is to apply to the vacuum state unitary operators constructed from fields and/or conjugate momenta smeared within the region $`r<R`$ at $`t=0`$. In particular, if $`h`$ is an hermitian operator constructed from fields and/or conjugate momenta smeared within $`r<R`$ at $`t=0`$, then $`U=e^{ih}`$ is such a unitary operator, and
$$|\psi >=U|0>=e^{ih}|0>$$
(13)
is a pure quantum state that has precisely the vacuum expectation values in the region $`D`$. This result can be seen formally from the fact that any operator $`O`$ confined to the region $`D`$ (the four-dimensional region $`r>R+|t|`$) that is causally disconnected from the ball $`B`$ (the three-dimensional region $`rR`$ on the $`t=0`$ hypersurface) commutes with the operators $`h`$ and $`U`$ that are confined to that hypersurface, $`[O,h]=[O,U]=0`$, so
$$tr(O\rho )=<\psi |O|\psi >=<0|U^1OU|0>=<0|U^1UO|0>=<0|O|0>.$$
(14)
If $`\{h_i\}`$ is a set of hermitian operators that each are confined to the ball $`B`$ (i.e., are constructed from fields and momenta that are smeared only over that region), and if $`\{q_i\}`$ is a set of positive numbers that sum to unity, then
$$\rho =\underset{i}{}q_ie^{ih_i}|0><0|e^{ih_i}$$
(15)
is a more general vacuum-outside-$`R`$ state, since this density matrix gives vacuum expectation values, $`tr(O\rho )=<0|O|0>`$, for any operator $`O`$ confined to the region $`D`$ that is causally disconnected from $`B`$ (i.e., having no causal curves, either timelike or lightlike, intersecting both the ball $`B`$ of $`rR`$ at $`t=0`$ and the region $`D`$ with $`r>R+|t|`$).
Here let us consider the simple example in which $`i`$ takes only the two values 1 and 2, and $`h_1=0`$ and $`h_2=h`$. Let $`q_1=1q`$ and $`q_2=q`$, and let
$$e^{ih}|0>=U|0>=|\psi >=c|0>+s|1>$$
(16)
in terms of a decomposition of $`|\psi >`$ into the two orthonormal states $`|0>`$ and $`|1>=(|\psi ><0|\psi >|0>)/\sqrt{1|<0|\psi >|^2}`$, so
$$c=<0|\psi >=<0|U|0>=<0|e^{ih}|0>,$$
(17)
$$s=\sqrt{1|<0|U|0>|^2}=\sqrt{1|c|^2}.$$
(18)
Note that $`|1><1|`$ by itself is not generically a vacuum-outside-$`R`$ state.
Now Eq. (15) gives the density matrix as
$`\rho `$ $`=`$ $`(1q)|0><0|+q|\psi ><\psi |`$ (19)
$`=`$ $`(1qs^2)|0><0|+qcs|0><1|+q\overline{c}s|1><0|+qs^2|1><1|,`$
a density matrix in the two-dimensional space of pure states spanned by the two orthonormal pure states $`|0>`$ and $`|1>`$. The two eigenvalues of this density matrix are, say, $`p`$ and $`1p`$ (since their sum is $`tr\rho =1`$), with product
$`yp(1p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{[p+(1p)]^2[p^2+(1p)^2]\}={\displaystyle \frac{1}{2}}\{[tr(\rho )]^2[tr(\rho ^2)]\}`$ (20)
$`=`$ $`q(1q)(1|<0|U|0>|^2)=q(1q)s^2.`$
The expectation value of the energy of this mixed state is, since I have assumed $`H|0>=0`$,
$$E=tr(H\rho )=q<\psi |H|\psi >=q<0|U^1HU|0>.$$
(21)
Then
$$x2\pi RE=2\pi Rq<0|U^1HU|0>.$$
(22)
The von Neumann entropy of this mixed state is
$`S_{\mathrm{vN}}(y)`$ $`=`$ $`tr(\rho \mathrm{ln}\rho )=p\mathrm{ln}p(1p)\mathrm{ln}(1p)`$ (23)
$`=`$ $`{\displaystyle \frac{y}{\frac{1}{2}(1+\sqrt{14y})}}\mathrm{ln}{\displaystyle \frac{1}{y}}+\sqrt{14y}\mathrm{ln}{\displaystyle \frac{1}{\frac{1}{2}(1+\sqrt{14y})}}`$
$``$ $`y[(1+y+2y^2)\mathrm{ln}{\displaystyle \frac{1}{y}}+(1{\displaystyle \frac{1}{2}}y{\displaystyle \frac{5}{3}}y^2)],`$
a monotonically increasing function of $`yq(1q)s^21/4`$, where the last approximate equality of Eq. (23) applies for $`y1`$.
As $`q`$ and/or $`h`$ is reduced toward zero, $`x`$, $`y`$, and $`S`$ also decrease toward zero, but whereas $`x`$ and $`y`$ asymptotically decrease linearly with $`q`$, the dominant term of $`S`$ has an extra logarithmic factor that grows with the reduction of $`y`$, so the ratio,
$$B\frac{S_{\mathrm{vN}}}{x}\frac{S_{\mathrm{vN}}}{2\pi RE}\frac{y}{x}\left(\mathrm{ln}\frac{1}{y}+1\right)$$
(24)
when $`y1`$, increases without limit as $`y`$ is reduced toward zero. Therefore, when $`y`$ is made sufficiently small (e.g., by making $`q`$ sufficiently small), Bekenstein’s conjectured bound for the definition of $`R`$, $`E`$, and $`S`$ used here is violated.
## 4 Free Quantum Field Theory Examples
Let us consider a specific example for the Hermitian operator $`h`$ that is constructed from field operators confined to the ball $`B`$ of radius $`R`$ on an initial flat hypersurface of Minkowski spacetime. Take the quantum field theory to be that of a single massless scalar field operator $`\varphi `$. Consider the smeared linear Hermitian field operator
$$\chi =d^3x[F(𝐱)\varphi (t=0,𝐱)+G(𝐱)\dot{\varphi }(t=0,𝐱)],$$
(25)
where $`F(𝐱)`$ and $`G(𝐱)`$ are real functions of the spatial location $`𝐱`$ that are zero for $`|𝐱|>R`$, so that $`\chi `$ is made up of operators confined to the ball $`B`$, $`rR`$ at $`t=0`$. Then for real parameters $`\alpha `$ and $`\beta `$, let
$$h=\alpha \chi +\beta \chi ^2,$$
(26)
which is thus also an Hermitian operator confined to the ball $`B`$.
Then by expanding out $`\varphi (𝐱)`$ and $`H`$ in terms of creation and annihilation operators, one can show, after a certain amount of algebra that will not be repeated here, that
$$c<0|\psi ><0|U|0><0|e^{ih}|0>=(12i\beta X)^{1/2}\mathrm{exp}\left(\frac{\frac{1}{2}\alpha ^2X}{12i\beta X}\right),$$
(27)
$$s^21|<0|U|0>|^21|c|^2=1(1+4\beta ^2X^2)^{1/2}\mathrm{exp}\left(\frac{\alpha ^2X}{1+4\beta ^2X^2}\right),$$
(28)
and
$$<\psi |H|\psi ><0|U^1HU|0><0|e^{ih}He^{ih}|0>=\alpha ^2Y+4\beta ^2XY=(\alpha ^2X+4\beta ^2X^2)Z,$$
(29)
where
$$X<0|\chi ^2|0>=d^3xd^3y\frac{F(𝐱)F(𝐲)+G(𝐱)G(𝐲)}{4\pi ^2|𝐱𝐲|^2}$$
(30)
and
$$YXZ<0|\chi H\chi |0>=\frac{1}{2}d^3x[|F(𝐱)|^2+|G(𝐱)|^2].$$
(31)
Incidentally, I have not included individual higher powers of $`\chi `$ in $`h`$, because then expanding $`e^{ih}`$ into a power series in $`\chi `$ and taking the expectation values gives divergent series when one uses the key intermediate results
$$<0|\chi ^m|0>=\{\begin{array}{cc}(m1)!!X^{m/2}\hfill & m\text{ even}\hfill \\ 0\hfill & m\text{ odd}\hfill \end{array},$$
(32)
and
$$<0|\chi ^mH\chi ^n|0>=\{\begin{array}{cc}mn(m+n3)!!X^{(m+n)/2}Z\hfill & m+n\text{ even}\hfill \\ 0\hfill & m+n\text{ odd}\hfill \end{array}.$$
(33)
(The divergences arise from the rapid growth of the double factorials with their arguments. These double factorials arise from the counting of the number of pairings of the creation and annihilation operators in the powers of the $`\chi `$’s and in the Hamiltonian $`H`$ for the massless scalar field $`\varphi `$.)
If one takes $`F(𝐱)`$ and $`G(𝐱)`$ to be spherically symmetric, say
$$F(𝐱)=R^2f\left(\frac{|𝐱|}{R}\right)R^2f(u)$$
(34)
and
$$G(𝐱)=R^1g\left(\frac{|𝐱|}{R}\right)R^1g(u)$$
(35)
with $`f`$ and $`g`$ being dimensionless functions of the dimensionless radial variable (hereafter to be called $`u`$ or $`v`$) that vanish when the latter variable is greater than unity (corresponding to points outside the sphere of radius $`R`$), then
$$X=_0^1𝑑u_0^1𝑑v\left\{2uv\mathrm{ln}\left|\frac{u+v}{uv}\right|f(u)f(v)+\left[(u^2+v^2)\mathrm{ln}\left|\frac{u+v}{uv}\right|2uv\right]g^{}(u)g^{}(v)\right\}$$
(36)
and
$$YXZ=\frac{2\pi }{R}_0^1u^2𝑑u[f^2(u)+g^2(u)],$$
(37)
where the prime on the function $`g`$ denotes a derivative with respect to the argument (the dimensionless radius $`u|𝐱|/R`$ or $`v|𝐲|/R`$).
Now, if we take a density matrix of the form (19), let us try to maximize the product of the two nonzero eigenvalues of the density matrix,
$$yp(1p)=q(1q)s^2=q(1q)\left[1(1+4\beta ^2X^2)^{1/2}\mathrm{exp}\left(\frac{\alpha ^2X}{1+4\beta ^2X^2}\right)\right],$$
(38)
and hence maximize $`S_{\mathrm{vN}}(y)`$ given by Eq. (23), for fixed
$$x2\pi RE=2\pi Rq<0|U^1HU|0>=2\pi RZq(\alpha ^2X+4\beta ^2X^2).$$
(39)
Note that for fixed $`RZ=RY/X`$, the three quantities $`\alpha `$, $`\beta `$ (the coefficients of $`\chi `$ and of $`\chi ^2`$ in the hermitian operator $`h=\alpha \chi +\beta \chi ^2`$), and $`X<0|\chi ^2|0>`$ enter into this $`x`$ and $`y`$ only in the two nonnegative combinations $`a\alpha ^2X`$ and $`b4\beta ^2X^2`$, and $`x`$ depends only on $`q(a+b)`$. Then it is easy to see that for fixed $`q`$ and fixed $`a+b`$, $`y`$ decreases monotonically with $`b`$, so to maximize $`y`$ and $`S_{\mathrm{vN}}(y)`$ for fixed $`x`$, we should set $`\beta =0`$ in order to get $`b=0`$, $`h=\alpha \chi `$,
$$a\alpha ^2X=<0|h^2|0>,$$
(40)
$$x=2\pi RZqa,$$
(41)
and
$$y=q(1q)(1e^a).$$
(42)
Next, in our attempt to maximize $`y`$ as a function of $`x`$, we may continue to hold
$$\gamma 2\pi RZ$$
(43)
fixed and hence maximize $`y`$ for fixed $`zx/\gamma =aq`$. Then one can easily calculate that $`y=q(1q)(1e^{z/q})`$ has its maximum at fixed $`z`$ when
$$q=\frac{e^a1a}{2e^a2a}\frac{1}{2}a(1\frac{2}{3}a),$$
(44)
giving
$$x=\gamma aq=\frac{\gamma a(e^a1a)}{2e^a2a}\frac{1}{2}\gamma a^2(1\frac{2}{3}a)$$
(45)
and
$$y=[1(1+a)e^a]\left(\frac{e^a1}{2e^a2a}\right)^2\frac{1}{2}a^2(1\frac{5}{3}a)\frac{x}{\gamma },$$
(46)
where all the approximate equalities apply for $`a1`$. For fixed $`\gamma `$, $`a`$ is given implicitly as a function of $`x`$ by Eq. (45), and then inserting Eq. (46) for $`y`$ into Eq. (23) for $`S_{\mathrm{vN}}(y)`$ gives the entropy (so far maximized over $`\beta `$ and $`q`$) of the density matrix (19) explicitly as a function of $`a=<0|h^2|0>`$ and hence implicitly as a function of $`x`$. In fact, for $`x\gamma `$, we get the asymptotic relation
$$S_{\mathrm{vN}}(y)\frac{x}{\gamma }\left(\mathrm{ln}\frac{\gamma }{x}+1\right),$$
(47)
which of course exceeds $`x`$ for sufficiently small $`x`$.
As the final step in the maximization of the von Neumann entropy $`S_{\mathrm{vN}}(y)`$ of a density matrix of the particular form (19) for fixed $`x2\pi RE=\gamma z=\gamma aq`$, we note that maximizing $`y`$ (and hence $`S_{\mathrm{vN}}(y)`$) for fixed $`x`$ is equivalent to minimizing $`x`$ for fixed $`y`$. Therefore, we need to minimize $`\gamma 2\pi RZ2\pi RY/X`$. By looking at Eqs. (30) and (31) for $`X`$ and $`Y`$, we see that the ratio $`ZY/X`$ is invariant under any constant rescaling of the functions $`F(𝐱)`$ and $`G(𝐱)`$ that appear as smearing functions for $`\varphi (t=0,𝐱)`$ and $`\dot{\varphi }(t=0,𝐱)`$ in the defining Eq. (25) for the linear hermitian field operators $`\chi `$ and $`h=\alpha \chi `$ \[now that we have set $`\beta =0`$ to drop the nonlinear term for $`h`$ in Eq. (26)\]. The quantity $`\gamma `$ is also invariant under a rescaling of the radius $`R`$ if $`F(𝐱)`$ and $`F(𝐱)`$ depend only on $`𝐱/R`$ and on some overall constant factor that can depend on $`R`$.
Minimizing $`\gamma =2\pi RY/X`$ is thus equivalent to maximizing $`X`$ for fixed $`Y`$, which by Eq. (31) is half the integral of the sum of the squares of $`F(𝐱)`$ and of the gradient of $`G(𝐱)`$. Because the double integral (30) for $`X`$ is also quadratic in $`F(𝐱)`$ and in $`G(𝐱)`$ but has a positive-definite nonlocal kernel, maximizing $`X`$ at fixed $`Y`$ is best done with fairly smooth functions $`F(𝐱)`$ and $`G(𝐱)`$. In particular, if $`F(𝐱)`$ is expanded in spherical harmonics, one can readily see that the maximum is obtained by keeping only the spherically symmetric ($`\mathrm{}=0`$) terms. Also, since it is the dot product of $`G(𝐱)`$ and $`G(𝐲)`$ that enters into Eq. (30), which generically dilutes its contribution relative to that of $`F(𝐱)F(𝐲)`$ for the same values of the integrals of $`|F(𝐱)|^2`$ and of $`|G(𝐱)|^2`$ in Eq. (31), one can readily see that the maximum for $`X`$ at fixed $`Y`$ is obtained by setting $`G(𝐱)=0`$, as well as choosing a spherically symmetric $`F(𝐱)=f(|𝐱|/R)/R^2`$ as given by Eq. (34). Then one gets that
$$\gamma =\frac{4\pi ^2_0^1u^2𝑑uf^2(u)}{_0^1𝑑u_0^1𝑑v\mathrm{\hspace{0.25em}2}uv\mathrm{ln}\left|\frac{u+v}{uv}\right|f(u)f(v)}.$$
(48)
One then sees that the minimum value for $`\gamma `$ is
$$\gamma =\frac{4\pi ^2}{\lambda },$$
(49)
where $`\lambda `$ is the largest eigenvalue of the weakly singular linear Fredholm integral equation of the third kind,
$$_0^1𝑑v\mathrm{\hspace{0.25em}2}\mathrm{ln}\left|\frac{u+v}{uv}\right|w(v)=\lambda w(u)$$
(50)
for $`0u1`$, with $`w(u)=uf(u)`$ being the eigenfunction.
For $`F(𝐱)=f(|𝐱|/R)/R^2`$ to be a smooth function of $`𝐱`$, $`f`$ should be a smooth even function of its argument $`u|𝐱|/R`$. This means that $`w(u)`$ should be a smooth odd function of $`u`$, so we can expand it as an infinite sum of odd Legendre polynomials $`P_{2m1}(u)`$,
$$w(u)=\underset{m=1}{\overset{\mathrm{}}{}}c_mP_{2m1}(u).$$
(51)
This expansion converts the integral eigenvalue Eq. (50) into the matrix eigenvalue equation
$$\underset{n=1}{\overset{\mathrm{}}{}}A_{mn}c_n=\lambda \underset{n=1}{\overset{\mathrm{}}{}}B_{mn}c_n,$$
(52)
where the matrix components are
$`A_{mn}`$ $`=`$ $`{\displaystyle _0^1}𝑑u{\displaystyle _0^1}𝑑v\mathrm{\hspace{0.25em}2}\mathrm{ln}\left|{\displaystyle \frac{u+v}{uv}}\right|P_{2m1}(u)P_{2n1}(v)`$ (53)
$`=`$ $`{\displaystyle \frac{2}{[14(mn)^2](m+n)(m+n1)}}`$
and
$$B_{mn}=_0^1𝑑u_0^1𝑑vP_{2m1}(u)P_{2n1}(v)=\frac{\delta _{mn}}{4m1}.$$
(54)
\[Actually, I cheated slightly in obtaining the explicit expression above for the matrix components $`A_{mn}`$. I calculated $`A_{11}=1`$ by hand, but when I tried to calculate the general $`A_{mn}`$, I got finite sums that I did not readily see how to simplify. Therefore, I resorted to Maple. I did not quickly see how to get it to give me a simple general expression for $`A_{mn}`$ either, but in one afternoon I was able to get it to give me all the values for $`m<10`$, $`n<10`$ (45 different terms, since $`A_{mn}=A_{nm}`$). The form of these terms was sufficiently simple that part way through their rather slow evaluation I was able to deduce the simple expression given in Eq. (53), which indeed fit all 45 terms. So although I have not bothered to find a rigorous proof that Eq. (53) is correct for all $`m`$ and $`n`$ not both smaller than 10, the fact that it is a very simple formula that works for all 45 smaller values strongly suggests that it is exact for all values of $`m`$ and $`n`$. I could say that the proof is left as an exercise for the reader.\]
Maple readily solved the matrix eigenvalue Eq. (52) for various truncations of the infinite matrices $`A_{mn}`$ and $`B_{mn}`$. For example, 40-digit precision for $`70\times 70`$, $`80\times 80`$, $`90\times 90`$, $`100\times 100`$, $`110\times 110`$, and $`200\times 200`$ truncations all gave the largest eigenvalue agreeing to 13 digits:
$$\lambda 3.132010216749.$$
(55)
A 20-digit calculation of the $`60\times 60`$ case gave the last digit 8 instead of 9 but was used to get the following approximate expansion of the eigenfunction corresponding to the largest eigenvalue:
$`w(u)`$ $`+`$ $`P_1(u)0.3968319408P_3(u)+0.0102661635P_5(u)`$ (56)
$``$ $`0.0070631137P_7(u)0.0032552106P_9(u)0.0018849293P_{11}(u)`$
$``$ $`0.0011814233P_{13}(u)0.0007878354P_{15}(u)0.0005510316P_{17}(u)`$
$``$ $`0.0004002345P_{19}(u)0.0002997366P_{21}(u)0.0002302203P_{23}(u)`$
$``$ $`0.0001806208P_{25}(u)0.0001442924P_{27}(u)0.0001170806P_{19}(u)`$
$`+`$ $`\mathrm{terms}\mathrm{with}\mathrm{coefficients}\mathrm{less}\mathrm{than}0.0001.`$
One can notice that only $`P_1(u)=u`$ and $`P_3(u)=1.5u+2.5u^3`$ give large contributions to the eigenfunction, so one can get a fairly accurate estimate of the largest eigenfunction by taking even just the $`2\times 2`$ truncation of the matrices, which gives the eigenvalue
$$\lambda _2=\frac{5(15+\sqrt{57})}{36}3.131921449343,$$
(57)
which is smaller than the actual largest eigenvalue for the infinite matrices by less than one part in 35 283. An even simpler, but rather ad hoc, approximation is to change the coefficient of the $`P_3(x)`$ term above to 0.4 and drop all the higher terms. Dividing this trial function for $`w(u)`$ by $`u`$ (and multiplying by 8 to avoid fractions in the answer) gives $`f(u)=85u^2`$, which may be inserted into Eqs. (48) and (49) to give another estimate for $`\lambda `$,
$$\lambda _{\mathrm{est}}=\frac{1757}{561}3.131907308378,$$
(58)
which has almost 16% more error than $`\lambda _2`$, though this is still only a tiny error, being smaller than the actual largest eigenvalue for the infinite matrices by less than one part in 30 434. Even the very crude constant trial function for $`f(u)`$ gives an eigenvalue estimate, $`\lambda _{\mathrm{crude}}=3`$, that is smaller than the actual largest eigenvalue for the infinite matrices by only about 4.215%, or less than one part in 23.
Using the approximation Maple gave for $`\lambda `$, the largest eigenvalue of the infinite matrices, Eq. (49) then gives
$$\gamma =\frac{4\pi ^2}{\lambda }12.604817632215,$$
(59)
which can be used in Eq. (47) to get the asymptotic behavior of $`S_{\mathrm{vN}}(x)`$ at sufficiently small $`x`$. One can then see that this gives $`S_{\mathrm{vN}}(x)>x`$ for
$$x<\gamma e^{1\gamma }0.000115,$$
(60)
or alternatively for
$$S_{\mathrm{vN}}<\gamma e^{1\gamma }0.000115.$$
(61)
Thus if the dimensionless energy, $`x2\pi RE`$, and the von Neumann entropy, $`S=S_{\mathrm{vN}}tr\rho \mathrm{ln}\rho `$, are sufficiently small, then with the definitions used here for these quantities, they can violate Bekenstein’s conjectured entropy bound (1), $`Sx`$, though admittedly the range of $`x`$ and $`S`$ for which this happens is very narrow.
We may now use the value of $`\gamma `$, given by Eq. (59), in Eqs. (23), (45), and (46) to get $`BS_{\mathrm{vN}}/x`$ as a precise implicit function purely of $`x`$, or, alternatively, to get both $`x`$ and $`B`$ as explicit functions of $`a=<0|h^2|0>`$. Of course, this is merely for one simple example of a one-parameter family of mixed states given by Eq. (19), with $`q`$ given by Eq. (44) and $`|\psi >`$ given by Eq. (16) with $`\beta =0`$ so $`h=\alpha \chi `$ and with $`\chi `$ given by Eq. (25) with $`G(𝐱)=0`$ and Eq. (34) giving $`F(𝐱)=f(|𝐱|/R)/R^2`$ with $`w(u)=uf(u)`$ being an eigenvector corresponding to the largest eigenvalue, $`\lambda `$, of the homogeneous linear Fredholm integral equation (50). Therefore, it is not likely to give the maximum possible $`B`$ as a function of $`x`$, which was called $`B_\mathrm{N}(x)`$ in Eq. (9). However, this $`B(x)`$ does give at least a lower bound on $`B_\mathrm{N}(x)`$ for a single massless scalar field.
## 5 Conjectures for Entropy Bounds of Vacuum-Outside-R States
I would conjecture that asymptotically at small $`x2\pi RE`$, the density matrix (19), with all the entropy-maximization procedures given above for a density matrix of this form, gives $`B(x)`$ that does asymptotically approach the unknown global maximum function $`B_\mathrm{N}(x)`$ for a single massless scalar field. Therefore, if we divide $`x`$ into the asymptotic form of $`S_{\mathrm{vN}}`$ for small $`x`$ that is given by Eq. (47), I would conjecture that this gives the asymptotic form of the true upper bound, $`B_\mathrm{N}(x)`$, for very small $`x`$,
$$B_\mathrm{N}(x)\frac{1}{\gamma }\left(\mathrm{ln}\frac{\gamma }{x}+1\right),$$
(62)
with Eq. (59) giving $`\gamma 12.604817632215`$. One might expect a similar formula for other free massless fields, though perhaps each with a different value of $`\gamma `$.
When $`x`$ is not small, it is certainly not the case that the density matrix of fixed $`x`$ needed to maximize the entropy is approximately of the simple rank-two form given by Eq. (19). One would surely need a more general vacuum-outside-$`R`$ state, with a density matrix obeying Eq. (2) (giving vacuum expectation values for all operators not in causal contact with the ball $`rR`$ at $`t=0`$), such as that given by Eq. (15), most likely with an infinite sum of terms and an infinite rank. I do not know how to proceed toward finding such a density matrix obeying Eq. (2) that would maximize $`B(x)S/x`$ at finite $`x`$ that is neither asymptotically small or large. However, one might try using in Eq. (15) $`h_i`$’s that have the form given in Eq. (26), with the $`F(𝐱)`$’s and $`G(𝐱)`$’s of Eq. (25) being suitable eigenfunctions of the three-dimensional version of the integral equation (50). Even this wide class of examples may not be sufficient, since one could imagine instead constructing the Hermitian operators $`h_i`$ from smeared functions of the field $`\varphi (t=0,𝐱)`$ and of its time-derivative (or conjugate momentum) $`\dot{\varphi }(t=0,𝐱)`$ that are not merely linear as is the $`\chi `$ given by Eq. (25). Going to nonlinear Hermitian operators (other than the relatively simple $`\chi ^2`$ considered above) leads to such a wealth of possibilities that I do not presently know how to proceed to obtain a true maximum for $`B(x)`$ at fixed finite $`x`$, the postulated function $`B_\mathrm{N}(x)`$.
In the more usual case of boundary conditions on the field, quantum states obeying these boundary conditions may be coherently superposed (i.e., the corresponding wavefunctions added, not just the density matrices) to get other states that also obey the boundary conditions. Then one can look for superpositions that diagonalize the Hamiltonian (i.e., energy eigenstates). From these, one can form a Gibbs ensemble to maximize the von Neumann entropy at a fixed expectation value of the energy.
However, for the vacuum-outside-$`R`$ states considered here, coherent superpositions of pure vacuum-outside-$`R`$ states are generically not vacuum-outside-$`R`$ states. (Of course, positive-weight combinations of vacuum-outside-$`R`$ density matrices are still vacuum-outside-$`R`$ density matrices when normalized, since the vacuum-outside-$`R`$ condition, that all expectation values outside the ball $`rR`$ at $`t=0`$ are the same as the vacuum, is homogeneous and linear in the density matrix, though not in the wavefunction.) Therefore, the procedure for diagonalizing the Hamiltonian for such states fails.
Indeed, one can see that none of the vacuum-outside-$`R`$ states, except for the vacuum itself, can be an energy eigenstate. This is because any state which is non-vacuum in a finite region at some time will inevitably have that region spread with time. For fields with linear field equations, the perturbations of the field itself will spread. But even for fields with self-coupling which allow classical field solitons that do not spread with time, any quantum state of the field which is non-vacuum at some time will inevitably have that region spread with time as a result of the quantum uncertainty principle. For example, suppose that there is some definition of the location and momentum of the soliton, such that the velocity of the location is proportional to the momentum. Then the position-momentum uncertainty principle will prevent one from having that the location remain, with certainty, within any finite region for an infinite amount of time; the quantum uncertainty of the position, if initially confined to a finite region, will inevitably spread to extend all over space. Therefore, the confined configuration cannot be stationary and hence cannot be an energy eigenstate.
One can test my conjecture that Eq. (62) is the correct asymptotic form of the true upper bound on the entropy per $`x2\pi RE`$ by examining some other simple density matrices of the form given by Eq. (15) that allow explicit evaluation of the von Neumann entropy. For example, the rank-three density matrix
$$\rho =(1q)|0><0|+(q/2)e^{i\alpha \chi }|0><0|e^{i\alpha \chi }+(q/2)e^{i\alpha \chi }|0><0|e^{i\alpha \chi }$$
(63)
has $`zx/\gamma =aq`$, just like rank-two density matrix (19) when $`\beta =0`$, and it has the three nonzero eigenvalues
$`p_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}[1{\displaystyle \frac{1}{2}}q(1e^{2a})]+{\displaystyle \frac{1}{2}}\sqrt{[1{\displaystyle \frac{1}{2}}q(1e^{2a})]^22q(1q)(1e^a)^2}`$ (64)
$``$ $`1q(1e^a)+{\displaystyle \frac{1}{2}}q^2(1e^a)^21qa+{\displaystyle \frac{1}{2}}(q+q^2)a^2,`$
$`p_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}q(1e^{2a})`$ (65)
$``$ $`qa(1a),`$
and
$`p_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}[1{\displaystyle \frac{1}{2}}q(1e^{2a})]{\displaystyle \frac{1}{2}}\sqrt{[1{\displaystyle \frac{1}{2}}q(1e^{2a})]^22q(1q)(1e^a)^2}`$ (66)
$``$ $`{\displaystyle \frac{1}{2}}q(1q)(1e^a)^2{\displaystyle \frac{1}{2}}q(1q)a^2,`$
where the approximate equations apply for very small $`a\alpha ^2X`$. If one chooses $`q`$ to maximize the von Neumann entropy of this mixed state,
$$S_{\mathrm{vN}}tr(\rho \mathrm{ln}\rho )=p_1\mathrm{ln}p_1p_2\mathrm{ln}p_2p_3\mathrm{ln}p_3,$$
(67)
for very small $`a`$ (equivalently, very small $`z`$), one gets that $`q(1/3)(1+4a/3)`$, so $`z=aq(a/3)+4(a/3)^2`$, which may be inverted to give $`a3z(14z)`$ and $`q(1/3)(1+4z)`$. This then gives
$$S_{\mathrm{vN}}z[1+z(1z)\mathrm{ln}z]\frac{x}{\gamma }\left(\mathrm{ln}\frac{\gamma }{x}+1\right),$$
(68)
which has the same asymptotic form for small $`x`$ as Eq. (47) for the rank-two density matrix (19).
Another example would be to consider the rank-five density matrix
$`\rho =(1q)|0><0|`$ $`+`$ $`(q/4)e^{i\alpha \chi +i\beta \chi ^2}|0><0|e^{i\alpha \chi i\beta \chi ^2}`$ (69)
$`+`$ $`(q/4)e^{i\alpha \chi i\beta \chi ^2}|0><0|e^{i\alpha \chi +i\beta \chi ^2}`$
$`+`$ $`(q/4)e^{i\alpha \chi +i\beta \chi ^2}|0><0|e^{i\alpha \chi i\beta \chi ^2}`$
$`+`$ $`(q/4)e^{i\alpha \chi i\beta \chi ^2}|0><0|e^{i\alpha \chi +i\beta \chi ^2},`$
which gives $`zx/\gamma =q(a+b)`$, where, as above, $`a\alpha ^2X`$ and $`b4\beta ^2X^2`$. Even though the eigenvalue equation is now a fifth-order polynomial equation, it appears that one may be able to use the symmetries of the problem to find the eigenvectors and eigenvalues explicitly, as functions of $`q`$, $`a`$, and $`b`$, without requiring any roots higher than square roots. However, this seems to be messier than is worth doing here, so it shall be left as another exercise for the reader.
Nevertheless, one can show that when $`a1`$ and $`b1`$, there is one eigenvalue near unity, one near $`qa`$, one near $`qb/2`$, and the remaining two are smaller by factors of the order of $`z=q(a+b)`$. Therefore, in this limit only the three largest eigenvalues contribute significantly to the von Neumann entropy, giving
$$S_{\mathrm{vN}}(1qaqb/2)\mathrm{ln}(1qaqb/2)qa\mathrm{ln}(qa)(qb/2)\mathrm{ln}(qb/2).$$
(70)
When this is maximized at fixed $`z=q(a+b)`$, one finds that the first three eigenvalues need to be approximately in a geometric series (as, e.g., are all the eigenvalues of the thermal density matrix for an harmonic oscillator), giving $`qb/2(qa)^2(1+qa)`$. Solving for $`qa`$ and $`qb`$ in terms of $`z`$ and inserting this back into Eq. (70) gives the maximum von Neumann entropy
$$S_{\mathrm{vN}}z[1+\frac{1}{2}z\mathrm{ln}z]\frac{x}{\gamma }\left(\mathrm{ln}\frac{\gamma }{x}+1\right)$$
(71)
for the rank-five density matrix (69) at fixed tiny $`zx/\gamma `$. This entropy is just slightly larger, by an amount roughly $`z^2[\mathrm{ln}(1/z)(1/2)]`$, than the corresponding maximum entropy (68) for the rank-three density matrix (63) at tiny $`z=x/\gamma `$, but it has the same asymptotic limit (given after the $``$ sign).
Therefore, although of course these three simple examples of finite-rank density matrices do not begin to exhaust the infinite set of possibilities, they give some support to the conjecture given above that Eq. (62) is the correct asymptotic form of the upper bound for $`S/(2\pi RE)`$ when the denominator of this expression, $`x2\pi RE`$, is much smaller than unity.
In the opposite limit, when $`x2\pi RE`$ is much greater than unity, we would expect that, at least for a scale-invariant field so that the energy $`E`$ is large with respect to all relevant parameters with the same dimension ($`1/R`$ being the only relevant one if the field does not have a rest mass or other parameter setting a higher energy scale), the maximum entropy is given by Eq. (11) for high-temperature thermal radiation, giving
$$B_\mathrm{N}(x)\frac{\beta }{x^{1/4}},$$
(72)
with now
$$\beta =\left[\frac{2^6}{3^65}(n_b+\frac{7}{8}n_f)\right]^{\frac{1}{4}},$$
(73)
no longer the $`\beta `$ of Eq. (26) that we have subsequently set to zero to maximize $`B(x)`$ for the density matrix (19). It is tempting to combine this asymptotic formula for $`x1`$ with the asymptotic formula (62) conjectured above for $`x1`$ to conjecture that a reasonably good approximate formula for $`B_\mathrm{N}(x)`$, as a function of any $`x2\pi RE`$, for a quantum field theory with a given set of massless fields, is
$$B_\mathrm{N}(x)\frac{4}{\gamma }\mathrm{ln}\left(1+\frac{\beta \gamma }{4x^{1/4}}\right),$$
(74)
where the constants $`\beta `$ and $`\gamma `$ would depend upon the massless fields in the theory. For the single massless real scalar field that has been considered here, $`n_b=1`$ and $`n_f=0`$, so Eq. (26) would give
$$\beta =\left(\frac{2^6}{3^65}\right)^{\frac{1}{4}}0.364016115028,$$
(75)
and Eq. (59) has already given $`\gamma 12.604817632215`$ for the single real massless scalar field.
Another way to state this conjecture is to write
$$B_\mathrm{N}(x)=\frac{4}{\gamma }\mathrm{ln}\left(1+\frac{\beta \gamma }{4x^{1/4}}\right)C(x),$$
(76)
where $`C(x)`$ is a correction factor yet to be found, and then conjecture that $`C(x)`$ tends asymptotically to unity for both very small and very large $`x`$, and perhaps further to conjecture that $`C(x)`$ is always relatively close to unity (e.g., say within a factor of two). This conjecture would then imply a conjectured entropy bound,
$$S_{\mathrm{vN}}\frac{8}{\gamma }\mathrm{ln}\left(1+\frac{\beta \gamma }{4(2\pi RE)^{1/4}}\right)2\pi RE.$$
(77)
An improved formula might be to write
$$B_\mathrm{N}(x)=\frac{4}{\gamma }\mathrm{ln}\left(\frac{1+Ax^{1/4}+Bx^{1/2}}{1+Cx^{1/4}}\right)\stackrel{~}{C}(x)$$
(78)
with $`B/C=(e\gamma )^{1/4}`$ to fit the final 1 in the asymptotic formula (62) for $`x1`$, $`AC=\beta \gamma /4`$ to fit the asymptotic formula (72) for $`x1`$, and $`2BA^2+C^2=\gamma \delta /2`$ to fit the following two-term improvement to Eq. (72):
$$B_\mathrm{N}(x)\frac{\beta }{x^{1/4}}+\frac{\delta }{x^{1/2}}.$$
(79)
I have not tried to work out what $`\delta `$ is. It would be straightforward to calculate, if it were the same as for the thermal state of a massless scalar field inside a sphere with Dirichlet boundary conditions on the field at the boundary $`r=R`$, but it is not obvious to me whether or not it is the same.
Since Eq. (78) with the correction factor $`\stackrel{~}{C}(x)`$ omitted should give a better asymptotic fit to $`B_\mathrm{N}(x)`$ than Eq. (74), I would expect that $`\stackrel{~}{C}(x)`$ would generally be closer to unity than the corresponding correction factor $`C(x)`$ of Eq. (76) \[not to be confused with the coefficient $`C`$ in Eq. (78)\]. But whether this is true over the entire infinite range of $`x`$ remains to be seen.
For free massive quantum fields, for fixed entropy one would expect that the energy would have to be higher, so an upper bound for a set of free massless quantum fields should also give an upper bound for a corresponding set of free massive quantum fields. Therefore, I would conjecture that for any given free quantum field theory, one can find a $`\beta `$ and $`\gamma `$ (presumably with $`\beta `$ obeying Eq. (75), and $`\gamma `$ some combination of eigenvalues of the appropriate integral equations) such that the inequality (77) holds with that value of $`\beta `$ and $`\gamma `$. The conjecture might even be true for any reasonable interacting quantum field theory that is causal in the sense defined above, though then one might need a different value of $`\beta `$.
## 6 Possibilities for Trying to Retain Bekenstein’s Proposed Bound
Returning to a consideration of how Bekenstein’s proposed bound (1) fits with the results derived and conjectured here, one first notes that the results here violate (1) for the von Neumann entropy $`S_{\mathrm{vN}}`$ of a vacuum-outside-$`R`$ mixed state with sufficiently small energy expectation value $`E`$, e.g., for $`x2\pi RE<0.000115`$ in the example above. However, even if one accepts the use of vacuum-outside-$`R`$ states for defining a finite size $`R`$, one might still object that the Bekenstein bound is not intended to be applied to the definition of $`E`$ and/or $`S`$ being used here.
For example, Schiffer and Bekenstein refer to “quantum states accessible to the field system with energy up to and including $`E`$.” It could be objected that since the vacuum-outside-$`R`$ states considered above are not energy eigenstates, they are actually composed of states with energy both lower and higher than the energy expectation value that I have used as the definition of $`E`$. If one takes $`E`$ to be the energy of one of the energy eigenstates that is sufficiently higher than the expectation value, then the Bekenstein bound (1) may be obeyed even in the examples I have given above that violate the bound when $`E`$ is taken to be the energy expectation value.
But if one takes this approach, it is hard to see how to give any content to the proposed bound for the vacuum-outside-$`R`$ states. Presumably not only is it the case that any vacuum-outside-$`R`$ state is not an energy eigenstate (since it is not stationary), but also it is surely the case that if any vacuum-outside-$`R`$ state is decomposed into energy eigenstates, it will include energy eigenvalues of arbitrarily large value. However, using an arbitrarily large value of $`E`$ in the bound (1) makes it trivial, entropy less than or equal to infinity. Therefore, for the bound to have any content, we need to have a definition of $`E`$ that gives finite values. The definition given by Eq. (3) above, $`Etr(H\rho )`$, is surely the simplest, though others could be proposed.
For example, one could propose instead that for a vacuum-outside-$`R`$ density matrix of the form (15), $`E`$ could be defined as the maximum value of the expectation value of the Hamiltonian $`H`$ in any of the normalized states $`e^{ih_i}|0><0|e^{ih_i}`$ whose sum, weighted by the $`q_i`$’s, forms $`\rho `$. However, using this definition would not avoid violations of Bekenstein’s bound (1). For example, one could use the density matrix (63) with $`q=1`$ so that the first term vanishes, and then the remaining two terms are of the form (15) with $`h_1=\alpha \chi `$ and $`h_2=\alpha \chi `$, and with $`q_1=q_2=1/2`$. Each of the two nonzero terms of the density matrix then gives the same energy expectation value, and one can calculate that for small $`z=x/\gamma =a\alpha ^2X`$ one gets
$$S_{\mathrm{vN}}z[1\frac{1}{2}z(1z)\mathrm{ln}z]\frac{x}{\gamma }\left(\mathrm{ln}\frac{\gamma }{x}+1\right),$$
(80)
again violating Bekenstein’s bound (1) for $`x<0.000115`$.
Yet another proposal would be that $`E`$ be defined as $`E_{\mathrm{max}}`$, the maximum expectation value of the Hamiltonian in all of the eigenstates of the density matrix. In the example just discussed,
$$\rho =p_1|1><1|+p_2|2><2|$$
(81)
with $`p_1=(1+e^{2a})/2`$, $`p_2=(1e^{2a})/2`$, and orthonormal eigenstates
$$|1>=\frac{\mathrm{cos}(\alpha \chi )|0>}{\sqrt{p_1}}$$
(82)
and
$$|2>=\frac{\mathrm{sin}(\alpha \chi )|0>}{\sqrt{p_2}}.$$
(83)
Then by using Eq. (33), one can calculate not only that
$`<0|e^{ih_1}He^{ih_1}|0><0|e^{i\alpha \chi }He^{i\alpha \chi }|0>`$ $`=`$ $`<0|e^{ih_2}He^{ih_2}|0><0|e^{i\alpha \chi }He^{i\alpha \chi }|0>`$ (84)
$`=`$ $`\alpha ^2Y\alpha ^2XZaZ`$
as given by Eq. (29), but also
$`<0|e^{ih_1}He^{ih_2}|0><0|e^{i\alpha \chi }He^{i\alpha \chi }|0>`$ $`=`$ $`<0|e^{ih_2}He^{ih_1}|0><0|e^{i\alpha \chi }He^{i\alpha \chi }|0>`$ (85)
$`=`$ $`aZe^a.`$
From these results and from the form of the orthonormal eigenstates given by Eqs. (82) and (83), one readily obtains
$$H_{11}<1|H|1>=\frac{aZ(1e^a)}{1+e^{2a}}\frac{1}{2}a^2Z$$
(86)
and
$$H_{22}<2|H|2>=\frac{aZ}{1e^a}Z$$
(87)
Then if one takes $`E_{\mathrm{max}}`$, the larger of $`H_{11}`$ and $`H_{22}`$, namely $`H_{22}`$, as the definition of $`E`$, one sees that it has the positive lower limit $`Z=\gamma /(2\pi R)`$ as one takes $`a`$ (and hence $`S_{\mathrm{vN}}`$) to zero, so with this definition one does not get a violation of Bekenstein’s bound (1) in this example. In particular, for this example
$`B_{E_{\mathrm{max}}}{\displaystyle \frac{S_{\mathrm{vN}}}{2\pi RE_{\mathrm{max}}}}`$ $`=`$ $`{\displaystyle \frac{1e^a}{\gamma a}}\{{\displaystyle \frac{1}{2}}(1+e^{2a})\mathrm{ln}[{\displaystyle \frac{1}{2}}(1+e^{2a})]`$ (88)
$``$ $`{\displaystyle \frac{1}{2}}(1e^{2a})\mathrm{ln}[{\displaystyle \frac{1}{2}}(1e^{2a})]\}<1.`$
Whether this definition of $`B_{E_{\mathrm{max}}}`$ always gives a result less than unity for all vacuum-outside-$`R`$ states, thus agreeing with Bekenstein’s bound, remains to be proven, but the extremely meagre evidence that I have does not seem to contradict this conjecture.
Another objection that might be made against the violations of Bekenstein’s bound (1) using vacuum-outside-$`R`$ states to define $`R`$, using the expectation value of the Hamiltonian given by Eq. (3) as the definition of the energy $`E`$, and using the von Neumann entropy given by Eq. (4) as the definition of the entropy $`S`$, is to demand that the entropy instead be given by a microcanonical ensemble rather than by Eq. (4) applied to any mixed state. In other words, instead of allowing a generic vacuum-outside-$`R`$ mixed state or density matrix $`\rho `$ obeying Eq. (2) for all operators $`O`$ completely confined to the region $`D`$, $`r>R+|t|`$, one might propose that Bekenstein’s conjectured bound should only be applied to density matrices made up of equal mixtures of $`n`$ orthogonal vacuum-outside-$`R`$ pure states. (These are rank-$`n`$ density matrices with precisely $`n`$ nonzero eigenvalues, all equal to $`1/n`$, and with the corresponding set of $`n`$ orthonormal eigenvectors all being vacuum-outside-$`R`$ pure states.) The entropy of such a density matrix whose nontrivial part is proportional to the identity matrix in the $`n`$ nontrivial dimensions is then $`S=\mathrm{ln}n`$.
Again, I do not have evidence that Bekenstein’s conjectured bound (1) is violated for such restricted vacuum-outside-$`R`$ density matrices. However, it is a rather severe limitation to restrict the discussion to such a small subset of vacuum-outside-$`R`$ density matrices, a subset of measure zero in the space of all such density matrices. Furthermore, it appears rather difficult to find many explicit examples of precisely orthogonal vacuum-outside-$`R`$ pure states.
For example, for each fixed choice of the two functions $`F(𝐱)`$ and $`G(𝐱)`$ in Eq. (25), the constants $`\alpha `$ and $`\beta `$ in Eq. (26) give a two-parameter family of vacuum-outside-$`R`$ pure states of the form $`e^{i\alpha \chi +i\beta \chi ^2}|0>`$, and then Eq. (27) and its trivial generalization gives the inner product between any two states among this two-parameter family. However, none of these inner products are zero for finite $`\alpha `$’s and $`\beta `$’s (and hence for finite expectation values of the energy), so none of these states are orthogonal for fixed $`F(𝐱)`$ and $`G(𝐱)`$.
Of course, if one combines various ones of these nonorthogonal pure state density matrices to get a mixed density matrix and then finds the eigenvectors of that density matrix, they will form an orthonormal set of density matrices, such as the set $`|1><1|`$ and $`|2><2|`$ of Eqs. (82) and (83). However, these density matrices are not by themselves vacuum-outside-$`R`$ states, but only when they are combined with the particular eigenvalues $`p_1`$ and $`p_2`$ given just before Eq. (82). Hence they cannot be used in a different linear combination (e.g., with $`p_1=p_2=1/2`$) to get a vacuum-outside-$`R`$ state that is an equal-weight combination of $`n`$ orthonormal pure-state vacuum-outside-$`R`$ states.
The only explicit vacuum-outside-$`R`$ state orthogonal to the vacuum itself (the trivial vacuum-outside-$`R`$ state) that I have found so far is the pure state
$$|\psi >=\mathrm{exp}\left(i\alpha e^{\beta \chi ^2}\right)|0>$$
(89)
with real $`\alpha `$ and $`\beta `$ (not the same $`\alpha `$ and $`\beta `$ used elsewhere in this paper) chosen so that
$$<0|\psi >=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(i\alpha )^n}{n}(1+2n\beta X)^{1/2}=0.$$
(90)
Using Maple, I found a numerical solution at
$$\alpha 4.727048274,\beta 1.536994796/X.$$
(91)
I have not worked out the expectation value of the energy, $`<\psi |H|\psi >`$, of this pure state, but I suspect that it is greater than $`\mathrm{ln}2/(\pi R)`$, so that the entropy of an equal mixture of this state and of the vacuum, $`\mathrm{ln}2`$, would be less than $`2\pi R`$ times the expectation value of the energy in this mixed state (half the expectation value of the energy of the pure state $`|\psi ><\psi |`$, since the vacuum half of the mixed state contributes zero to the expectation value of the energy). If so, then this example would not be a counterexample to Bekenstein’s conjectured bound (1) restricted to microcanonical ensembles that are equal mixtures of orthogonal pure vacuum-outside-$`R`$ states.
It would be interesting to find the lowest-energy vacuum-outside-$`R`$ state orthogonal to the vacuum itself and see whether its energy is indeed not more than $`\mathrm{ln}2/(\pi R)`$, but I do not see how to do this at present. More generally, one would like to find, for each positive integer $`n`$, the set of $`n`$ mutually orthogonal vacuum-outside-$`R`$ states (possibly, but not necessarily, including the vacuum itself) such that the sum of the $`n`$ energy expectation values, say $`E_s`$, is minimized. Then if one finds that $`E_s(n\mathrm{ln}n)/(2\pi R)`$ for each $`n`$, then Bekenstein’s conjectured bound (1) will be obeyed for these microcanonical ensembles of vacuum-outside-$`R`$ states.
One way to look for other pure vacuum-outside-$`R`$ states orthogonal to the vacuum would be to choose some non-hermitian operator, say $`\kappa `$, that is confined to the ball $`rR`$ at $`t=0`$, and consider the one-complex-parameter (two-real-parameter) set of states
$$|\psi (C;\kappa )>=e^{iC\kappa i\overline{C}\kappa }|0>$$
(92)
for the complex parameter $`C`$. For a generic such $`\kappa `$, $`<0|\psi (C;\kappa )>`$ would be a complex function of $`C`$ (not analytic, since both $`C`$ and its complex conjugate $`\overline{C}`$ appear in the definition of $`|\psi (C;\kappa )>`$), and a simple parameter-counting argument suggests that there should be discrete complex values of $`C`$ at which $`<0|\psi (C;\kappa )>=0`$, giving a pure vacuum-outside-$`R`$ states orthogonal to the vacuum, though of course for particular $`\kappa `$’s, the number of such discrete solutions for $`C`$ may be zero. One might extend this method to try to find $`n`$ mutually orthogonal vacuum-outside-$`R`$ states; this would require $`n`$ different operators $`\kappa _i`$ and $`n(n+1)/2`$ complex parameters. The obvious problem for carrying out this procedure explicitly is that for most sets of operators $`\kappa _i`$, the inner products (functions of the complex parameters) would be difficult to evaluate.
Perhaps a compromise to the stringent requirement of a microcanonical ensemble of $`n`$ equally-weighted orthogonal pure vacuum-outside-$`R`$ density matrices is simply to use the von Neumann entropy $`S_{\mathrm{vN}}tr\rho \mathrm{ln}\rho `$, which equals $`\mathrm{ln}n`$ for a microcanonical ensemble, but require that it be at least as large as $`\mathrm{ln}2`$, the minimum nontrivial value for a microcanonical ensemble. Then one might conjecture that the bound (1) is correct for vacuum-outside-$`R`$ states such that $`S=S_{\mathrm{vN}}\mathrm{ln}2`$.
Alternatively, one might replace Bekenstein’s conjectured bound (1) with the weaker conjectured bound
$$S2\pi ER+\mathrm{ln}2,$$
(93)
still using $`S=S_{\mathrm{vN}}tr\rho \mathrm{ln}\rho `$, $`Etr(H\rho )`$, and restricting to vacuum-outside-$`R`$ density matrices $`\rho `$ obeying $`tr(O\rho )=<0|O|0>`$ for all operators $`O`$ totally confined to the region $`D`$, $`r>R+|t|`$, that is not in causal contact with the ball $`B`$, $`rR`$ at $`t=0`$ (no causal curves connecting these two regions). (Equivalently, one may require that $`tr(O\rho )=<0|O|0>`$ for all operators $`O`$ that commute with all operators defined totally on the ball $`B`$.)
## 7 Other Ways to Define a Radius R
So far I have been considering only the new proposal to define $`R`$ by restricting to vacuum-outside-$`R`$ states. However, one might ask whether there are other ways to define a radius $`R`$ for a class of states for which one is seeking a bound on the entropy $`S`$ as a function of R and of the energy $`E`$.
One proposal that is very close to my proposal of vacuum-outside-$`R`$ states is a proposal for what might be called stressless-outside-$`R`$ states, states such that on the $`t=0`$ flat hyperplane of the Minkowski spacetime that I have always been assuming so far in this paper, the expectation value of the regularized stress-energy tensor operator, $`T_{\mu \nu }`$, is zero everywhere outside the radius $`R`$ (as it is everywhere for the vacuum state),
$$\tau _{\mu \nu }(𝐱)tr(T_{\mu \nu }(t=0,𝐱)\rho )=0$$
(94)
for all $`|𝐱|>R`$. Alternatively, one might restrict to what might be called energyless-outside-$`R`$ states,
$$\epsilon (𝐱)\tau _{00}tr(T_{00}(t=0,𝐱)\rho )=0$$
(95)
for all $`|𝐱|>R`$, for which the expectation value $`\epsilon (𝐱)`$ of the regularized energy density operator, $`T_{00}`$, at $`t=0`$ vanishes outside the radius $`|𝐱|=R`$.
Of course, all vacuum-outside-$`R`$ states are also stressless-outside-$`R`$ states, and all stressless-outside-$`R`$ states are also energyless-outside-$`R`$ states, but I do not know whether the converses of these statements are true. If they are not both true, there would exist energyless-outside-$`R`$ states, and possibly also stressless-outside-$`R`$ states, that are not also vacuum-outside-$`R`$ states. If there is indeed a broader class of states than vacuum-outside-$`R`$ states, whether stressless-outside-$`R`$ states and/or energyless-outside-$`R`$ states, then the corresponding entropy maximization function $`\sigma _{\mathrm{vN}}(R,E)`$ would be expected to be larger for the broader class of states.
One can try to define $`R`$ for even broader classes of states, not by requiring that some expectation values vanish for $`r>R`$ at $`t=0`$, but instead by using the spatial distribution of some quantity to define an effective radius $`R`$. One way that first comes to mind is to use some spatially-dependent real weight function $`W(𝐱)`$ coming from the quantum state to define $`R`$ as an rms value of $`r|𝐱|`$:
$$R_W^2=\frac{d^3xW(𝐱)r^2}{d^3xW(𝐱)}.$$
(96)
Of course, the weight function should be such that both the numerator and the denominator are finite and have the same sign (which without loss of generality will be assumed to be positive), at least for the class of states to be considered.
An obvious simple choice of the weight function is the energy density expectation value $`\epsilon (𝐱)`$. Then the denominator of Eq. (96) for $`R_W^2`$ is the total energy, which is positive for a nontrivial state. However, the numerator is not positive for all nontrivial states, as one can see from the following argument: Motivated by the state with locally negative energy density given by Kuo and Ford , consider the state
$$|\psi >=\alpha |0>+\beta |2>,$$
(97)
where $`|\alpha |^2+|\beta |^2=1`$ and $`|2>`$ is a two-quantum state of energy
$$E_2<2|H|2>=d^3x<2|T_{00}|2>$$
(98)
and with mode functions that are sufficiently localized that
$$R_2^2E_2d^3x<2|T_{00}|2>r^2$$
(99)
is finite. Because $`T_{00}`$ is the regularized (e.g., normal-ordered) energy density operator, $`<0|T_{00}|0>=0`$ and $`d^3x<2|T_{00}|0>=d^3x<0|T_{00}|2>=0`$. For a generic two-quantum state $`|2>`$,
$$Cd^3x<2|T_{00}|0>r^2$$
(100)
will be a nonzero complex number. Then Eq. (96) for $`W(𝐱)=\epsilon (𝐱)`$ gives
$$R_W^2=\frac{d^3x<\psi |T_{00}|\psi >r^2}{d^3x<\psi |T_{00}|\psi >r^2}=R_2^2+\mathrm{}\left(\frac{2C}{E_2}\frac{\alpha }{\beta }\right).$$
(101)
Since $`\alpha /\beta `$ can be an arbitrary complex number, $`R_W^2`$ can take any real value if $`C0`$, including zero and negative values. Even if one restricted to states for which $`R_W^2`$ is positive, this quantity can be made arbitrarily small, and then if such a state has finite energy and positive entropy (e.g., by being a mixture of the vacuum state $`|0><0|`$ and of $`|\psi ><\psi |`$), it can make $`BS/(2\pi RE)`$ arbitrarily large.
If one did want to use $`\epsilon (𝐱)`$ as the weight function $`W(𝐱)`$ in Eq. (96), one would have to restrict the states so that $`R_W^2`$ cannot be too small for states of finite energy and nonzero entropy. One way that might work would be to restrict the states to those in which $`\epsilon (𝐱)`$ is nonnegative everywhere, unlike the state $`|\psi ><\psi |`$ for sufficiently large $`C\alpha /\beta `$.
Another option that might work for all sufficiently localized states would be to choose a weight function $`W(𝐱)`$ that is nonnegative for all states. Examples of this would be $`\epsilon ^2`$, $`_{\mu =0}^4_{\nu =0}^4(\tau _{\mu \nu })^2`$, $`(\tau _\mu ^\mu )^2`$, $`(\tau _\nu ^\mu \tau _\mu ^\nu )^2`$, $`(\tau _\nu ^\mu \tau _\rho ^\nu \tau _\mu ^\rho )^2`$, $`(\tau _\nu ^\mu \tau _\rho ^\nu \tau _\sigma ^\rho \tau _\mu ^\sigma )^2`$, $`(tr(\varphi (t=0,𝐱)\rho ))^2`$, $`(tr(\dot{\varphi }(t=0,𝐱)\rho ))^2`$, $`(tr(:\varphi ^2(t=0,𝐱):\rho ))^2`$, $`(tr(:\dot{\varphi }^2(t=0,𝐱):\rho ))^2`$, $`(tr(:\varphi ^{;\mu }(t=0,𝐱)\varphi _{;\mu }(t=0,𝐱):\rho ))^2`$, etc., and positive powers of these positive functions of $`𝐱`$ at $`t=0`$.
The fourth quantity above is the square of $`\tau _\nu ^\mu \tau _\mu ^\nu `$, which itself usually seems to be positive everywhere, though it might be some restriction of states for this to be true everywhere for all states in the class.
If one used one of these positive quantities, or one of the usually positive quantities just for states in which it is everywhere positive, as a weight function $`W(𝐱)`$ for defining $`R`$ by Eq. (96), one would again presumably get some upper bound on the entropy $`S`$ as a function of $`R`$ and of $`E`$ (depending on how one defined $`S`$ and $`E`$ and what further restrictions one puts on the states). However, I have not investigated what these relations might be.
## 8 Difficulties with Entropy Bounds in Quantum and Semiclassical Gravity
So far I have been restricting attention to nongravitational quantum field theories in flat Minkowski spacetime. However, since the original motivation for Bekenstein’s conjectured entropy bound came from quantum considerations of gravitating black holes, it is interesting to consider whether a similar entropy bound can be applied in quantum gravity.
Here I must admit that I see serious problems in attempting to apply the bound to quantum gravity. Assuming that the quantum part of quantum gravity is sufficiently similar to the ordinary quantum theory of nongravitational systems, the entropy $`S`$ might still be a well-defined quantity, at least for a complete system, such as the entire universe (though if the ultimate quantum gravity theory specifies a unique quantum state, there may be no option as to what the entropy is). The energy $`E`$ is more problematic, at least if the universe is not asymptotically flat, though if one can restrict to quantum states in which the universe is asymptotically flat, then $`E`$ also might have a good definition in quantum gravity. However, what I don’t see how to give a good precise definition for is the size $`R`$.
The main problem is that states of quantum gravity should be coordinate invariant, so it is hard to see how to say that some state is confined to a radius $`R`$ or has this size. How would one define the center with respect to which the state is within a distance $`R`$? Furthermore, if the state is an asymptotically flat one with energy $`E`$, the gravitational field of this energy should extend all the way out to spatial infinity, so in that sense it seems that the state cannot be confined to be within radius $`R`$ or have vacuum properties outside that radius.
It is not that I have a rigorous proof that an entropy bound such as Eq. (1) cannot be applied in quantum gravity, but I just don’t see how it can be applied.
The situation seems somewhat more hopeful in semiclassical gravity, in which one has quantum field theory for nongravitational fields on a classical curved spacetime (perhaps whose Einstein tensor is proportional to the regularized expectation value of the stress-energy tensor of the nongravitational quantum fields). In this case one can imagine defining the equivalent of vacuum-outside-$`R`$ states of energy $`E`$ in the following way:
Take an asymptotically flat spherically symmetric spacetime which has a totally geodesic Cauchy hypersurface (with zero extrinsic curvature, say at $`t=0`$), about which it has time-reflection symmetry ($`tt`$). Use a Schwarzschildean radial coordinate $`r`$, the circumference/$`(2\pi )`$ of each symmetrical sphere. Outside the $`r=R`$ sphere on this $`t=0`$ hypersurface, assume that the spatial metric and the expectation value of all operators confined to this region are the same as that of the static spherically symmetric asymptotically Schwarzschild semiclassical metric with ADM mass $`E`$ and quantum state that is the semiclassical version of the zero-temperature Boulware state for this metric. The entropy $`S`$ can then be the von Neumann entropy $`S_{\mathrm{vN}}tr\rho \mathrm{ln}\rho `$ of the quantum state of the nongravitational quantum field in the classical curved metric.
If one applies this definition to all possible states of this form, then it seems one can easily violate a bound of the form (1) by a state with arbitrarily large entropy by having the $`r=R`$ sphere be a neck separating the asymptotically flat exterior with an interior on the $`t=0`$ hypersurface that is almost an entire three-sphere of arbitrarily large size filled with thermal radiation. In other words, take the interior of the $`r=R`$ two-sphere to be the moment of maximum expansion of an almost-complete large $`k=1`$ radiation Friedman-Robertson-Walker model, and take the exterior to be a moment of time-symmetry of a nearly empty approximately Schwarzschild metric. If the interior three-sphere radius is $`aR`$, giving an interior volume going as $`a^3`$, the semiclassical Einstein equations imply that the radiation energy density must go as $`a^2`$ in Planck units at the moment of maximum expansion, so the temperature $`T`$ goes as $`a^{1/2}`$ and the entropy density goes as $`T^3`$ or as $`a^{3/2}`$. When this is multiplied by the volume, one gets an entropy going as $`a^{3/2}`$, which can be made arbitrarily large for fixed $`R`$ and $`E`$ (the asymptotic ADM mass) by making the interior size $`a`$ arbitrarily large.
One might seek to avoid this violation of (1) by restricting the states to which it is conjectured to apply to exclude this example of a huge interior universe separated from an asymptotically flat exterior by a relatively small neck.
One way to do that would be to demand that inside the $`r=R`$ two-sphere on the hypersurface of time symmetry, there are no round two-spheres of radius greater than $`R`$ (i.e., topological two-spheres with intrinsic two-metrics that are those of the standard unit round two-sphere multiplied by a constant $`r^2`$ that is larger than $`R^2`$). This would exclude an interior that is approximately a large round three-sphere of radius $`aR`$, since such an interior region would have round two-spheres of radii $`ra`$. However, it still appears to allow a very long throat of radius near $`R`$, which by its arbitrarily great length could have arbitrarily large volume and hence arbitrarily large entropy $`S`$ for fixed $`R`$ and $`E`$.
Another way to restrict the states so that they might possibly obey a bound similar to (1) is to demand that the evolution of the semiclassical geometry give a nonsingular metric over the whole of $`𝐑^4`$. This would exclude the examples of a large internal approximate three-sphere and also the long internal throat, since these examples would be expected to collapse gravitationally to singularities. Only in cases in which the metric is not too much different from flat spacetime would one expect that no singularities develop from gravitational collapse, and in these cases one might expect an entropy bound not too different from its flat spacetime form. However, for the restricted states to be sufficiently broad to encompass most of the semiclassical gravity generalizations of the allowed states (e.g., vacuum-outside-$`R`$ states) in the nongravitational theory, the semiclassical Einstein equations should give nonsingular evolution in these cases of sufficiently weak gravity. Since the semiclassical Einstein equations are of higher order in time than the ordinary Einstein equations with a classical source, it is not clear that this will be the case, and so one might need to find a suitable semiclassical gravity theory before trying to apply entropy bounds.
## 9 Conclusions and Acknowledgments
In conclusion, we have found that one can formulate precise definitions for entropy bounds of a complete quantum field system (i.e., one not restricted to the interior of some boundary) by giving precise definitions for the size $`R`$ of the system, at least when the metric is classical so that sizes can be unambiguously defined. In particular, $`R`$ may be defined for vacuum-outside-$`R`$ states as the largest round two-sphere on a suitable $`t=0`$ hypersurface, outside of which all of the operators have the same expectation values as in a suitable vacuum state (e.g., the ordinary vacuum state for nongravitational fields in Minkowski spacetime, or a Boulware-type quantum state in semiclassical gravity). Other values of $`R`$ may also be defined, such as the rms value of $`r`$ with a suitable weight function dependent upon the quantum state of the field. On the other hand, for a fully quantum gravity theory, it appears to be difficult to give an unambiguous definition of a size $`R`$ of a system, so it is not clear there how to define a bound for the entropy $`S`$ in terms of the energy $`E`$ and a size $`R`$.
Discussions with Valeri Frolov, Jonathan Oppenheim, and L. Sriramkumar have been helpful. This work was supported in part by the Natural Sciences and Engineering Research Council of Canada.
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# 1 Introduction
## 1 Introduction
The field equations
$`R_{\mu \nu }=\stackrel{~}{R}_{\mu \nu }F_{\mu \nu }=\stackrel{~}{F}_{\mu \nu }`$ (1.1)
pertain to many aspects of physics and mathematical physics: self-dual field theory, string theory, instantons and monopoles, and the classification of four-manifolds.
The $`N=2`$ worldsheet supersymmetric string is unique among string theories as its critical dimension is four . Its full spectrum and its exact (in $`\alpha ^{}`$) classical field equations have been identified to be merely those of self-dual gravity and self-dual Yang-Mills theory (for a review up to 1992, see ). Several quantum field theory formulations of the latter theories point to a non-vanishing S-matrix at the quantum level , a fact which is, however, in contrast to the claims of zero quantum S-matrix for the $`N=2`$ string . Apparently, we witness a quantum discrepancy between two theories which are classically equivalent. In this work we address this question by calculating the $`N=(2,2)`$ closed-string genus-one amplitudes in the RNS formulation and identifying the target spacetime theory which gives rise to these amplitudes.
Self-duality in $`d=2+2`$ dimensions (or in $`d=4+0`$) is implemented in the field equations of gravity and Yang-Mills by (1.1), of which the only known Lorentz covariant Lagrangian formulation employs Lagrange multipliers and is given by <sup>1</sup><sup>1</sup>1 Sub- and superscripts $`\alpha \{+,\}`$ and $`\stackrel{\text{.}}{\alpha }\{\stackrel{\text{.}}{+},\stackrel{\text{.}}{}\}`$ are spinor indices of $`SL(2,)`$ (or $`SU(2)`$).
$`=\mathrm{Tr}G^{\alpha \beta }F_{\alpha \beta }=\mathrm{Tr}\rho ^{\alpha \beta }R_{\alpha \beta },`$ (1.2)
where $`F`$ and $`R`$ are the self-dual projections of the field-strength and Riemann tensor for a gauge and spin connection vector, respectively. . The Lagrangian (1.2) involves two fields, related to the two polarization states of a gauge field, yet only one appears as an asymptotic state.
Alternatively, fixing a light-cone gauge in (1.1) allowed Leznov and Plebanski to reduce the self-duality equations to a pair of second-order equations
$`\mathrm{}\varphi +\frac{g}{2}[_+{}_{}{}^{\stackrel{\text{.}}{\alpha }}\varphi ,_{+\stackrel{\text{.}}{\alpha }}\varphi ]=0\mathrm{}\psi +\frac{\kappa }{2}_+{}_{}{}^{\stackrel{\text{.}}{\alpha }}_{+}^{}{}_{}{}^{\stackrel{\text{.}}{\beta }}\psi _{+\stackrel{\text{.}}{\alpha }}_{+\stackrel{\text{.}}{\beta }}\psi =0`$ (1.3)
for scalar prepotentials $`\varphi `$ (to $`F`$) and $`\psi `$ (to $`R`$) which extremize the respective Lorentz non-covariant gauge-fixed actions belonging to
$`=\mathrm{Tr}\varphi \left(\frac{1}{2}\mathrm{}\varphi +\frac{g}{6}[_+{}_{}{}^{\stackrel{\text{.}}{\alpha }}\varphi ,_{+\stackrel{\text{.}}{\alpha }}\varphi ]\right)`$ (1.4)
and
$`=\psi \left(\frac{1}{2}\mathrm{}\psi +\frac{\kappa }{6}_+{}_{}{}^{\stackrel{\text{.}}{\alpha }}_{+}^{}{}_{}{}^{\stackrel{\text{.}}{\beta }}\psi _{+\stackrel{\text{.}}{\alpha }}_{+\stackrel{\text{.}}{\beta }}\psi \right).`$ (1.5)
Further Lorentz non-covariant formulations of self-dual quantum field theories can be found by solving the gauge constraints in (1.1) differently . As these one-field actions share a coupling constant of positive length dimension they are all power-counting non-renormalizable.
The Lorentz-covariant two-field actions are much better behaved in this respect. In light-cone gauge, their Lagrangians are
$`=\mathrm{Tr}\stackrel{~}{\varphi }\left(\mathrm{}\varphi +\frac{g}{2}[_+{}_{}{}^{\stackrel{\text{.}}{\alpha }}\varphi ,_{+\stackrel{\text{.}}{\alpha }}\varphi ]\right)`$ (1.6)
and
$`=\stackrel{~}{\psi }\left(\mathrm{}\psi +\frac{\kappa }{2}_+{}_{}{}^{\stackrel{\text{.}}{\alpha }}_{+}^{}{}_{}{}^{\stackrel{\text{.}}{\beta }}\psi _{+\stackrel{\text{.}}{\alpha }}_{+\stackrel{\text{.}}{\beta }}\psi \right)`$ (1.7)
which allows no scattering beyond one-loop, because the multiplier fields go with $`1/\mathrm{}`$. The one- and two-field theories both generate the maximally helicity violating (MHV) scattering at one-loop and the vanishing next-to-MHV amplitudes at tree-level , and the latter theories are one-loop exact perturbatively.
To compare with, the $`N=2`$ superstring has been shown (modulo contact term ambiguities) to possess trivial scattering in its critical dimension . This indicates the presence of an anomaly in the string, or a target-space interpretation different from self-dual gravity or gauge theory. A possible anomaly interpretation behind the $`d=3+1`$ MHV amplitudes in gauge theory was initially pointed out in in the context of the conserved symmetries of the field equations.
Until now, the $`N=2`$ string quantum amplitude has never been computed in the traditional RNS formalism (but functional methods for the quantization at higher genera have been developed ). However, by embedding the $`N=2`$ string in an $`N=4`$ topological string it was demonstrated that, up to contact terms, these amplitudes vanish to all loops . Linearized symmetry arguments have also formally shown this in the RNS formulation in . In order to compare with the field-theory results and to find the root of this discrepancy, an explicit traditional computation at genus one is worthwhile. In the present work we perform this calculation. We find that the $`N=2`$ string loop dynamics appears to be reduced to two dimensions. Based on earlier one-loop computations of the partition function and the three-point function , Marcus already identified the technical origin of this dimensional mismatch. Here, we confirm his observation and extend it to the full quantum dynamics by evaluating the one-loop four-point scattering. As the $`N=2`$ string has a critical dimension of four, with four (real) target spacetime coordinates, this calculation indicates that it represents a ghost system in the MHV sector of gauge theory. In order to make the above explicit we render the scattering manifestly Lorentz invariant by normalizing the vertex operators and incorporating the gauge invariance through the use of spinor helicity techniques.
A further unexpected relation arises between the one-loop MHV amplitude in pure gravity regulated to two dimensions and the next-to-MHV amplitude in IIB supergravity evaluated in ten dimensions.<sup>2</sup><sup>2</sup>2 This dimension-shifting relation involving a change in the number of supersymmetries was initially found in . The absence of one-loop divergences in the massless sector of IIB supergravity in ten dimensions within dimensional reduction explains the vanishing of the two-dimensional MHV result and thereby the triviality of the field theory limit of the $`N=2`$ string scattering. Alternatively, a relation is found between two string theories: the IIB superstring in ten dimensions and the $`N=2`$ closed string in four dimensions. Such a relation may originate in an integrable structure in the ultra-violet regime for the massless modes of the string (a supergravity analog of the Regge kinematical limit of Yang-Mills theory). If this connection extends to multi-loops, the vanishing theorems of the $`N=2`$ string at higher genera deserve further study.
The outline of this work is as follows. In section 2 we review and discuss the properties of the gauge theory MHV amplitudes in different dimensions. Section 3 implements gauge invariance directly into the $`N=2`$ string scattering through the incorporation of spinor helicity techniques and a normalization of the vertex operators. We also analyze contact term subtleties in the scattering. In section 4 the three-point genus-one closed-string amplitude is obtained and compared with the field theoretic one. In section 5 we finally compute a modular-integral expression for the genus-one closed-string four-point amplitude, carefully taking into account the superconformal ghost structure. From this result, we extract the field theory limit by taking $`\alpha ^{}0`$ before summing over spin structures (which then trivializes). The answer is zero, and the comparison with the MHV amplitudes is made. In section 6 we explicitly perform the spin structure summation on the full modular integrand before taking the field theory limit, with identical (vanishing) result. A discussion and an Appendix on Jacobi theta functions conclude the paper.
## 2 Review of MHV Gauge Theory Amplitudes
Recent developments in techniques in gauge theory calculations <sup>3</sup><sup>3</sup>3 For a review at tree-level see and at loop-level . have made possible the calculation of closed analytic forms of several infinite sequences of one-loop gauge theory amplitudes. The maximally helicity violating (MHV) amplitudes are described by scattering of gauge fields of identical helicity, either in Yang-Mills theory or in gravity. One of the features of these amplitudes is that in a supersymmetric theory they are identically zero to infinite loop order; this implies that at tree-level the amplitudes are identically zero. The amplitudes closest to MHV are simpler to calculate, and the self-dual description has lead to reformulations and improved diagrammatic techniques in calculating gauge theory amplitudes as well as second-order formulations for incorporating fermions <sup>4</sup><sup>4</sup>4 In a self-dual non-abelian background, fermions may be bosonized, and fields become spin independent. . In this section we briefly review these maximally helicity violating amplitudes and describe their relations to both self-dual field theory and string theory. The continuation of the four-point MHV gravity amplitude and its conjectured form to $`n`$-point order to arbitrary dimensions is directly related to the zero-slope limit of the $`N=2`$ closed string.
At one-loop in four dimensions, the leading-in-color partial amplitude for the scattering of $`n`$ gluons of identical out-going helicity in Yang-Mills theory is
$`A_{n;1}^{[1]}(k_i)={\displaystyle \frac{i}{48\pi ^2}}{\displaystyle \underset{1i<j<k<ln}{}}{\displaystyle \frac{ij[jk]kl[li]}{1223\mathrm{}n1}}`$ (2.1)
where the superscript $`[J]`$ represents the spin of the internal state (gluon, complex scalar or Weyl fermion give the same result up to a minus sign for half-integral spin). The amplitude is written in color-ordered form ; the leading-in-color group theory structure,
$`N^2\mathrm{Tr}\mathrm{T}^{a_1}\mathrm{T}^{a_2}\mathrm{}\mathrm{T}^{a_n},`$ (2.2)
has been extracted from the kinematics in accord with Chan-Paton assignments in open string theory and gauge theory. In (2.1) we have decomposed each lightlike momentum vector $`k_i`$ into two momentum Weyl spinors and defined two different inner products,
$`k_i^{\alpha \stackrel{\text{.}}{\alpha }}=k_i^\alpha k_i^{\stackrel{\text{.}}{\alpha }}\mathrm{and}ij=k_i^\alpha k_{j,\alpha }[ij]=k_i^{\dot{\alpha }}k_{j,\dot{\alpha }}.`$ (2.3)
In $`2+2`$ dimensions, $`ij`$ is not the complex conjugate of $`[ij]`$; the Lorentz group is $`SL(2,)\times SL(2,)^{}`$ as opposed to $`SL(2,)`$ in $`3+1`$ dimensions (and there are poles in (2.1) in the self-dual plane parameterized by the $`SL(2,)`$ half of the Lorentz group).
The analogous result for the all-plus gravitational amplitude ,
$`A_4^{[2]}(k_i)=i\left({\displaystyle \frac{\kappa }{2}}\right)^4{\displaystyle \frac{1}{120(4\pi )^2}}\left({\displaystyle \frac{s_{12}s_{23}}{12233441}}\right)^2(s_{12}^2+s_{23}^2+s_{13}^2),`$ (2.4)
and its $`n`$-point form , together with the three-point vertex, describe the scattering of all-plus helicity gravitons to one-loop order. The amplitude in (2.1) and its gravitational analog have a number of features in common with $`N=2`$ string scattering. They are channel-dual in the sense that exchange of any two legs gives the same form. Furthermore, they have only two-particle poles (in one $`SL(2,)`$ factor of the Lorentz group), which signals integrable characteristics related to the infinite number of symmetries in the self-dual field equations.
The $`n`$-point gauge theory amplitude in (2.1) has been found by constraining the functional form based on analyticity as well as through a direct calculation with a fermion in the loop . The amplitude in (2.1) also arises in a one-loop S-matrix element for self-dual Yang-Mills theory. This happens for the Lorentz-covariant two-field theory (one-loop exact) as well as, to a factor of two, for the one-field (Leznov) formulation , although the latter is not Lorentz covariant (or one-loop exact). The same story occurs in gravity where (2.4) and its generalizations describe quantum self-dual gravity at one-loop . One might expect to find similar non-vanishing scattering amplitudes at the one-loop level in the $`N=2`$ string, as both the string and the self-dual field theory share the same classical field equations. However, this expectation is not borne out by our calculation below.
The $`d`$ dimensional generalization of the Yang-Mills result in (2.1) has been found in up to six-point (together with a conjectured form at $`n7`$ point), and we list here the form of these amplitudes. At four-point one has
$`A_{4;1}^{[1]}(k_i)={\displaystyle \frac{2i}{12233441}}{\displaystyle \frac{(4d)(2d)}{4(4\pi )^2}}s_{12}s_{23}I_4^{4+d}(s,t),`$ (2.5)
where the box diagram $`I_4^{4+d}`$ is the integral function
$`I_4^p(s,t)={\displaystyle \frac{d^p\mathrm{}}{(2\pi )^p}\frac{1}{\mathrm{}^2(\mathrm{}k_1)^2(\mathrm{}k_1k_2)^2(\mathrm{}+k_4)^2}},`$ (2.6)
continued from $`p`$ to $`d+4`$ dimensions but with the external vectors in $`d`$ dimensions. The generalization of the series in (2.1) arises by keeping the external kinematics and polarizations in four dimensions and analytically continuing the scalar integral functions. In a Schwinger proper-time formulation of the integrals this amounts to inserting additional factors of $`\tau _2`$ in the integral over the proper time. In four dimensions, the $`8`$-dimensional box diagram in (2.5) relevant to the amplitude is UV divergent, but the result is finite because the $`d4`$ prefactor extracts the residue. In two dimensions the $`6`$-dimensional box diagram with external massless kinematics is both IR and UV finite, but the prefactor forces the result in (2.5) to be identically zero. The MHV amplitude thus vanishes upon continuation to $`d=2`$, without recourse to spacetime supersymmetry.
The five- and six-point amplitudes and their dimensional form have the same properties as the expression (2.5), as does the conjectured $`n`$-point form at one-loop. The five-point amplitude,
$`A_{5;1}^{[1]}(k_i)`$ $`=`$ $`{\displaystyle \frac{i}{1223344551}}{\displaystyle \frac{(4d)(2d)}{4(4\pi )^{d/2}}}[s_{23}s_{34}I_4^{d+4}+s_{34}s_{45}I_4^{d+4}`$ (2.7)
$`+`$ $`s_{45}s_{51}I_4^{d+4}+s_{51}s_{12}I_4^{d+4}+s_{12}s_{23}I_4^{d+4}+4idϵ_{\mu \nu \rho \sigma }k_1^\mu k_2^\nu k_3^\rho k_4^\sigma I_5^{d+6}],`$
is zero when continued to two dimensions because the six-dimensional box and eight-dimensional pentagon in (2.7) are finite and the pre-factor vanishes in $`d=2`$. The gauge theory result at six-point is similar and is described in eqs. (16) and (17) of reference .
In (spacetime) supersymmetric gauge or gravitational theory, the MHV one-loop amplitudes vanish because of a cancellation between the contributions stemming from different spin states running inside the loop . Concretely,
$`A^{[1]}=A^{[0]}=A^{[\frac{1}{2}]}=A^{[2]},`$ (2.8)
for a gauge boson, complex scalar, Weyl fermion, or graviton, so that amplitudes need to be computed for only one conveniently chosen spin value.
The all-$`n`$ conjectured form of the MHV Yang-Mills amplitude relates to a $`d+4`$ $`𝒩=16`$ supersymmetric non-MHV amplitude as follows,
$`A_{n;1}^{[0]}(k_i)|_d={\displaystyle \frac{(4d)(2d)}{2}}(4\pi )^2{\displaystyle \frac{1}{12^4}}A_{n;1}^{𝒩=16}(k_1^{},k_2^{},k_3^+,\mathrm{},k_n^+)|_{d+4},`$ (2.9)
where for definiteness we denote it for an internal complex scalar, with $`[J=0]`$. The factor of $`12^4`$ gives the left-hand side the appropriate spinor weight to describe the negative-helicity gluons on legs one and two. Curiously, the prefactor in (2.9) is negative for $`2d4`$. Again, the finiteness of the amplitude on the right-hand side of (2.9) in $`d+4=6`$ translates into the vanishing of the MHV amplitude in $`d=2`$. For $`d+4=8`$ the UV singularity of $`A_n^{𝒩=16}`$ reproduces (2.1).
The explicit result in $`d`$ dimensions for the four-point one-loop maximally helicity violating Einstein-Hilbert gravitational amplitude is
$`A_4^{[2]}(k_i)|_d={\displaystyle \frac{(4d)(2d)d(2+d)}{8}}(4\pi )^4{\displaystyle \frac{1}{12^8}}A_4^{𝒩=32}(1^{},2^{},3^{++},4^{++})|_{d+8}`$ (2.10)
where the relation is between an MHV amplitude in $`d`$ dimensions to a non-MHV amplitude in $`d+8`$ dimensions and in the $`𝒩=32`$ (maximally) supersymmetric theory. Similar to the Yang-Mills case, the additional $`12^4`$ gives the MHV amplitude the proper helicity weight (the graviton has twice the spin) and dimensions.
For $`d=2`$ the amplitude on the right-hand side in (2.10) is to be evaluated in ten dimensions. In this case no counterterms occur in the amplitude calculation in dimensional regulation since the divergences at four-point are proportional to
$`\left({\displaystyle \frac{1}{d10}}\right)(s+t)+\left({\displaystyle \frac{1}{d10}}\right)(t+u)+\left({\displaystyle \frac{1}{d10}}\right)(u+s)`$ (2.11)
which is zero on-shell, forcing the MHV result in (2.10) in $`d=2`$ to vanish. Parallel to the relation in (2.9) and generalizing (2.10), the conjectured $`d`$-dimensional gravitational MHV amplitude at arbitrary $`n`$-point coincides with the $`𝒩=32`$, $`d+8`$ next-to-MHV amplitude. The absence of a counterterm at $`n`$-point in the dimensionally regularized/reduced form of IIB supergravity in ten dimensions means that, due to the prefactor in the $`n`$-point generalization of (2.10), the MHV result for graviton scattering in two dimensions is zero at arbitrary $`n`$-point order at one-loop.
Two-dimensional gravity and Yang-Mills theory are topological and have no dynamical degrees of freedom. The scattering in these theories is trivial in topologically trivial spacetime, which explains the vanishing of the amplitudes not only at one-loop but also to infinite loop order. A possible relation between the reduced form of the scattering in $`d=2`$ and that in $`d=10`$ implies further non-trivial structure in the ultra-violet of IIB supergravity.
In the following we shall relate the above $`d=2`$ result in the gravitational case to the scattering obtained in the RNS formulation of the closed $`N=2`$ superstring in the zero-slope limit. Given the holomorphic/anti-holomorphic factorization of the string integrand, this relation might persist to the open string as well.
## 3 N=2 String Vertex Operators
In this section we review the relevant facts of the closed $`N=2`$ string and its tree-level scattering amplitudes. We pay particular attention to the its vertex operators, for two reasons: First, the representation of the vertex operators affects possible contact interactions and their contributions to scattering amplitudes. Second, the normalization of the vertex operators translates to the choice of external leg factors which are crucial to achieve a manifestly gauge-invariant representation of the amplitudes via spinor helicity techniques. For a brief review, the reader may consult and references therein.
3.1 Generalities
From the worldsheet point of view, critical closed $`N=2`$ strings in flat Kleinian space $`^{2,2}`$ are a theory of $`N=(2,2)`$ supergravity on a $`1+1`$ dimensional (pseudo) Riemann surface, coupled to two chiral $`N=(2,2)`$ massless matter multiplets $`X^a`$, $`a=1,2`$. The latter’s components are complex scalars $`x`$ (the four string coordinates), $`SO(1,1)`$ Dirac spinors $`\psi `$ (their four NSR partners) and complex auxiliaries $`F`$,
$`X^a=x^a+\theta ^{}\psi ^{+a}+\theta ^+\psi ^{\overline{a}}+\theta ^+\theta ^{}F^a`$ (3.1)
with arguments $`yz+\theta ^+\theta ^{}`$. Complex conjugation reads
$`z^{}=z(\theta ^+)^{}=\theta ^{}(x^a)^{}=x^{\overline{a}}(\psi ^{+a})^{}=\psi ^{\overline{a}}`$ (3.2)
while chiral conjugation exchanges right- and left-movers via
$`z\overline{z}\theta ^\pm \overline{\theta }^\pm x^ax^a\psi ^{+a}\overline{\psi }^{+a}.`$ (3.3)
The extended worldsheet supersymmetry has induced a spacetime complex structure which reduces the global Lorentz symmetry,
$`\mathrm{Spin}(2,2)=SU(1,1)\times SU(1,1)^{}U(1)\times SU(1,1)^{}U(1,1).`$ (3.4)
In superconformal gauge, however, manifest $`SO(2,2)`$ symmetry is restored in the worldsheet action, which is given by
$`S={\displaystyle d^2zd^2\theta d^2\overline{\theta }K(X,\overline{X})}={\displaystyle d^2z\eta _{a\overline{a}}[x^a\overline{}x^{\overline{a}}+\psi ^{+a}\overline{}\psi ^{\overline{a}}+\overline{\psi }^{+a}\overline{\psi }^{\overline{a}}]}`$ (3.5)
where $`\eta _{a\overline{a}}=\mathrm{diag}(+)`$ is the flat metric in $`^{1,1}`$, and the auxiliary fields have been integrated out.
Although the above notation makes transparent the local R symmetry properties of the fields (for instance, $`x`$ is neutral while $`\psi ^\pm `$ is not), it is not convenient for our computations. The interrelation (3.2) with complex conjugation allows us to change it,
$`x^ax^{+a}x^{\overline{a}}x^a\psi ^{+a}\psi ^{+a}\psi ^{\overline{a}}\psi ^a,`$ (3.6)
so that the $`SO(2,2)`$ invariant scalar product reads
$`kx=\frac{1}{2}(k^+x^{}+k^{}x^+)=\frac{1}{2}(k^{+1}x^1k^{+2}x^2+k^1x^{+1}k^2x^{+2})`$ (3.7)
where the dot is also used to denote the $`SU(1,1)^{}`$ invariant scalar product. There exist three antisymmetric $`SU(1,1)^{}`$ invariant products,
$`k^+x^+`$ $`=`$ $`ϵ_{ab}k^{+a}x^{+b}=k^{+1}x^{+2}k^{+2}x^{+1}`$
$`k^+x^{}`$ $`=`$ $`\frac{1}{2}(k^+x^{}k^{}x^+)=\frac{1}{2}(k^{+1}x^1k^{+2}x^2k^1x^{+1}+k^2x^{+2})`$
$`k^{}x^{}`$ $`=`$ $`ϵ_{ab}k^ax^b=k^1x^2k^2x^1`$ (3.8)
which feature prominently in the following.
The $`N=(2,2)`$ supergravity multiplet defines a gravitini and a Maxwell bundle over the worldsheet Riemann surface. The topology of the total space is labeled by the Euler number $`\chi `$ of the punctured Riemann surface and the first Chern number (instanton number) $`M`$ of the Maxwell bundle. It is notationally convenient to replace the Euler number by the “spin”
$`J:=2\chi =2n4+4(\mathrm{\#}\mathrm{handles})2.`$ (3.9)
The action (3.5) is to be considered for string worldsheets of a given topology.<sup>5</sup><sup>5</sup>5 Of course, the Lagrangian in (3.5) is in general not correct globally. The first-quantized string path integral for the $`n`$-point function $`A_n`$ includes a sum over worldsheet topologies $`(J,M)`$, weighted with appropriate powers in the string couplings $`(\kappa ,e^{i\theta })`$:
$`A_n(\kappa ,\theta )={\displaystyle \underset{J=2n4}{\overset{\mathrm{}}{}}}\kappa ^{J/2}A_n^J(\theta )={\displaystyle \underset{J=2n4}{\overset{\mathrm{}}{}}}{\displaystyle \underset{M=J}{\overset{+J}{}}}\kappa ^{J/2}e^{iM\theta }A_n^{J,M}`$ (3.10)
where the instanton sum has a finite range because bundles with $`|M|>J`$ do not contribute. The presence of Maxwell instantons breaks the explicit $`U(1)`$ factor in (3.4) but the $`SU(1,1)`$ factor (and thus the whole $`\mathrm{Spin}(2,2)`$) is fully restored if we let $`\kappa ^{1/4}(e^{i\theta /2},e^{i\theta /2})`$ transform as an $`SU(1,1)`$ spinor. As a consequence, the string couplings depend on the $`SO(2,2)`$ Lorentz frame, and we may choose a convenient one for calculations. We call the choice $`\theta =0`$ a ‘Leznov frame’ and name an averaging over $`\theta `$ a ‘Yang frame’. The partial amplitudes $`A_n^{J,M}`$ are integrals over the metric, gravitini, and Maxwell moduli spaces. The integrands may be obtained as correlation functions of vertex operators in the $`N=(2,2)`$ superconformal field theory on the worldsheet surface of fixed shape (moduli) and topology.
The vertex operators produce from the (first-quantized) vacuum state the asymptotic string states in the scattering amplitude under consideration. They correspond to the physical states of the $`N=2`$ closed string and carry their quantum numbers. Being representatives of the (semi-chiral) BRST cohomology, they are unique only up to BRST-trivial terms and normalization. The physical subspace of the $`N=2`$ string Fock space in a covariant quantization scheme turns out to be surprisingly small : Only the ground state $`|k`$ remains, a scalar on the massless level, i.e. for center-of-mass momentum $`k^{\pm a}`$ with $`kk=0`$. The dynamics of this string “excitation” is described by a massless scalar field,
$`\mathrm{\Phi }(x)={\displaystyle d^4ke^{ikx}\stackrel{~}{\mathrm{\Phi }}(k)},`$ (3.11)
whose self-interactions are determined on-shell from the (amputated) string scattering amplitudes at tree-level,
$`\stackrel{~}{\mathrm{\Phi }}(k_1)\mathrm{}\stackrel{~}{\mathrm{\Phi }}(k_n)_{\mathrm{tree},\theta }^{\mathrm{amp}}=:A_n^{2n4}(k_1\mathrm{}k_n;\theta )=:\delta _{k_1+\mathrm{}+k_n}\stackrel{~}{A}_n^{2n4}(k_1\mathrm{}k_n;\theta ).`$ (3.12)
Interestingly, it has been shown that all tree-level $`n`$-point functions vanish on-shell, except for the three-point amplitude ,
$`\stackrel{~}{A}_3^2(k_1,k_2,k_3;\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[e^{i\theta }k_1^+k_2^+2k_1^+k_2^{}e^{i\theta }k_1^{}k_2^{}\right]^2`$ (3.13)
with $`k_ik_j=0`$ due to $`_nk_n=0`$. Note that $`\stackrel{~}{A}_3^2`$ is totally symmetric in all momenta. Expanding the square, one reads off $`\stackrel{~}{A}_3^{2,M}`$ for $`M=2,\mathrm{},+2`$. However, using the on-shell relations
$`k_1^+k_2^{}=h(k)^{}k_1^+k_2^+=h(k)k_1^{}k_2^{}`$ (3.14)
with the phase
$`h(k):={\displaystyle \frac{k^{+1}}{k^2}}={\displaystyle \frac{k^{+2}}{k^1}}=1/h(k)^{}`$ (3.15)
(identical for all three momenta), the three-point amplitude simplifies to
$`\stackrel{~}{A}_3^2(k_1,k_2,k_3;\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4}}[h(k)^{1/2}e^{i\theta /2}h(k)^{1/2}e^{i\theta /2}]^4(k_1^+k_2^{})^2`$ (3.16)
$`=`$ $`{\displaystyle \frac{1}{4}}e^{2i\theta }[1h(k)^1e^{i\theta /2}]^4(k_1^+k_2^+)^2.`$
We see that the $`\theta `$ dependence factorizes, and the contributions from different instanton sectors differ only by powers of the leg factor $`h(k)`$. After switching to real $`SL(2,)\times SL(2,)^{}`$ spinor coordinates, it is easy to see that this three-point tree-level amplitude exactly coincides, in the Leznov frame, with the one obtained from Plebanski’s second equation (1.3) for the prepotential $`\psi `$. In the Yang frame, one makes contact with Plebanski’s first equation. Furthermore, after including the appropriate leg factors the result becomes identical to covariant gauge scattering.
The above structure of the $`\theta `$ dependence is not a speciality of the tree-level three-point function but actually a generic property. One may localize the Maxwell instantons at the worldsheet punctures and thereby define vertex operators $`V^M`$, $`M=J,\mathrm{},J`$ for various instanton sectors which create an asymptotic string state together with a Maxwell instanton out of the $`M=0`$ vacuum. Yet, it turns out that any two such operators are proportional to each other, differing merely by (momentum-dependent) normalization,
$`V^M(k)=h(k)^MV(k)`$ (3.17)
where $`V(k)`$ is the vertex operator in the zero-instanton sector. It follows that the partial amplitudes (tree or loop) in the various instanton sectors are related by simple leg factors, and that knowledge of a particular $`A^{J,M}`$ is sufficient. For this reason, we shall be content to perform our calculations in the zero-instanton sector, except in section four where we employ a Leznov frame.
3.2 Avoiding Contact Terms
The canonical computation of one-loop amplitudes entails the use of the integrated ground state vertex operator in the $`(0,0;0,0)`$ superconformal ghost picture. Its standard representative is
$`\stackrel{~}{V}(k)`$ $`=`$ $`{\displaystyle d^2zd^2\theta d^2\overline{\theta }\mathrm{exp}(ikX)}`$
$`=`$ $`{\displaystyle d^2z(k^{[+}x^]ik^{}\psi ^{}k^{}\psi ^+)(k^{[+}\overline{}x^]+ik^+\overline{\psi }^{}k^{}\overline{\psi }^+)e^{ikx}}.`$
The use of this vertex operator in amplitude calculations gives rise to delta functions (and squares of delta functions) on the string worldsheet because of holomorphic/antiholomorphic Wick contractions
$`x^{+a}(z_1)\overline{}\overline{x}^b(z_2)=\eta ^{ab}\delta ^{(2)}(z_1z_2).`$ (3.19)
These contact terms are usually dropped in perturbation theory, but care must be taken to ensure that these terms do not contribute to the scattering in any representation.<sup>6</sup><sup>6</sup>6 These contact terms are proportional, after the incorporation of helicity techniques, to inner products $`ϵ_i\overline{ϵ}_j`$ which vanish manifestly in the MHV amplitudes. It is possible to completely avoid such contact terms by changing the vertex operator representative. Adding the total derivative term
$$i\left[(k^{[+}\overline{}x^]+ik^+\overline{\psi }^{}k^{}\overline{\psi }^+)e^{ikx}\right]i\overline{}\left[(k^{[+}x^]ik^{}\psi ^{}k^{}\psi ^+)e^{ikx}\right]\overline{}e^{ikx}$$
(3.20)
and using
$$e^{ikx}=k^{(+}x^)e^{ikx}$$
(3.21)
we arrive at
$`V(k)={\displaystyle d^2z(2k^+x^{}ik^+\psi ^{}k^{}\psi ^+)(2k^+\overline{}x^{}+ik^+\overline{\psi }^{}k^{}\overline{\psi }^+)e^{ikx}}`$ (3.22)
which contains $`x^+`$ only in the exponent and therefore precludes not only $`x\overline{}x`$ but also $`xx`$ and $`\overline{}x\overline{}x`$ contractions. In the following we shall use this vertex operator. There is one drawback, however. Since $`V(k)`$ in (3.22) is no longer invariant under complex conjugation, our computations will not produce holomorphic squares, making chiral splitting impossible.
Next we derive the unintegrated weighted generating functional (Koba-Nielsen form) for $`n`$-point amplitudes. The bosonic portion is
$$\underset{j=1}{\overset{n}{}}d\theta _jd\overline{\theta }_j𝑑\mu _n\mathrm{exp}\left[d^2zd^2\stackrel{~}{z}J^+(z)G(z,\stackrel{~}{z})J^{}(\stackrel{~}{z})\right].$$
(3.23)
Here, $`\theta _j`$ correspond to an exponentiation
$`k^+x^{}e^{kx}=\mathrm{exp}[kx+\theta k^+x^{}]|_{\mathrm{multi}\mathrm{linear}}`$ (3.24)
of the pre-factor in the vertex operator from which subsequently (after functional integration) the multi-linear part is extracted to obtain the correlation. For the chirally non-split form $`V`$ in (3.22),<sup>7</sup><sup>7</sup>7 The real form $`\stackrel{~}{V}`$ of the vertex operator in (3.22) leads to $`J^{}(z)`$ being the complex conjugate of (3.25). the currents are
$`J^+(z)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}\left[ik_j^+\delta ^{(2)}(zz_j)+\theta _jk_j^+\delta ^{(2)}(zz_j)+\overline{\theta }_jk_j^+\overline{}\delta ^{(2)}(zz_j)\right]`$ (3.25)
$`J^{}(z)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}ik_j^{}\delta ^{(2)}(zz_j).`$ (3.26)
The sum in (3.23) may be evaluated to
$$\underset{j=1}{\overset{n}{}}d\theta _jd\overline{\theta }_j𝑑\mu _n\underset{ij}{}\mathrm{exp}\left[k_ik_jG_{ij}+i\theta _ik_i^+k_j^{}G_{ij}+i\overline{\theta }_ik_i^+k_j^{}\overline{}G_{ij}\right]$$
(3.27)
where $`G_{ij}=x^+(z_i,\overline{z}_i)x^{}(z_j,\overline{z}_j)`$, the bosonic two-point function on the torus, and $`d\mu _n`$ denotes the measure to integrate over the general punctured super-Riemann surface. The global $`N=2`$ superspace form generalizing that in (3.27) is
$$\underset{j=1}{\overset{n}{}}d\theta _jd\overline{\theta }_j𝑑\mu _n^s\underset{i<j}{}\mathrm{exp}\left[k_ik_jG_{ij}+i\theta _ik_i^+k_j^{}D_i^+G_{ij}+i\overline{\theta }_ik_i^+k_j^{}D_i^{}G_{ij}\right]|_{\mathrm{multi}\mathrm{linear}}$$
(3.28)
where $`d\mu _n^s`$ is the superspace measure and $`D^\pm `$ the $`N=2`$ superspace derivatives. The form in (3.28) is covariantized in the next section.
3.3 Gauge Invariance and Reference Momenta
In this subsection we describe the transversality of the amplitude at the level of the vertex operators and introduce the calculational tool of reference momenta in order make manifest the gauge invariance of the amplitudes. These instruments will allow us to compare the integrand with that of IIB superstring and gravity loop amplitudes. Spinor helicity is a useful tool in gauge theory calculations and implicitly has been incorporated in the $`N=2`$ string, although obscured in previous representations. Here, we find it convenient to switch to a real $`SL(2,)\times SL(2,)^{}`$ notation
$`v^{\alpha \dot{\alpha }}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}v^{+1}+v^1iv^{+2}+iv^2& iv^{+1}+iv^1+v^{+2}+v^2\\ iv^{+1}iv^1+v^{+2}+v^2& v^{+1}+v^1+iv^{+2}iv^2\end{array}\right)`$ (3.29)
for vectors and coordinates and rewrite the $`U(1,1)`$ scalar product as <sup>8</sup><sup>8</sup>8 Note that the $`\pm `$ superscripts have different meaning on left- and right-hand sides.
$`2v^+w^{}=ϵ_{\stackrel{\text{.}}{\alpha }\stackrel{\text{.}}{\beta }}\left[v^{+\stackrel{\text{.}}{\alpha }}w^{\stackrel{\text{.}}{\beta }}v^{\stackrel{\text{.}}{\alpha }}w^{+\stackrel{\text{.}}{\beta }}iv^{+\stackrel{\text{.}}{\alpha }}w^{+\stackrel{\text{.}}{\beta }}iv^{\stackrel{\text{.}}{\alpha }}w^{\stackrel{\text{.}}{\beta }}\right].`$ (3.30)
For a light-like momentum vector $`k^{\alpha \stackrel{\text{.}}{\alpha }}=k^\alpha k^{\stackrel{\text{.}}{\alpha }}`$, we have the freedom to choose the spinor $`q=q(k)`$ like <sup>9</sup><sup>9</sup>9 The matrix is degenerate and not related to the identity by a similarity transformation.
$`\left(\begin{array}{c}q_+\\ q_{}\end{array}\right)=\left(\begin{array}{cc}1& i\\ i& 1\end{array}\right)\left(\begin{array}{c}k_+\\ k_{}\end{array}\right)\mathrm{hence}q_+=iq_{},`$ (3.31)
which permits us to express
$`k^+v^{}=\frac{1}{2}q_\beta k_{\stackrel{\text{.}}{\beta }}v^{\beta \stackrel{\text{.}}{\beta }}`$ (3.32)
in $`SO(2,2)`$ covariant form. Two different spinors $`q_1^\alpha `$ and $`q_2^\alpha `$ related to momenta $`k_1`$ and $`k_2`$ further satisfy
$`q_1^+={\displaystyle \frac{k_1^+ik_1^{}}{k_2^+ik_2^{}}}q_2^+.`$ (3.33)
A representation of the two physical polarization vectors $`ϵ_{\alpha \stackrel{\text{.}}{\beta }}^\pm `$ in terms of products of spinors is
$`ϵ_{\alpha \dot{\beta }}^+(k;q)=i{\displaystyle \frac{q_\alpha k_{\dot{\beta }}}{q^\gamma k_\gamma }}ϵ_{\alpha \dot{\beta }}^{}(k;q)=i{\displaystyle \frac{q_{\dot{\beta }}k_\alpha }{q^{\dot{\gamma }}k_{\dot{\gamma }}}}`$ (3.34)
and has the following properties:
$`ϵ_{\alpha \dot{\beta }}^\pm (k;\stackrel{~}{q})=ϵ_{\alpha \dot{\beta }}^\pm (k;q)+f(\stackrel{~}{q},q;k)k_{\alpha \dot{\beta }},`$ (3.35)
$`ϵ^+ϵ^+=0ϵ^+ϵ^{}=1.`$ (3.36)
Because the choice of $`q`$ is arbitrary in any gauge-invariant calculation, it can be chosen to force many inner products to vanish, considerably reducing the amount of algebra in intermediate steps of the calculation.
For example, in an MHV amplitude calculation the individual reference momenta $`q_i`$ may be taken to coincide: $`q_i=q`$. This choice eliminates all inner products of polarization vectors,
$`ϵ^+(k_1;q)ϵ^+(k_2;q)=0.`$ (3.37)
Since individual diagrams must, by dimensional analysis, contain at least one inner product of two polarization vectors, the vanishing of the tree-level MHV (and next to MHV) amplitudes follows immediately. Because the next-to-MHV amplitudes describe the self-dual scattering at tree-level, this also shows the classical triviality of self-dual field theory scattering . At the loop-level it also allows a direct comparison between the N=2 string amplitude calculations and those in the field theory because no $`\overline{}G_{ij}`$ arises in the integral form in (3.28).
In order to compare we normalize the $`i^{\mathrm{th}}`$ vertex operator with an additional line factor,
$`V^{}(k_i)=\left({\displaystyle \frac{1}{q_i^\alpha k_{i,\alpha }}}\right)^2V(k_i)`$ (3.38)
with $`q_i`$ satisfying (3.31). By this step, $`V^{}`$ takes the same form as the type IIB superstring gravitational vertex operator,
$`V^{}(k,ϵ)={\displaystyle d^2zϵ_{\alpha \dot{\alpha }}^+ϵ_{\beta \dot{\beta }}^+(x^{\alpha \stackrel{\text{.}}{\alpha }}i\psi ^{\alpha \stackrel{\text{.}}{\alpha }}k^{}\psi ^+)(\overline{}x^{\beta \stackrel{\text{.}}{\beta }}+i\overline{\psi }^{\beta \stackrel{\text{.}}{\beta }}k^{}\overline{\psi }^+)e^{ikx}},`$ (3.39)
and is clearly Lorentz covariant due to the reference momenta property in (3.35). The graviton polarization in four dimensions ($`d=2+2`$) is identified after adjoining $`ϵ_{\alpha \dot{\alpha },\beta \dot{\beta }}^{++}(k)=ϵ_{\alpha \dot{\alpha }}^+(k)ϵ_{\beta \dot{\beta }}^+(k)`$. Since by (3.17) the vertex operator in a non-zero instanton sector is related to the one in (3.22) by a leg factor only, covariant versions of vertex operators can be given for any instanton sector by an appropriate choice of reference momenta.
The reference momenta defined in (3.31) for the different vertex operators satisfy
$`q_i^\alpha q_{j,\alpha }=0,`$ (3.40)
which means that this choice automatically nullifies all the different inner products $`ϵ^+(k_i;q_i)ϵ^+(k_j;q_j)=0`$. Other choices of reference momenta, e.g. $`q_j^\alpha =q^\alpha `$ for all external lines, may be obtained by a gauge transformation of the vertex operator after normalizing the external lines; they correspond to adding a longitudinal component in (3.35) and yield the same on-shell S-matrix elements.
With the representation in (3.39) the integrand is identical to the Koba-Nielsen representation of the IIB superstring, apart from the spin structure dependence,
$$𝑑\mu _n\underset{ij}{}\mathrm{exp}\left(k_ik_jG_{ij}\right)\underset{ij}{}\left|\mathrm{exp}\left[ϵ_{[i}k_{j]}_iG_{ij}+ϵ_iϵ_j_i_jG_{ij}+ϵ_i\overline{ϵ}_j_i\overline{}_jG_{ij}\right]\right|_{\mathrm{multi}\mathrm{linear}}^2$$
(3.41)
where the label ‘multi-linear’ means that the integrand is expanded in powers of the polarizations, keeping only the terms linear in each polarization ($`ϵ_j`$ or $`\overline{ϵ}_j`$). The $`N=1`$ superspace form has
$`_iG_{ij}D_+^iG_{ij}_i_jG_{ij}D_+^iD_+^jG_{ij}_i\overline{}_jG_{ij}D_+^iD_{}^jG_{ij}.`$ (3.42)
This procedure accounts for the $`\theta `$ integrations in the preceeding form in (3.27), after choosing the reference momenta such that all $`ϵ_iϵ_j=0`$, $`ϵ_i\overline{ϵ}_j=0`$ and $`\overline{ϵ}_i\overline{ϵ}_j=0`$. The reference momenta that occur naturally in the vertex operator for the $`N=2`$ string in (3.31) force all inner products in (3.41) $`ϵ_iϵ_j=0`$ and $`ϵ_i\overline{ϵ}_j=0`$ via (3.37) and we regain (3.28), although an arbitrary choice of $`q_i`$ demonstrates the covariance in (3.41).
## 4 Three-point Genus One
In this section we calculate the genus-one closed-string three-point amplitude originally derived (the $`M=0`$ part) in and compare the result with field theory, i.e. self-dual gravity. As mentioned in the previous section, in the Leznov frame the tree-level expression $`A_3^{J=2}(\theta =0)`$ from (3.16) exactly produces the field-theory result generated from the Lagrangians (1.5) or (1.7),
$`A_3^{J=2}(\theta =0)=A_3^{\mathrm{tree}}=(ϵ_{\stackrel{\text{.}}{\alpha }\stackrel{\text{.}}{\beta }}k_1^{+\stackrel{\text{.}}{\alpha }}k_2^{+\stackrel{\text{.}}{\beta }})^2`$ (4.1)
where we switched to real spinor notation again. Other formulations of self-dual gravity are related by appropriately normalizing the external lines. In the gauge choice of (1.3) and without the external line factors required for covariance, the field-theoretic one-loop expression $`A_3^{1\mathrm{loop}}`$ is, by dimensional analysis, constrained to be
$`A_3^{1\mathrm{loop}}=(k_1^{+\stackrel{\text{.}}{\alpha }}k_{2\stackrel{\text{.}}{\alpha }}^+)^6\overline{A}_3^{\mathrm{SDG}}.`$ (4.2)
This fixes the tensor structure. The remaining proportionality factor $`\overline{A}_3^{\mathrm{SDG}}`$ in the amplitude then boils down to a field-theoretic triangle integral.
The triangle integrals appearing below are infra-red divergent as on-shell kinematics require $`k_i^2=k_i^+k_i^{}=0`$. The field-theory loop calculation can also be performed by keeping $`k_3^20`$ until after the integration, which generates the infra-red divergence as $`k_3^20`$. Direct comparison with the on-shell string scattering is independent of this limit.
After introducing Feynman parameters and Schwinger time, the three-point on-shell one-loop amplitude becomes
$`\overline{A}_3^{1\mathrm{loop}}`$ $`=`$ $`{\displaystyle \frac{d^d\mathrm{}}{(2\pi )^d}_0^{\mathrm{}}𝑑TT^2_0^1𝑑a_1𝑑a_2𝑑a_3\delta (1a_1a_2a_3)}`$ (4.3)
$`\times \mathrm{exp}\left[T\left(a_1\mathrm{}^2+a_2(\mathrm{}k_1)^2+a_3(\mathrm{}+k_3)^2\right)\right]`$
$`\times (\mathrm{}^{+\stackrel{\text{.}}{\alpha }}k_{1\stackrel{\text{.}}{\alpha }}^+)^2\left((\mathrm{}k_1)^{+\stackrel{\text{.}}{\alpha }}k_{2\stackrel{\text{.}}{\alpha }}^+\right)^2\left(\mathrm{}^{+\stackrel{\text{.}}{\alpha }}k_{3\stackrel{\text{.}}{\alpha }}^+\right)^2`$
which, after shifting
$`\mathrm{}=\mathrm{}^{}+a_2k_1a_3k_3,`$ (4.4)
takes the form of (4.2), with
$`\overline{A}_3^{\mathrm{SDG}}={\displaystyle \frac{d^d\mathrm{}}{(2\pi )^d}_0^{\mathrm{}}𝑑TT^2_0^1𝑑a_1𝑑a_2𝑑a_3a_1^2a_2^2a_3^2\delta (1\underset{j=1}{\overset{3}{}}a_j)\mathrm{exp}\left[T\mathrm{}^2\right]}.`$ (4.5)
Integrating over the loop momentum in (unregulated) $`d=4`$ real dimensions and restoring the tensor structure gives
$`A_3^{1\mathrm{loop}}=(k_1^{+\stackrel{\text{.}}{\alpha }}k_{2\stackrel{\text{.}}{\alpha }}^+)^6\times {\displaystyle \frac{1}{16\pi ^2}}\times {\displaystyle \frac{1}{15\times 5!}}{\displaystyle 𝑑T}.`$ (4.6)
The integral is IR divergent,<sup>10</sup><sup>10</sup>10 It vanishes in dimensional reduction or regularization. and we regulate it by imposing a Schwinger proper-time cutoff at $`T=T_{\mathrm{max}}`$; the unregulated results for both the field theory and string theory may be compared without referring to a regulator.
The three-point function in (4.6) is to be compared with the $`N=2`$ string result found next. The string-theory calculation in the Leznov frame confirms the same tensor structure as in the field theory,
$`A_3^{J=6}(\theta =0)=(k_1^{+\stackrel{\text{.}}{\alpha }}k_{2\stackrel{\text{.}}{\alpha }}^+)^6\overline{A}_3^{N=2}`$ (4.7)
which differs from the three-point scattering found in only by normalization ($`\theta =0`$ instead of $`M=0`$). We may therefore take over their result,
$`\overline{A}_3^{N=2}={\displaystyle \frac{d^2\tau }{\tau _2^2}E_3(\tau ,\overline{\tau })},`$ (4.8)
where the non-holomorphic Eisenstein series is defined as
$`E_3(\tau ,\overline{\tau })={\displaystyle \underset{(m,n)(0,0)}{}}{\displaystyle \frac{\tau _2^3}{|m+n\tau |^6}}`$ (4.9)
and satisfies
$`\tau _2^2_\tau _{\overline{\tau }}E_3(\tau ,\overline{\tau })=6E_3(\tau ,\overline{\tau }).`$ (4.10)
Using (4.10) the integral in (4.8) can be evaluated,
$`\overline{A}_3^{N=2}={\displaystyle d^2\tau (_{\tau _1}^2+_{\tau _2}^2)E_3(\tau ,\overline{\tau })}=\frac{1}{6}_{\tau _2}E_3(\tau ,\overline{\tau })|_{\tau _2=\kappa }`$ (4.11)
where the integral has been regulated by cutting its large-$`\tau _2`$ region at $`\tau _2=\kappa `$. The over the boundary term at the small-$`\tau _2`$ end of the fundamental keyhole domain
$`|\tau |1\mathrm{and}|\tau _1|{\displaystyle \frac{1}{2}}`$ (4.12)
is zero. For $`\tau _2\mathrm{}`$, $`E_3`$ has the asymptotic form
$`E_3=2\zeta (6)\tau _2^3+\sqrt{\pi }\zeta (5)\mathrm{\Gamma }(5/2)\tau _2^2+O(e^{2\pi \tau _2}),`$ (4.13)
yielding for (4.11) the expression
$`\overline{A}_3^{N=2}=\zeta (6)\tau _2^2|_\kappa `$ (4.14)
together with terms that vanish as $`\kappa \mathrm{}`$. Bringing back the tensor structure, we end up with
$`A_3^{J=6}(\theta =0)=(k_1^{+\stackrel{\text{.}}{\alpha }}k_{2\stackrel{\text{.}}{\alpha }}^+)^6\zeta (6)\tau _2^2|_\kappa ,`$ (4.15)
with $`\zeta (6)=\pi ^6/945`$.
The three-point functions in (4.6) and in (4.15) agree after redefining the string proper time $`\tau _2^2=T`$. The regulator
$`T_{\mathrm{max}}=\kappa ^2`$ (4.16)
together with a normalization that can be absorbed in (4.16) gives the match.
The two integrals (4.6) and (4.8) differ by a factor of $`\tau _2`$ or, in the field theory interpretation, a shift in dimension . The matching of the scattering at three-point order is simply a redefinition of the Schwinger proper-time or the cutoff. This is inconsequential at three-point order because the results are both infra-red divergent. However, for the finite higher-point amplitudes such a redefinition is not possible, and the mismatch by a factor of $`\tau _2`$ makes for a crucial difference between the two theories.
## 5 Four-Point Genus One
5.1 String Integrand
In this section we analyze the measure for the integration of the four-point (and higher-point) amplitudes for the $`N=2`$ closed string in the critical dimension $`d=2+2`$ and compute the integrand in terms of the bosonic and fermionic worldsheet correlators. The field-theory limit is taken in order to compare with the self-dual field theory and one-loop maximally helicity violating amplitudes in gravity. The comparison between the measure factors in the string and field theory persists to multi-genus.
The $`N=2`$ superconformal algebra has as its generators the energy-momentum tensor $`T`$, two supercurrents $`G^\pm `$, and the $`U(1)`$ current $`J`$. The associated ghost structure consists of the $`(b,c)`$ diffeomorphism ghosts, the $`(\beta ^{},\gamma ^\pm )`$ local supersymmetry ghosts, and an additional $`(b^{},c^{})`$ ghost system for the local $`U(1)`$ invariance or R symmetry. Each chiral $`N=2`$ matter multiplet $`X=(x,\psi )`$ and each ghost system contributes a (modular invariant) determinant factor to the one-loop string integration measure (continued to $`d`$ dimensional target spacetime)
$`Z_d[\genfrac{}{}{0pt}{}{\alpha }{\beta }](\tau ,\overline{\tau })=Z_x(\tau ,\overline{\tau })Z_\psi [\genfrac{}{}{0pt}{}{\alpha }{\beta }](\tau ,\overline{\tau })Z_{bc}(\tau ,\overline{\tau })Z_{\beta \gamma }[\genfrac{}{}{0pt}{}{\alpha }{\beta }](\tau ,\overline{\tau })Z_{b^{}c^{}}(\tau ,\overline{\tau }),`$ (5.1)
with the respective factors being
$`Z_x(\tau ,\overline{\tau })=\tau _2^{d/2}|\eta (\tau )|^{2d}Z_\psi [\genfrac{}{}{0pt}{}{\alpha }{\beta }](\tau ,\overline{\tau })=|\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](0,\tau )|^d|\eta (\tau )|^d,`$ (5.2)
$`Z_{bc}(\tau ,\overline{\tau })=\tau _2|\eta (\tau )|^4Z_{\beta \gamma }[\genfrac{}{}{0pt}{}{\alpha }{\beta }](\tau ,\overline{\tau })=|\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](0,\tau )|^4|\eta (\tau )|^4`$ (5.3)
and, for the one associated with the local $`U(1)`$ symmetry,
$`Z_{b^{}c^{}}(\tau ,\overline{\tau })=\tau _2|\eta (\tau )|^4.`$ (5.4)
The building blocks are the Jacobi theta functions (featured in the Appendix) with continuous characteristic $`[\genfrac{}{}{0pt}{}{\alpha }{\beta }]`$ equal to spin structure and the Dedekind eta function
$`\eta (\tau )=q^{1/24}{\displaystyle \underset{n0}{}}(1q^n)\mathrm{where}q=e^{2\pi i\tau },`$ (5.5)
with $`\tau `$ denoting the modular parameter of the torus. For general $`d`$ the product of all determinant factors combines into
$`Z_d[\genfrac{}{}{0pt}{}{\alpha }{\beta }](\tau ,\overline{\tau })=\tau _2^{\frac{(d4)}{2}}|\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](0,\tau )|^{d4}|\eta (\tau )|^{3(d4)},`$ (5.6)
and equals unity in four real dimensions . This point trivializes the spin structure summation for the one-loop partition function and signals the absence of a tachyonic mode otherwise arising from the $`q`$-expansion of eta functions.
Superconformal gauge fixing of the worldsheet $`N=(2,2)`$ supergravity produces not only constraints and their ghost systems but also reduces the supergravity path integral to one over the associated finite-dimensional moduli spaces. After explicitly performing the fermionic moduli integrals, which generate picture-raising insertions, one is left with reparametrization and Maxwell moduli. Both come in two varieties: moduli encoding the shape of the $`U(1)`$ bundle over the worldsheet, and moduli describing the locations and $`U(1)`$ monodromies of the vertex operators. In the torus case, the former are $`(\tau ,\overline{\tau })`$ and $`[\genfrac{}{}{0pt}{}{\alpha }{\beta }]`$ while the latter comprise $`\{(z_i,\overline{z}_i)\}`$ and twist angles $`\{(\rho _i,\overline{\rho }_i)\}`$ interpolating between NS- and R-type puncture.<sup>11</sup><sup>11</sup>11 Isometries fix the reparametrization and Maxwell moduli of one of the punctures. Since for genus one the Jacobian torus of spin structures is isomorphic to the worldsheet itself we may parametrize it by an additional torus variable,
$`u=(\frac{1}{2}\alpha )\tau +(\frac{1}{2}\beta ).`$ (5.7)
The modular invariant integration measures are
$`{\displaystyle \frac{d^2\tau }{\tau _2^2}}\mathrm{and}{\displaystyle \frac{d^2u}{\tau _2}}`$ (5.8)
on the fundamental domain $``$ of $`PSL(2,)`$ and the torus $`𝒯`$, respectively. Due to spectral flow, the integrand is independent of the twist angles, whose integration thus results merely in a constant volume factor for each puncture. The integration over the puncture locations, however, are nontrivial but modular invariant in the combination $`_𝒯d^2zV(z,\overline{z})`$.
Putting everything together, the scattering amplitude is given by
$`A_n(k_j)={\displaystyle _{}}{\displaystyle \frac{d^2\tau }{\tau _2^2}}{\displaystyle _𝒯}{\displaystyle \frac{d^2u}{\tau _2}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle _𝒯}d^2z_j{\displaystyle \underset{i<j}{}}e^{k_ik_jG_{ij}}K_{KN}(z_i,\overline{z}_i;u,\overline{u};\tau ,\overline{\tau }),`$ (5.9)
with $`K_{KN}`$ labeling the contractions between the vertex fields.
The expansion of $`K_{KN}`$ has a suggestive form after the grouping of terms that we now turn to. Each term with fermionic contractions can be paired with a purely bosonic term. This property is a consequence of worldsheet $`N=2`$ superconformal invariance and can also be used to prove the vanishing of the corresponding tree-level amplitudes.
We break the contractions into three groups of terms and analyze the contributions from the string scattering when the reference momenta are chosen to agree with (3.31). These holomorphic and anti-holomorphic mirror terms are depicted graphically in Figures 1 and 2. The bosonic propagator is
$`G_{ij}=\mathrm{ln}E(z_iz_j)\mathrm{ln}E(\overline{z}_i\overline{z}_j)+{\displaystyle \frac{2\pi }{\tau _2}}\left[\mathrm{Im}(z_iz_j)\right]^2`$ (5.10)
where $`E`$ is the prime form on the torus,
$`E(z,\tau )={\displaystyle \frac{\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )}{\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](0,\tau )}},`$ (5.11)
and the latter term in (5.10) subtracts the bosonic zero mode from the kernel.<sup>12</sup><sup>12</sup>12 Our convention is that $`G_{ij}`$ marks the full propagator including holomorphic, anti-holomorphic and zero-mode term; the holomorphic piece, $`\mathrm{ln}E(z_iz_j)`$, will be denoted by $`G(z_{ij})`$, explicitly displaying the holomorphic coordinate. The holomorphic half of the fermionic propagator is the Szegö kernel for general continuous monodromies,
$`S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )={\displaystyle \frac{\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](0,\tau )}{\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](0,\tau )\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )}},`$ (5.12)
except for the $`\alpha =\beta =1/2`$ periodic sector in which an additional zero mode develops. Expansions of the propagators are given in the Appendix.
The first type of term in $`K_{KN}`$ is
$`I^{(1234)}`$ $`=`$ $`k_1^+k_2^{}k_2^+k_3^{}k_3^+k_4^{}k_4^+k_1^{}`$ (5.13)
$`\times \left(_1G_{12}_2G_{23}_3G_{34}_4G_{41}S_{12}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]S_{23}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]S_{34}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]S_{41}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]\right).`$
Its reverse ordering $`(4321)`$ is denoted by $`I^{(4321)}`$. The latter gives the complex conjugated contribution via $`k^+k^{}`$. In addition we need the remaining permutations, $`I^{(1324)}`$, $`I^{(1243)}`$, and their reverse orderings. This set is closed under permutation of any two indices.
Next we have the three terms
$`I^{(12)(34)}`$ $`=`$ $`k_1^+k_2^{}k_2^+k_1^{}k_3^+k_4^{}k_4^+k_3^{}`$ (5.14)
$`\times \left(_1G_{12}_2G_{21}_3G_{34}_4G_{43}S_{12}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]S_{21}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]S_{34}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]S_{43}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]\right),`$
together with the orderings $`I^{(14)(23)}`$ and $`I^{(24)(13)}`$. The terms in (5.14) are products of pairs of Szegö kernels as opposed to the cyclic combinations in (5.13).
The remaining terms are paired so that there are products of only two Szegö kernels (in a cyclic fashion),
$`I^{(12)}`$ $`=`$ $`k_1^+k_2^{}k_2^+k_1^{}\left(_1G_{12}_2G_{21}S_{12}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]S_{21}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]\right)`$ (5.15)
$`\times \left(k_3^+k_1^{}_3G_{31}+k_3^+k_2^{}_3G_{32}+k_3^+k_4^{}_3G_{34}\right)`$
$`\times \left(k_4^+k_1^{}_4G_{41}+k_4^+k_2^{}_4G_{42}+k_4^+k_3^{}_4G_{43}\right),`$
together with its permutations: $`I^{(34)}`$, $`I^{(14)}`$, $`I^{(23)}`$, $`I^{(24)}`$, and $`I^{(13)}`$. Terms with three fermion pairs contracted cancel.
The gauge-invariant vertex operators normalized as in (3.41) produce the same set of terms as in (5.13), (5.14) and (5.15) but with the modification
$`k_i^+k_j^{}ϵ_ik_j`$ (5.16)
everywhere. Before and after integrating over spin structures, and with the choice of reference momenta $`q_i=q`$, this substitution shows that the zero-slope limit of the closed-string amplitude reproduces the Feynman rules of gravity one-loop amplitudes without any $`ϵ_iϵ_j`$ or $`ϵ_i\overline{ϵ}_j`$ terms (i.e. MHV structure).
Let us analyze the structure of $`K_{KN}`$ in (5.9) given the boson/fermion pairing of the various terms in the expansion. For the periodic spin structure $`[\genfrac{}{}{0pt}{}{1/2}{1/2}]`$,
$`S_{ij}[\genfrac{}{}{0pt}{}{\alpha }{\beta }]_iG_{ij}`$ (5.17)
for each Szegö kernel, and each set of terms in eqs. (5.13), (5.14) and (5.15) vanishes identically. Furthermore, at generic values of $`[\genfrac{}{}{0pt}{}{\alpha }{\beta }]`$ the integrand vanishes at coincident points $`z_iz_j0`$, making contact with the vanishing tree-level result via worldsheet degeneration. More explicitly, in the short-distance limit of coincident points one gets
$`G(z_iz_j)=\mathrm{ln}(z_iz_j),S(z_iz_j)={\displaystyle \frac{1}{z_iz_j}}=G(z_iz_j),`$ (5.18)
and the integrand is zero pointwise before integration over the vertex operators. This cancellation can be explained in a number of ways. First, in field theory this is due to the fact that every tree diagram in gauge theory (Yang-Mills or gravity) contains at least one contraction $`ϵ_iϵ_j`$, and the identical reference momenta choice for all external lines in an MHV helicity configuration nullifies these terms. Second, at tree-level, target spacetime supersymmetric Ward identities in a supersymmetric gauge theory force the MHV amplitudes to be identically zero (in a supersymmetric gauge theory the tree-level graviton or gauge theory scattering amplitude does not contain internal fermion lines). Third, although the $`N=2`$ string is not spacetime supersymmetric, the worldsheet $`N=2`$ superconformal invariance of the vertex operator forces the tree-level amplitude to be zero.
5.2 Comparison with Field Theory at Zero-slope
In this subsection we take the zero-slope limit of the amplitude obtained in the previous subsection and compare it with the field-theory computation obtained in self-dual gravity at one-loop (2.10). Since the integration over the spin structures may be performed before or after the $`\alpha 0`$ limit and it is not a priori obvious whether the ordering matters (because of singularities at the periodic spin structure), we will examine both orderings: field-theory limit first in the present section, spin structure integration first in the next one. The results will turn out to be the same.
The amplitude from the string differs from self-dual gravity amplitudes in $`d=2+2`$ because of the $`(b^{},c^{})`$ ghost system associated to the $`U(1)`$ R symmetry. Thanks to it, the integrand contains an additional $`\tau _2`$ factor when compared to the integrand of type IIB superstring theory projected onto the self-dual sector of gravity in four dimensions (for example, by toroidal compactification on $`T^6`$ to the non-supersymmetric sector). Quite generally, a factor of <sup>13</sup><sup>13</sup>13 The $`n`$ factors of $`\tau _2`$, one for every vertex operator, arise from the mapping of the torus to the unit square by $`z_i=x_i+\tau y_i`$, with $`x_i,y_i[0,1]`$.
$`{\displaystyle \frac{d\tau _2}{\tau _2}\tau _2^{nd/2}e^{\tau _2f(k_i)}}`$ (5.19)
is associated with writing an $`n`$-point $`\varphi ^3`$ Feynman diagram in $`d`$ dimensions as
$$\frac{d^d\mathrm{}}{(2\pi )^d}\underset{j=1}{\overset{n}{}}\frac{1}{(\mathrm{}p_j)^2}=(4\pi )^{d/2}\left(\underset{j=1}{\overset{n}{}}_0^1𝑑a_j\right)\delta (1\mathrm{\Sigma }_{j=1}^na_j)_0^{\mathrm{}}𝑑TT^{n1d/2}e^{Tf(k_i,a_i)}$$
(5.20)
where $`f(k_i,a_i)=(_kp_ka_k)^2+_kp_k^2a_k`$.
The field-theory limit of the string amplitude is obtained by transforming the string worldsheet coordinates for the vertex operators into a Schwinger proper-time form. From the field-theory point of view, the higher-$`q`$ terms correspond to the exchange of massive modes (which are absent in the $`N=2`$ string). We briefly examine the full analytic structure in the limit. Following in the analytic extraction of poles, we introduce new variables $`w_{ij}`$ satisfying $`|w_{ij}|1`$ and defined by
$`w_{ij}=\{\begin{array}{cc}e^{2\pi iz_{ij}}\hfill & \hfill \mathrm{for}\mathrm{Im}z_{ij}>0\\ qe^{2\pi iz_{ij}}\hfill & \hfill \mathrm{for}\mathrm{Im}z_{ij}<0\end{array}`$ (5.23)
with $`z_{ij}z_iz_j`$. We also make use of the standard parametrization of the vertex insertion points in terms of the real variables $`\alpha _i`$ and $`u_i`$
$`\begin{array}{cc}u_1=y_1\hfill & \alpha _1=2\pi (x_1+u_1\tau _1)\hfill \\ u_2=y_2y_1\hfill & \alpha _2=2\pi (x_2x_1+u_2\tau _1)\hfill \\ u_3=y_3y_2\hfill & \alpha _3=2\pi (x_3x_2+u_3\tau _1)\hfill \\ u_4=1y_3\hfill & \alpha _4=2\pi \tau _1\alpha _1\alpha _2\alpha _3\hfill \end{array},`$ (5.28)
where $`u_1+u_2+u_3+u_4=1`$ and $`\alpha _1+\alpha _2+\alpha _3+\alpha _4=2\pi \tau _1`$. This can be achieved by using the translational symmetry of the torus to fix the position of one vertex operator insertion point.
Since only logarithmic derivatives of the prime form multiply the Koba-Nielsen term, multiplying $`E`$ with a $`z`$-independent factor produces only a constant shift in the Koba-Nielsen exponent which vanishes as a result of momentum conservation. We are therefore entitled to neglect constant factors in $`E`$ and simplify its product representation,
$`E(z_{ij})={\displaystyle \frac{\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z_{ij},\tau )}{\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](0,\tau )}}\dot{=}e^{\pi iz_{ij}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(1q^ne^{2\pi iz_{ij}}\right)\left(1q^{n+1}e^{2\pi iz_{ij}}\right)`$ (5.29)
and we define
$`(w_{ij})={\displaystyle \underset{ij}{}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}|1w_{ij}q^n|^{s_{ij}}`$ (5.30)
which is the component of $`e^{\frac{1}{2}s_{ij}G(z_{ij})}`$ that contains all the infinite products from the expansion in (5.29). The remaining contributions from the Koba-Nielsen terms $`_{i<j}e^{\frac{1}{2}s_{ij}G(z_{ij})}`$ that stem from the $`e^{\pi iz_{ij}}`$ part of the prime forms and from the zero-mode subtractions in the bosonic Green’s functions can be combined into the expression $`|q|^{(su_1u_3+tu_2u_4)}`$.
The full amplitude for a given spin structure $`[\genfrac{}{}{0pt}{}{\alpha }{\beta }]`$ may then be rewritten as $`A_4[\genfrac{}{}{0pt}{}{\alpha }{\beta }](s,t)+A_4[\genfrac{}{}{0pt}{}{\alpha }{\beta }](t,u)+A_4[\genfrac{}{}{0pt}{}{\alpha }{\beta }](u,s)`$, with
$`A_4[\genfrac{}{}{0pt}{}{\alpha }{\beta }](s,t)`$ $`=`$ $`{\displaystyle _{}}d^2\tau \tau _2^2{\displaystyle \underset{i=1}{\overset{4}{}}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\alpha _i}{2\pi }}\delta (2\pi \tau _1\mathrm{\Sigma }_j\alpha _j)`$ (5.32)
$`\times {\displaystyle \underset{i=1}{\overset{4}{}}}{\displaystyle _0^1}du_i\delta (1\mathrm{\Sigma }_ju_j)|q|^{(su_1u_3+tu_2u_4)}(w_{ij})K_{KN}`$
and the obvious permutations. We note that, for a given spin structure, $`K_{KN}`$ is identical to the MHV kinematic factor in IIB superstring theory in (3.41). The function $``$, defined from the product expansion of the $`\vartheta `$-functions as in (5.30), may be expanded in an infinite series as follows:
$`(w_{ij})={\displaystyle \underset{i=1}{\overset{4}{}}}\left|1e^{i\alpha _i}|q|^{u_i}\right|^{s_i}{\displaystyle \underset{n_i=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{|\nu _i|n_i}{}}P_{\{n_i\nu _i\}}^{(4)}(s,t){\displaystyle \underset{i=1}{\overset{4}{}}}|q|^{n_iu_i}e^{i\nu _i\alpha _i}.`$ (5.33)
Here, $`s_i=s`$ for $`i`$ even, $`s_i=t`$ for $`i`$ odd, and $`P_{\{n_i\nu _i\}}^{(4)}(s,t)`$ are polynomials in $`s`$ and $`t`$ that may be generated recursively. Consider now the identity
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\alpha }{2\pi }}e^{i\eta \alpha }\left|1xe^{i\alpha }\right|^s=x^r{\displaystyle _0^{\mathrm{}}}𝑑\beta x^\beta \phi _{r\eta }(s;\beta ),`$ (5.34)
where
$`\phi _{r\eta }(s;\beta )={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}C_k(s)C_{k+|\eta |}(s)\delta (2k+r+|\eta |\beta )`$ (5.35)
is the inverse Laplace transform of a hypergeometric function, and
$`C_k(s)={\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}+k)}{\mathrm{\Gamma }(\frac{s}{2})\mathrm{\Gamma }(k+1)}}.`$ (5.36)
For $`x`$ being some power of $`|q|`$ we can take advantage of this identity and execute the integration over the angular variables $`\alpha _i`$. We observe that the zero-slope limit is equivalent to putting $`(w_{ij})=1`$ from the start, since the effect of a nontrivial function $``$ is felt only at higher order in $`\alpha ^{}`$. The analysis can be extended to contain the angular parameters associated with the kinematical factor $`K_{KN}`$ and justifies the substitution rules that follow from the closed-string context. The remaining propagator terms in $`K_{KN}`$ generate the Feynman parameters associated with the derivative couplings of the field-theory vertex in the zero-slope limit; the kinematical expression is identical to that obtained in the IIB superstring before summing over spin structures. This procedure has been systematized at one-loop and $`n`$-point through the Bern-Kosower string-motivated rules for calculating gauge theory scattering adapted to gravity .
If it were not for the non-holomorphic zero-mode part in (5.10), perfect bose-fermi cancellation in eqs. (5.13), (5.14), and (5.15) would occur in the field-theory limit. As it is, however, the remainder of the various pairings of the bosonic contractions with the fermionic ones is proportional to at least one factor of $`(2\pi /i\tau _2)\mathrm{Im}z_{ij}`$. The zero modes explicitly break the holomorphicity of the string scattering, and the MHV amplitude may be understood as a holomorphic anomaly in the zero-slope limit of the $`N=2`$ string.
Remaining in (5.32) is the single factor of $`\tau _2^2`$ multiplied by bosonic zero-mode contributions from $`K_{KN}`$. Comparing with (5.19) we see that this corresponds to a field-theory result in $`d=2`$ real dimensions . The $`\tau _2`$-dependence in the low-energy scattering of gravitons and the $`d=2`$ technical interpretation follows from either simple toroidal compactifications of the IIB superstring, or using the well-established string-inspired Feynman rules adapted to the case of perturbative gravity . The latter we briefly discuss next in order to map the kinematical structure $`K_{KN}`$ to the MHV one-loop gravity amplitudes.
The string-inspired generation of the graviton scattering amplitude, which in the zero-slope limit arises from the corners of the moduli space, involves the Feynman-parametrized form originating from a $`\varphi ^3`$ diagram
$`D=c_n{\displaystyle _0^1}𝑑x_{i_{n1}}\mathrm{}{\displaystyle _0^{x_{i_2}}}𝑑x_{i_1}{\displaystyle \frac{K_{\mathrm{red}}}{\left(_{a<b}^nP_{i_a}P_{j_b}x_{i_aj_b}(1x_{i_aj_b})\right)^{nd/2}}},`$ (5.37)
where, in $`d`$ dimensions,
$`c_n=\left(4\pi \right)^{2d/2}{\displaystyle \frac{\mathrm{\Gamma }(nd/2)}{16\pi ^2}}.`$ (5.38)
$`P_i`$ is the momentum flowing into the $`i^{\mathrm{th}}`$ leg of the $`n`$-gon $`\varphi ^3`$ diagram, $`x_{ij}=x_ix_j`$, and $`x_i`$ are Feynman parameters. All $`\varphi ^3`$ diagrams are to be considered with external trees attached where the external lines follow a cycle ordering; in the gravitational case we sum over all of the non-cyclic orderings without associated Yang-Mills color factors. The factor $`K_{\mathrm{red}}`$ comprises terms generated from the kinematical expression identical to the Koba-Nielsen term of the multi-graviton scattering amplitude,
$`\mathrm{exp}\left[(k_iϵ_jϵ_jk_i)\dot{G}_{ij}ϵ_iϵ_j\ddot{G}_{ij}\right]\times (\mathrm{anti}\mathrm{hol})|_{\mathrm{multi}\mathrm{linear}}.`$ (5.39)
In (5.39) the dotted $`G`$s represent worldline derivatives of the complete propagator, including bosonic zero modes.<sup>14</sup><sup>14</sup>14 This notation follows that of the first-quantized form of scattering amplitudes at one-loop, derived and motivated by string theory considerations. In (5.39) we also have a multiplicative factor of the mixing from holomorphic/anti-holomorphic,
$`\mathrm{exp}\left[(ϵ_i\overline{ϵ}_j+\overline{ϵ}_iϵ_j)H_{ij}\right],`$ (5.40)
where $`H_{ij}:=_i\overline{}_jG_{ij}`$ in the field-theory expression for the amplitude. This comes from the last term in (3.41) and is zero for the $`N=2`$ string because of the MHV-type condition that $`ϵ_i\overline{ϵ}_j=0`$ for all external legs $`i,j`$. The propagator in (5.39) is $`\dot{G}_{ij}=\frac{1}{2}\mathrm{sign}(x_{ij})+x_{ij}`$, and the usual Feynman parameters are related to $`x_i`$ via $`x_i=_{j=1}^ia_j`$. The point of the form as written in (5.37) is that the kinematical expression arises from the standard first-quantized form of a particle, as generated from integrating over the worldline with measure factor (5.19). Mapping the $`N=2`$ zero-slope limit to this expression removes the need to explicitly integrate over $`\tau _2`$ (including four-dimensional box integrals with up to eight insertions of loop momenta in the numerator) because the amplitudes are known . The Feynman-parametrized form of (5.37) and (5.39) produces the integral form of the zero-slope limit of the $`N=2`$ string expression but with $`d=2`$ as opposed to $`d=4`$, as we shall now demonstrate.
We begin by writing all possible $`\varphi ^3`$ diagrams, obtained by pinching together different sets of vertex operators. Then, after expanding the kinematical factor in (5.39) we collect sets of $`\dot{G}_{ij}`$ (and $`\dot{\overline{G}}_{ij}`$) terms in accord with the tree and loop rules (example diagrams are illustrated in Figure 3). The Bern-Kosower tree rule in the low-energy extraction involves substituting on an external leg $`1/(P_i+P_j)^2`$ for the occurence of a single power of $`\dot{G}_{ij}`$, from the outside of the diagram into the diagram, and then resubstituting $`i=j`$ in the remaining momentum flow of the tree-line as well as making the substitution in the remaining $`\dot{G}_{ij}`$ (and $`\dot{\overline{G}}_{ij}`$) factors. In the gravity analog we substitute the same when there is a single product of $`\dot{G}_{ij}\dot{\overline{G}}_{ij}`$. In the string-theory amplitude, this amounts to pinching a pole from the $`e^{k_ik_jG_{ij}}`$ kinematical factor (i.e. integrating near $`z_iz_j`$). After substituting the tree rules on an individual diagram we have a remaining kinematical factor on which we apply the loop rules. The tree rules do not depend on the spacetime spin of the particle being integrated out; rather, the loop rules map to the internal spacetime statistics of the particle.
The Bern-Kosower loop rules tell us to expand the holomorphic and anti-holomorphic terms in (5.39) for a given ordering of the four external lines and, after applying the tree rules, to substitute the factors of $`\dot{G}_{ij}`$ by the low-energy expansion of the propagators with an overall factor of two for combinatorics,
$`\dot{G}_{ij}{\displaystyle \frac{1}{2}}\mathrm{sign}(x_{ij})+x_{ij}.`$ (5.41)
These generate the uncyclic contributions and represent the bosonic portion of the worldsheet correlators in the zero-slope limit. Next, we examine the integrand for cyclic occurances of $`\dot{G}_{ij}`$, following the ordering of the legs exiting the loop (for example $`\dot{G}_{12}\dot{G}_{23}\dot{G}_{31}`$), and substitute as follows,
$`\dot{G}_{ij}\dot{G}_{ji}2\mathrm{and}\dot{G}_{i_1i_2}\dot{G}_{i_2i_3}\mathrm{}\dot{G}_{i_ni_1}1(n>2).`$ (5.42)
After having applied the cyclic substitutions in (5.42) some $`\dot{G}_{ij}`$ may be left unsubstituted; they are to be replaced according to (5.41). The outcome is the cyclic contributions and model the zero-slope limit of the fermionic correlations. The loop rules are applied separately on the $`\dot{G}_{ij}`$s and $`\dot{\overline{G}}_{ij}`$s in the case of gravity.
In the case of Einstein-Hilbert gravity, the rules in (5.41) and (5.42) generate a graviton as the state within the loop. In the field-theory limit, the cyclic and uncyclic substitutions arise from the fermionic and bosonic worldsheet propagators, respectively (when there are no $`\ddot{G}_{ij}`$ terms), and match with the form of $`\mathrm{K}_{KN}`$. (A worldline systematics in the case of spin $`[J1]`$ has also been analyzed in a number of works, including within the context of 1PI diagrams.)
There are further simplifications in the integrand of the MHV amplitudes that are beyond the naive collecting of terms obtained from expanding the Koba-Nielsen factor. The fact that integrating out a spacetime graviton or a spacetime complex scalar in the loop makes no difference for the amplitude implies cancellations of the cyclic terms obtained from the second rule in (5.42). In order to obtain the MHV amplitude with an internal complex scalar, only the first loop rule (5.41) needs to be implemented on both the holomorphic and anti-holomorphic sides.
In the MHV amplitudes, which satisfy (2.8), this means that integrating out the cyclic contributions in (5.42) gives identically zero ($`A^{[2]}=A^{[0]}`$). This fact has been noted in in the application to a gravitational four-point amplitude with helicity assignment $`(,++,++,++)`$. At the level of the superstring and the $`N=2`$ string this means that the worldsheet fermions do not contribute to the MHV amplitudes in the field-theory limit. Momentum conservation eliminates their total sum; this is demonstrated in the next section. We shall find an identical result when integrating over the spin structures prior to taking $`\alpha ^{}0`$.
The $`q`$-expansion of the derivative of the bosonic propagator involved in extracting $`K_{\mathrm{KN}}`$ is (see the Appendix)
$`\left[\mathrm{ln}|E(z)|^2{\displaystyle \frac{2\pi }{\tau _2}}(\mathrm{Im}z)^2\right]=i\pi {\displaystyle \frac{2\pi i}{1e^{i\alpha }|q|^u}}+{\displaystyle \frac{2i\pi }{\tau _2}}\mathrm{Im}z+𝒪(q),`$ (5.43)
where $`\mathrm{Im}z0`$. Effectively, after the angular integration , the middle term in (5.43) integrates to $`2\pi i`$ with higher-order in $`\alpha ^{}`$ effects (naively there appears to be a potential singularity), yielding the outcome
$`\left[\mathrm{ln}|E(z)|^2{\displaystyle \frac{2\pi }{\tau _2}}(\mathrm{Im}z)^2\right]i\pi +{\displaystyle \frac{2i\pi }{\tau _2}}\mathrm{Im}z=i\pi (12y),`$ (5.44)
in agreement with the rules in (5.41) and (5.42) and the tensor algebra of the one-loop diagram after angular integration.<sup>15</sup><sup>15</sup>15 In the field theory limit $`\mathrm{cot}2\pi zi`$.
Turning to the fermions, one observes that the field-theory limit of a Szegö kernel is independent of the spin structure $`[\genfrac{}{}{0pt}{}{\alpha }{\beta }]`$. In the $`q`$-expansion,
$`S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z_{i_1i_2},\tau )S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z_{i_2i_3},\tau )\mathrm{}S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z_{i_ni_1},\tau )(i\pi )^n+O(q)`$ (5.45)
after the integration over the angular coordinates $`\alpha _i`$. With this substitution, the second rule in (5.42) obtains for the integral expression of the amplitude. The zero-slope limit of the $`N=2`$ string reproduces individually all the diagrams of the gravity amplitude after a careful tracking of the indices of the $`\dot{G}_{ij}`$ which come in a specific order within the Koba-Nielsen form in (5.39) ($`\dot{G}_{ij}=\dot{G}_{ji}`$).
The primary difference between the integrands of the $`N=2`$ string and the IIB superstring truncated to obtain four-dimensional gravity lays in the integration measure. Concretely, the $`N=2`$ string scattering amplitude at $`n`$-point has an extra factor of $`\tau _2`$ compared to the amplitude obtained from a field-theory calculation using the Feynman rules of the self-dual gauge theory . As a result, the amplitude in (2.10) is obtained effectively in $`d=2`$ and not in $`d=4`$.
We compare now with the IIB superstring measure, continued to $`D`$ dimensions. On a torus with spin structure $`[\genfrac{}{}{0pt}{}{\alpha }{\beta }]`$, the NSR fermions and the supersymmetry ghosts $`(\beta ,\gamma )`$ produce the determinantal factors
$`Z_\psi [\genfrac{}{}{0pt}{}{\alpha }{\beta }]=|\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](0,\tau )|^D|\eta (\tau )|^DZ_{\beta \gamma }[\genfrac{}{}{0pt}{}{\alpha }{\beta }]=|\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](0,\tau )|^2|\eta (\tau )|^2`$ (5.46)
while the bosonic coordinates and the reparametrization ghosts $`(b,c)`$ yield
$`Z_x=\tau _2^{D/2}|\eta (\tau )|^{2D}Z_{bc}=\tau _2|\eta (\tau )|^4.`$ (5.47)
Together with the Weyl-Peterson measure $`d^2\tau /\tau _2^2`$, the product $`Z_\psi Z_{\beta \gamma }Z_xZ_{bc}`$ yields
$`{\displaystyle \frac{d^2\tau }{\tau _2^2}}\tau _2^{(D2)/2}|\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](0,\tau )|^{D2}|\eta (\tau )|^{3(D2)},`$ (5.48)
not taking into account the factors associated with the vertex operators. Upon compactification on $`T^{Dd}`$, it is modified by a lattice sum,
$`Z(\mathrm{\Gamma })=\tau _2^{(Dd)/2}{\displaystyle \underset{(P_L,P_R)\mathrm{\Gamma }}{}}e^{i\pi \tau P_LP_Li\pi \overline{\tau }P_RP_R}`$ (5.49)
where $`(P_L,P_R)`$ parametrize the $`(p,q)`$ signature lattice of dimension $`Dd`$ (consistency requires $`P_L^2P_R^22`$ and $`pq8`$). Furthermore, the individual vertex operators generate powers of $`\tau _2`$ after evaluating $`d^2z`$ (the volume $`d^2z=\tau _2`$).
The two measures, (5.48) times (5.49) for the compactified IIB superstring on one side and (5.6) times (5.8) for the $`N=2`$ string on the other, differ at zero-slope by a single factor of $`\tau _2`$ :
$`{\displaystyle \frac{d^2\tau }{\tau _2^2}}\tau _2^{(d2)/2}{\displaystyle \frac{d^2\tau }{\tau _2^2}}\tau _2^{(d4)/2}.`$ (5.50)
This calculation indicates the dimensional shift interpretation of the field-theory integration: the IIB superstring compactified on $`T^6`$ involves a $`d^2\tau /\tau _2^{3n}`$ at $`n`$-point (after inserting $`D=10`$ and $`d=4`$ in (5.48) and (5.49)). This is the same factor that the bosonic string in $`d=26`$ compactified on $`T^{22}`$ generates. In contrast, the $`N=2`$ string requires a $`d^2\tau /\tau _2^{2n}`$.
## 6 Spin Structure Summation
6.1 Torus Integrals of Elliptic Functions
This section pushes the expression for the full $`N=2`$ string scattering amplitude a step further and also provides an alternative calculation of its field-theory limit. Concretely, we explicitly evaluate the integrals over the monodromies of the worldsheet fermions, before taking the field theory limit. At four-point order this involves integrating over spin structures various products of up to four holomorphic Szegö kernels (those in (5.13) and (5.14)) together with the anti-holomorphic side with the measure in (5.8).
For a complex structure $`\tau `$ of the torus, we first define (suppressing $`\tau `$ dependence)
$`h_n(\{z_{ij}\};u):=S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z_{12})S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z_{23})\mathrm{}S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z_{n1}),`$ (6.1)
with
$`u=(\frac{1}{2}\alpha )\tau +(\frac{1}{2}\beta )`$ (6.2)
denoting the spin-structure dependent zero locus of the Szegö kernel. By inspecting the zeros and poles of (6.1) we learn how to rewrite this expression in terms of prime forms,
$`h_n={\displaystyle \frac{E(z_{12}u)E(z_{23}u)\mathrm{}E(z_{n1}u)}{[E(u)]^nE(z_{12})E(z_{23})\mathrm{}E(z_{n1})}},`$ (6.3)
which exposes the single $`n`$th-order pole in $`u`$ at the origin. The simplest case, $`n=2`$, yields
$`h_2(z_{12})={\displaystyle \frac{E(z_{12}u)^2}{E(u)^2E(z_{12})^2}}=\mathrm{}(z_{12})+\mathrm{}(u)=^2\mathrm{ln}E(z_{12})^2\mathrm{ln}E(u)`$ (6.4)
where $`\mathrm{}(z)`$ is the Weierstraß elliptic function. Furthermore in the coincidence limit $`z_{n1}0`$ one observes that $`h_nh_{n1}/z_{n1}`$. Since the spin structure has been encoded in an additional torus variable $`u`$, we have to integrate over $`u`$ (with correct measure) the functions $`h_n`$ times their anti-holomorphic relatives.
With $`h_0:=1`$ we define the integrals
$`f_{n,\overline{n}}(\{z_{ij},\overline{z}_{ij}\}):=h_n(\{z_{ij}\};u)\overline{h}_{\overline{n}}(\{\overline{z}_{kl}\};\overline{u})`$ (6.5)
with measure
$`\mathrm{}:={\displaystyle \frac{dud\overline{u}}{2i\tau _2}\mathrm{}}`$ (6.6)
which normalizes $`1=1`$. Explicit integration for (6.5) is made possible by the following theorem . As $`h_n(u)du`$ and $`\overline{h}_{\overline{n}}(\overline{u})d\overline{u}`$ are both closed one-forms with zero residue at $`u=0`$ we can express the surface integral in terms of period integrals over the $`a`$ and $`b`$ cycle,
$`f_{n,\overline{n}}={\displaystyle \frac{i}{2\tau _2}}\left[{\displaystyle _a}h_n(u)𝑑u{\displaystyle _b}\overline{h}_{\overline{n}}(\overline{u})𝑑\overline{u}{\displaystyle _b}h_n(u)𝑑u{\displaystyle _a}\overline{h}_{\overline{n}}(\overline{u})𝑑\overline{u}\right].`$ (6.7)
We next evaluate the period integrals.
It is a fact that an elliptic function with a single $`n`$th-order pole can be expressed as a linear combination of $`\mathrm{}`$ and its derivatives plus a constant. Hence, expanding $`h_n(u)`$ around the pole (no residue!) we obtain
$`h_n(u)`$ $`=`$ $`h_n^{(n)}u^n+\mathrm{}+h_n^{(3)}u^3+h_n^{(2)}u^2+h_n^{(0)}+𝒪(u)`$ (6.8)
$`=`$ $`\frac{()^n}{(n1)!}h_n^{(n)}\mathrm{}^{(n2)}(u)+\mathrm{}\frac{1}{2}h_n^{(3)}\mathrm{}^{}(u)+h_n^{(2)}\mathrm{}(u)+H_n^{(0)}`$
with Laurent coefficients $`h_n^{(k)}(\{z_{ij}\})`$, where we used $`\mathrm{}(u)=u^2+O(u^2)`$ and
$`\frac{()^k}{(k1)!}\mathrm{}^{(k2)}(u)=u^k+G_k\delta _{k\mathrm{even}}+O(u)`$ (6.9)
for $`k3`$. The holomorphic Eisenstein series
$`G_k={\displaystyle \underset{(m,n)(0,0)}{}}{\displaystyle \frac{1}{(m\tau +n)^k}}=2\zeta (k)+O(e^{2\pi i\tau })`$ (6.10)
occuring in (6.9) for even $`k`$ lead to a shift of the constant term in (6.8),
$`h_n^{(0)}H_n^{(0)}=h_n^{(0)}G_4h_n^{(4)}G_6h_n^{(6)}\mathrm{}G_{2[\frac{n}{2}]}h_n^{(2[\frac{n}{2}])}.`$ (6.11)
The virtue of the expression (6.8) is that the evaluation of its period integrals has become almost trivial. Indeed, since for $`k3`$ the antiderivative of $`\mathrm{}^{(k2)}(u)`$ is $`\mathrm{}^{(k3)}(u)`$, a doubly-periodic function, the integral of $`\mathrm{}^{(k2)}(u)`$ over a closed loop vanishes. This observation eliminates all period integrals except for the last two terms in (6.8). Hence, we only require the integrals
$`{\displaystyle _a}𝑑u=1{\displaystyle _b}𝑑u=\tau `$ (6.12)
as well as
$`{\displaystyle _a}𝑑u\mathrm{}(u)=2\eta _1=G_2{\displaystyle _b}𝑑u\mathrm{}(u)=2\eta _\tau =2\pi iG_2\tau `$ (6.13)
together with their complex conjugates, where we have introduced the “almost-modular” form (the regulated form of the divergent sum in (6.10) for $`k=2`$)
$`G_2(\tau )=4{\displaystyle \underset{(m,n)(0,0)}{}}{\displaystyle \frac{1}{(m\tau +n)^2(2m\tau +2n1)}}.`$ (6.14)
Via (6.5) and (6.7) this leaves us with only three basic non-vanishing spin structure averages,
$`1=1,\mathrm{}=G_2+{\displaystyle \frac{\pi }{\tau _2}},\mathrm{}\overline{\mathrm{}}=G_2\overline{G}_2{\displaystyle \frac{\pi }{\tau _2}}(G_2+\overline{G}_2).`$ (6.15)
With these averages we can now compute the integrals (6.5) as
$`f_{n,\overline{n}}`$ $`=`$ $`H_n^{(0)}\overline{H}_{\overline{n}}^{(0)}+H_n^{(0)}\overline{h}_{\overline{n}}^{(2)}\overline{\mathrm{}}+\overline{H}_{\overline{n}}^{(0)}h_n^{(2)}\mathrm{}+h_n^{(2)}\overline{h}_{\overline{n}}^{(2)}\mathrm{}\overline{\mathrm{}}`$
$`=`$ $`\left[H_n^{(0)}+h_n^{(2)}(G_2+\frac{\pi }{\tau _2})\right]\left[\overline{H}_{\overline{n}}^{(0)}+\overline{h}_{\overline{n}}^{(2)}(\overline{G}_2+\frac{\pi }{\tau _2})\right]h_n^{(2)}\overline{h}_{\overline{n}}^{(2)}(\frac{\pi }{\tau _2})^2`$
and observe that they do not split chirally. In the evaluation of the four-point function we require only the cases of $`(n,\overline{n})\{(0,0),(2,0),(2,2),(4,0),(4,2),(4,4)\}`$ together with the transposes.
It remains to list the coefficients $`h_n^{(k)}`$. In general $`h_n^{(1)}=0`$ and $`h_n^{(n)}=()^n`$. Here, we only need $`h_2^{(k)}`$ and $`h_4^{(k)}`$ for even $`k`$,
$`H_2^{(0)}`$ $`=`$ $`h_2^{(0)}=\mathrm{}(z_{12})=^2\mathrm{ln}E(z_{12})+G_2,`$ (6.17)
$`h_4^{(2)}`$ $`=`$ $`\frac{1}{2}T_4^1\stackrel{~}{}^2T_4+2G_2,`$ (6.18)
$`H_4^{(0)}`$ $`=`$ $`h_4^{(0)}G_4=\frac{1}{24}T_4^1\stackrel{~}{}^4T_4+G_2T_4^1\stackrel{~}{}^2T_4+2G_2^2,`$ (6.19)
The shorthand notation involving $`T_4`$ (generalizable to higher $`n`$ in this form) is
$`T_4`$ $`=`$ $`E_{12}E_{23}E_{34}E_{41},`$ (6.20)
$`\stackrel{~}{}^2T_4`$ $`=`$ $`E_{12}^{\prime \prime }E_{23}E_{34}E_{41}+E_{12}E_{23}^{\prime \prime }E_{34}E_{41}+E_{12}E_{23}E_{34}^{\prime \prime }E_{41}+E_{12}E_{23}E_{34}E_{41}^{\prime \prime }`$ (6.21)
$`+2E_{12}^{}E_{23}^{}E_{34}E_{41}+2E_{12}^{}E_{23}E_{34}^{}E_{41}+2E_{12}^{}E_{23}E_{34}E_{41}^{}`$
$`+2E_{12}E_{23}^{}E_{34}^{}E_{41}+2E_{12}E_{23}^{}E_{34}E_{41}^{}+2E_{12}E_{23}E_{34}^{}E_{41}^{}`$
and similarly for $`\stackrel{~}{}^4T_4`$, where we abbreviated $`E_{ij}=E(z_{ij})`$. The higher $`\stackrel{~}{}^k`$ represents the actions of $`k`$ derivatives spread out with respect to the insertion points $`z_{ij}`$.
For later reference, we present the first few spin structure integrals:
$`f_{2,0}`$ $`=`$ $`^2\mathrm{ln}E_{12}+\frac{\pi }{\tau _2}`$ (6.22)
$`f_{2,2}`$ $`=`$ $`|^2\mathrm{ln}E_{12}+\frac{\pi }{\tau _2}|^2(\frac{\pi }{\tau _2})^2`$ (6.23)
$`f_{4,0}`$ $`=`$ $`\frac{1}{24}T_4^1\stackrel{~}{}^4T_4+(G_2+\frac{\pi }{\tau _2})\frac{1}{2}T_4^1\stackrel{~}{}^2T_4+2G_2\frac{\pi }{\tau _2}.`$ (6.24)
The analysis above is generalizable to higher genus by employing the prime forms pertaining to the higher-genus Riemann surface.
6.2 Zero-slope Limit
We now analyze the field theory limit of the various terms obtained from summing over the spin structures. In the process of evaluating the ratios of prime forms $`E(z_{ij})`$ and their derivatives, the following can be implemented:
$`E(z){\displaystyle \frac{1}{2\pi }}\mathrm{sin}(2\pi z),`$ (6.25)
which in $`f_{n,\overline{n}}`$ leads to products of unity and $`\mathrm{cot}(2\pi z)`$. After the angular integration over $`z_r`$, $`\mathrm{cot}(2\pi z)i`$. Therefore, the analysis of the fermionic correlator terms $`f_{n,\overline{n}}`$ reduces to combinatoric factors and derivatives of (6.25) with respect to the $`z`$-coordinates, together with the appropriate $`\varphi ^3`$ diagram via pinching the $`|E(z_{ij})|^{\alpha ^{}s_{ij}}`$.
We consider first $`f_{2,0}(z)`$ and $`f_{2,2}(z,\overline{z})`$ given by (6.22) and (6.23), respectively,
$`f_{2,0}(z)(2\pi )^2\left[1+\mathrm{cot}^2(2\pi z)\right]+{\displaystyle \frac{\pi }{\tau _2}}{\displaystyle \frac{\pi }{\tau _2}},`$ (6.26)
$`f_{2,2}(z,\overline{z})\left|(2\pi )^2\mathrm{cot}^2(2\pi z)(2\pi )^2+{\displaystyle \frac{\pi }{\tau _2}}\right|^2\left({\displaystyle \frac{\pi }{\tau _2}}\right)^20,`$ (6.27)
where $`\overline{z}`$ is a variable independent of $`z`$. All of the contributions containing $`f_{2,0}`$ vanish after adding them up in the expansion of the Koba-Nielsen factor; we will show this after analyzing the remaining terms.
The remaining integrals $`f_{n,\overline{n}}`$ all involve at least either $`n=4`$ or $`\overline{n}=4`$. The field theory limit of such a term, exemplified in (6.24), does not vanish individually in the kinematical expression, but like terms add up to zero as we will show below. In the $`\tau _2\mathrm{}`$ limit the Eisenstein series simplify,
$`G_2{\displaystyle \frac{\pi ^2}{3}}\mathrm{and}G_4{\displaystyle \frac{\pi ^4}{45}}.`$ (6.28)
As displayed in (6.21) the $`\stackrel{~}{}^4`$ and $`\stackrel{~}{}^2`$ derivatives produce a large number of products of $`E_{ij}^{(k)}/E_{ij}`$. However, in the field theory limit no $`z_{ij}`$ dependence survives and
$`{\displaystyle \frac{E_{ij}^{(k)}}{E_{ij}}}(2\pi i)^k`$ (6.29)
which yields
$`T_4^1\stackrel{~}{}^4T_4+256(2\pi )^4T_4^1\stackrel{~}{}^2T_416(2\pi )^2.`$ (6.30)
Collecting all terms, the net limits of the remaining terms are
$`f_{4,0}`$ $``$ $`\stackrel{~}{f}_{4,0}=10(2\pi )^4{\displaystyle \frac{47}{6}}(2\pi )^2{\displaystyle \frac{\pi }{\tau _2}}`$ (6.31)
$`f_{4,2}`$ $``$ $`\stackrel{~}{f}_{4,2}=10(2\pi )^4{\displaystyle \frac{\pi }{\tau _2}}`$ (6.32)
$`f_{4,4}`$ $``$ $`\stackrel{~}{f}_{4,4}=100(2\pi )^8{\displaystyle \frac{470}{3}}(2\pi )^6{\displaystyle \frac{\pi }{\tau _2}}.`$ (6.33)
It is interesting that a $`1/\tau _2`$ appears in these terms, which indicates the modification necessary to obtain the MHV amplitudes. These terms terms do not produce a $`0/0`$ effect as they also vanish in four dimensions, being proportional to the difference between a scalar contribution and a graviton contribution to the MHV amplitude.
The spin structure integrals $`f_{n,\overline{n}}`$ are being multiplied by kinematical factors $`t_{n,\overline{n}}(\{ϵ_i,k_j\})`$ stemming from the contractions of polarization and momentum vectors on four vertex operators (3.39). Since in the field-theory limit all $`z`$ dependence has dropped from $`\stackrel{~}{f}_{n,\overline{n}}`$, the latter can be factored out from the remaining integrations. This fact allows one to combine directly the various permutations of a given kinematical factor.
We now analyze the kinematical factors $`t_{n,\overline{n}}=t_n\overline{t}_{\overline{n}}`$. We begin with the contractions of four pairs of fermions, which can happen in two distinct ways. The first possibility is a single cycle connecting all pairs,
$`t_4^{(1234)}=ϵ_1k_2ϵ_2k_3ϵ_3k_4ϵ_4k_1,`$ (6.34)
together with permutations $`(12)`$ and $`(23)`$. With arbitrary reference momenta $`q`$ chosen the same for all polarization vectors, we have
$`t_4^{(1234)}=[12][23][34][41]`$ (6.35)
which follows from the substitution
$`ϵ^{\alpha \dot{\alpha }}(k,q)=i{\displaystyle \frac{q^\alpha k^{\dot{\alpha }}}{q^\beta k_\beta }}.`$ (6.36)
Adding the three permutations produces
$`t_4^{(1234)}+t_4^{(2134)}+t_4^{(1324)}=[12][23][34][41]+[21][13][34][42]+[13][32][24][41].`$ (6.37)
Employing twice the Fierz identity
$`[AB][CD]=[AC][BD]+[AD][CB]`$ (6.38)
we find
$`t_4^{(1234)}+t_4^{(2134)}+t_4^{(1324)}`$ $``$ $`[12][23][34][41]+[13][24]\left([12][34]+[14][23]\right)`$ (6.39)
$`=`$ $`[12][23][34][41]+[13]^2[24]^2`$
$`=`$ $`[12][34]\left([24][31]+[21][43]\right)+[13]^2[42]^2`$
$`=`$ $`[12][24][43][31]+[12]^2[34]^2+[13]^2[42]^2.`$
Symmetrizing both sides of this equation and using the identity
$`[13]^2[42]^2+[12]^2[43]^2+[14]^2[32]^2=0`$ (6.40)
(again from (6.38)) one discovers that $`t_4^{(1234)}+t_4^{(2134)}+t_4^{(1324)}`$ equals minus itself. Thus, the sum in (6.37) vanishes.
The second option for contracting four pairs of fermions produces two cycles of two pairs each,
$`t_4^{(12)(34)}=ϵ_1k_2ϵ_2k_1ϵ_3k_4ϵ_4k_3,`$ (6.41)
together with its two permutations. Via (6.36) this equals
$`t_4^{(12)(34)}=[12]^2[34]^2,`$ (6.42)
which upon adding the three permutations and using (6.40) equals zero, too. Similar additions of the field theory limit for the fermionic terms add to zero in the three-point and two-point $`\varphi ^3`$ diagrams, where the momentum structure involves three and two independent momenta, respectively.
The previous analysis regarding the cyclic terms proportional to $`f_{4,0}`$ generalizes in a straightforward manner to the remaining cyclic terms multiplying $`f_{4,2}`$ and $`f_{4,4}`$. In the field-theory limit all of the $`z`$ dependence on the holomorphic half of the kinematical expression multiplying these functions is absent (the anti-holomorphic half multiplying $`f_{4,2}`$ includes bosonic zero modes which translate to Feynman parameters in the field theory limit). After summing over the different contributions on the holomorphic half, these contributions equal zero, as was shown in the preceeding paragraphs.
Finally we analyze the $`f_{2,0}`$ (and $`f_{2,2}`$) terms. The Wick contractions of two pairs of fermions yield a kinematical factor of
$`t_2^{(ij)}=ϵ_ik_jϵ_jk_i,`$ (6.43)
which multiplies the remaining kinematical structure from $`_kG_k\mathrm{}`$, as displayed by the solid lines in Figure 2. The relation $`A^{[0]}=A^{[1]}=A^{[2]}`$ enforces all these terms to be zero. This “supersymmetry identity” implies that the fermionic contractions associated with rule two all generate zero, as discussed in the previous section. In the string amplitude this means that the holomorphic sum of all of the world-sheet fermionic correlators with equal weight add to zero. We have already showed by momentum conservation (for the particular MHV helicity structure) that the four-fermion terms are zero, it follows that the $`t_2^{(ij)}`$ terms also add to zero (as they all have the same coefficient in (6.26)). This cancellation occurs separately for the four-fermion-pair contractions, i.e. the $`t^4`$ terms, as well as the two-fermion-pair contractions, in both two and four dimensions (as the supersymmetry idendity holds in both cases). The factors of $`\tau _2`$ in (6.26) and (6.31)–(6.33) cause a dimensional shift in the integration.
## 7 Discussion
In this work we have analyzed several aspects of the quantum scattering of the closed $`N=(2,2)`$ closed superstring at genus one. First, we have derived the zero-slope limit of the one-loop four-point function in the RNS formulation, and explicitly integrated over the spin structures of the worldsheet fermions. We have found agreement with the existing vanishing theorems in the literature. The mapping of the genus-one moduli space integrand to an MHV amplitude at $`n`$-point order is performed. Second, we have compared the one-loop integrated three- and higher-point string amplitudes with those of self-dual gravity. The disagreement (vanishing versus nonzero MHV) could be traced to a known difference in the integration measure whose origin is the local R symmetry of the $`N=2`$ string. Third, we have made manifest the Lorentz and coordinate invariance of the quantum (and classical) scattering by normalizing the vertices and incorporating spinor helicity techniques. Most of this analysis carries over straightforwardly to the open string.
A number of new features have arisen regarding the quantum amplitudes. The $`N=2`$ string has field equations of self-duality at the classical level, but at genus one its amplitudes are not directly found from self-dual field theory in four dimensions. Rather, the result appears in the loop integration as the dimensionally regulated version of the self-dual amplitudes, continued to two dimensions (with external kinematics in two complex dimensions). These field-theory amplitudes indeed vanish.
The two-dimensional nature of the $`N=2`$ string loop integration suggests that the effective dynamics of this string is only (real) two-dimensional. Then, the vanishing of two-dimensional gravity (and Yang-Mills) amplitudes may account for the all-order vanishing of the string amplitudes (like at genus one). Clearly, a string in two-dimensional target spacetime has no room for physical excitations (at generic momenta). In four-dimensional spacetime, however, the same situation can be arrived at by increasing the worldvolume dimension from two to four, since a spacetime-filling brane affords only topological degrees of freedom. Indeed, the analogy
$`TJ(b,c)(b^{},c^{})`$ (7.1)
and the fact that the $`N=2`$ string ghost systems remove two complex unphysical directions from the excitation spectrum (best seen for the RNS fermions) have led to the speculation that the $`N=2`$ string actually is a space-filling brane. It is tempting to interpret the $`U(1)^2`$ fibre associated with the local R symmetry as carrying the two additional dimensions, making for a total of four parametrizing the full bundle. The $`N=2`$ string formulation then amounts to a fibration of the $`2+2`$ dimensional world-volume over a Riemann surface.
Another avenue is to search for modifications in the string amplitude which resurrect the non-vanishing four-dimensional MHV scattering. A single factor of $`1/\tau _2`$ in the integrand of the closed string is required to extract the one-loop self-dual field-theory amplitudes in $`d=4`$. An insertion of an unintegrated zero-momentum vertex operator or a bosonic zero mode would already do the job, for example, through
$`\underset{k=0}{lim}\sqrt{g}x\overline{}xe^{ikx}\mathrm{or}\overline{}G(z,\overline{z})={\displaystyle \frac{2\pi }{\tau _2}}\delta ^{(2)}(z,\overline{z}).`$ (7.2)
This conformal anomaly may have a target spacetime interpretation as a $`\beta `$ function expansion around $`d=2`$. The insertion of the zero mode, breaking the worldsheet conformal invariance, expands the amplitude to those of one-loop self-dual field theory, and through perturbations of self-duality to gravity and Yang-Mills theory.
We have analyzed one-loop string amplitudes in the field-theory limit, i.e. calculated the leading term in the $`q`$-expansion of a string amplitude. As the vanishing theorems and the Ward identities of the $`N=2`$ string imply that the entire tower of $`q`$-expansion coefficients is zero, we expect the above two-dimensional interpretation to hold for the full $`N=2`$ string theory. A direct verification of the vanishing of the higher-$`q`$ components of the genus one string amplitude is outside the scope of the present paper; however, we have made important steps in that direction by providing the reader with an explicit expression of the string integrand after spin-structure summation. Whether the analysis remains feasible at $`\alpha ^{}0`$ at the level of the full $`q`$-expansion is to be shown.
Finally, the one-loop amplitudes generated by the closed $`N=2`$ string are related through an order $`ϵ=10d`$ identity to those of IIB supergravity in ten dimensions. Via a relation $`A_{N=2}ϵA_{IIB}`$, the zero-slope limit of the $`N=2`$ string captures the ultraviolet portion of the IIB amplitudes; the latter amplitudes are finite in a dimensionally regulated form in ten dimensions. As both amplitudes are low-energy limits of critical string theories, this suggests a relation between the $`N=2`$ and $`N=1`$ strings, which at multi-loop requires a similar relation between the MHV amplitudes and the non-MHV IIB amplitudes. It is interesting to note that membrane-string and string-string connections have been noted in the context of the heterotic $`(2,1)`$ formulation in relation to world-volumes of membranes .
Acknowledgements
The work of G.C. is supported in part by the U.S. Department of Energy, Division of High Energy Physics, Contract W-31-109-ENG-38. Support for O.L. and B.N. by the German National Science Foundation (DFG) under grant LE 838/5-2 is gratefully acknowledged. G.C. thanks the Institute for Theoretical Physics at the University of Hannover for hospitality and Warren Siegel for relevant correspondence. O.L. acknowledges fruitful discussions with Michael Flohr, Klaus Hulek, Jacob Nielsen and Jeroen Spandaw. B.N. is thankful to Klaus Jünemann for useful conversations.
## 8 Appendix : Theta Functions
We list in this Appendix some of the properties of the Jacobi theta functions and elliptic functions useful in this work. The theta function with $`(\alpha ,\beta )`$ characteristics is defined by the infinite sum
$$\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )=\underset{nZ}{}e^{\pi i\tau (n+\alpha )^2+2\pi i(n+\alpha )(z+\beta )}.$$
(8.1)
The theta function satisfies the identity
$`\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )`$ $`=`$ $`e^{\pi i\tau \alpha ^2+2\pi i\alpha (z+\beta )}\vartheta [\genfrac{}{}{0pt}{}{0}{0}](z+\tau \alpha +\beta ,\tau )`$ (8.2)
$`=`$ $`e^{\pi i\tau (\alpha ^21/2)^2+2\pi i(\alpha 1/2)(z+\beta )}\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z+(\alpha \frac{1}{2})\tau +(\beta \frac{1}{2}),\tau )`$
with the $`\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )=\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )`$. Abbreviating $`q=e^{2\pi i\tau }`$, the infinite product form of the odd theta function reads
$$\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )=iq^{1/8}e^{\pi iz}\underset{n=1}{\overset{\mathrm{}}{}}\left(1q^n\right)\underset{n=0}{\overset{\mathrm{}}{}}\left[\left(1q^ne^{2\pi iz}\right)\left(1q^{n+1}e^{2\pi iz}\right)\right],$$
(8.3)
and that of its $`z`$-derivative at $`z=0`$ is
$$\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](0,\tau )=2\pi q^{1/8}\underset{n=1}{\overset{\mathrm{}}{}}\left(1q^n\right)^3.$$
(8.4)
In the same manner, we may rewrite the prime form as
$$E(z,\tau )=\frac{\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )}{\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](0,\tau )}=\frac{e^{\pi iz}}{2\pi i}\frac{\underset{n=0}{\overset{\mathrm{}}{}}\left[\left(1q^ne^{2\pi iz}\right)\left(1q^{n+1}e^{2\pi iz}\right)\right]}{_{n=1}^{\mathrm{}}\left(1q^n\right)^2},$$
(8.5)
where inspection reveals that
$`E(z+1,\tau )=E(z,\tau )\mathrm{and}E(z+\tau ,\tau )=e^{\pi i\tau 2\pi iz}E(z,\tau ).`$ (8.6)
The chiral bosonic correlator (without the zero-mode part) is
$`G(z,\tau )=\mathrm{ln}E(z,\tau )=\mathrm{ln}{\displaystyle \frac{\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )}{\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](0,\tau )}}.`$ (8.7)
Once we insert the bosonic propagators into the expression for the four-point function, the $`z`$-independent factors will vanish as a result of momentum conservation. We also need the expanded version of $`G`$, where we define the parameter $`w=e^{2\pi iz}`$,
$`G(z,\tau )`$ $`=`$ $`\pi i{\displaystyle \frac{1+w^1}{1w^1}}+2\pi i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}q^n\left({\displaystyle \frac{w}{1q^nw}}{\displaystyle \frac{w^1}{1q^nw^1}}\right)`$ (8.8)
$`=`$ $`\pi \mathrm{cot}(\pi z)+2\pi i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}q^n\left({\displaystyle \frac{w}{1q^nw}}{\displaystyle \frac{w^1}{1q^nw^1}}\right).`$ (8.9)
The fermionic Szegö kernel for a given spin structure $`(\alpha ,\beta )(\frac{1}{2},\frac{1}{2})`$ is
$`S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )`$ $`=`$ $`{\displaystyle \frac{\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](0,\tau )}{\vartheta [\genfrac{}{}{0pt}{}{\alpha }{\beta }](0,\tau )\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )}}`$ (8.10)
$`=`$ $`e^{2\pi i(\alpha 1/2)z}{\displaystyle \frac{\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z+(\alpha \frac{1}{2})\tau +(\beta \frac{1}{2}),\tau )\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](0,\tau )}{\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}]((\alpha \frac{1}{2})\tau +(\beta \frac{1}{2}),\tau )\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )}}.`$
For the odd spin structure, the fermionic propagator (again ignoring the zero-mode part) is a derivative of the chiral bosonic Greens function,
$`S[\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )=G(z,\tau )=\mathrm{ln}E(z,\tau )={\displaystyle \frac{\vartheta ^{}[\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )}{\vartheta [\genfrac{}{}{0pt}{}{1/2}{1/2}](z,\tau )}}.`$ (8.11)
This relation between bosonic and fermionic propagator extends to the anti-holomorphic part and the zero-mode part as well,
$`{\displaystyle \frac{2\pi }{i\tau _2}}\mathrm{Im}z=_z{\displaystyle \frac{2\pi }{\tau _2}}\left[\mathrm{Im}z\right]^2.`$ (8.12)
We must integrate over all spin structures including the odd one; however, the odd spin-structure correlators do not contribute to any of the amplitudes derived in this work.
The Szegö kernels $`S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )`$ as well as $`G(z,\tau )`$ are singular as we take $`z0`$; however, the combination
$`[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )=S[\genfrac{}{}{0pt}{}{\alpha }{\beta }](z,\tau )G(z,\tau )`$ (8.13)
is finite in this limit.
The Weierstraß function $`\mathrm{}(z,\tau )`$ is the unique doubly periodic function with a single second-order pole at the origin and no constant term in its Laurent expansion,
$`\mathrm{}(z)`$ $`=`$ $`{\displaystyle \frac{1}{z^2}}+{\displaystyle \underset{(m,n)(0,0)}{}}\left({\displaystyle \frac{1}{(zm\tau n)^2}}{\displaystyle \frac{1}{z^2}}\right)`$ (8.14)
$`=`$ $`{\displaystyle \frac{1}{z^2}}+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(2k+1)G_{2k+2}(\tau )z^{2k},`$
with the modular form
$`G_{2k+2}(\tau )={\displaystyle \underset{(m,n)(0,0)}{}}(m\tau +n)^{(2k+2)}=2\zeta (2k+2)+O(q)`$ (8.15)
being known as the holomorphic Eisenstein function of weight $`2k+2`$, for $`k1`$. The antiderivative of the $`\mathrm{}`$ function is denoted a $`\zeta (z)`$; it takes the half-point values
$`\zeta (1/2)\eta _1=\frac{1}{2}G_2\mathrm{and}\zeta (\tau /2)\eta _\tau =\frac{1}{2}G_2\tau i\pi `$ (8.16)
where the failure of the sum in (8.15) to absolutely converge for $`k=0`$ necessitates a regularized definition of the “almost-modular” form
$`G_2(\tau )=4{\displaystyle \underset{(m,n)(0,0)}{}}{\displaystyle \frac{1}{(m\tau +n)^2(2m\tau +2n1)}}={\displaystyle \frac{4\pi }{i}}_\tau \mathrm{ln}\eta (\tau ).`$ (8.17)
The last equality makes contact with the logarithm of the Dedekind eta function,
$`\mathrm{ln}\eta (\tau )={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{ln}(1e^{2\pi in\tau })+i\pi \frac{\tau }{12},`$ (8.18)
and generalizes to the higher Eisenstein functions, e.g.
$`2G_2(\tau )^210G_4(\tau )=\left({\displaystyle \frac{4\pi }{i}}\right)^2_\tau ^2\mathrm{ln}\eta (\tau ).`$ (8.19)
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# FIRST ORDER PHASE TRANSITION OF THE Q-STATE POTTS MODEL IN TWO DIMENSIONS
## 1 Introduction
In many of the physical systems that exhibit the first order phase transition, the order of the transition changes to the second order by changing the parameter of the system. It is important to know how the quantities at the first order phase transition point diverge when the parameter approaches the point at which the order of the transition changes. The $`q`$-state Potts model$`^{\mathrm{?},\mathrm{?}}`$ in two dimensions gives a good place to investigate this subject. It exhibits the first order phase transition for $`q>4`$ and the second order one for $`q4`$. Many quantities of this model are known exactly at the phase transition point $`\beta =\beta _t`$ for $`q>4`$, including the latent heat$`^\mathrm{?}`$ and the correlation length in the disordered phase,$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ which vanishes and diverges, respectively, in the limit $`q4_+`$. One the other hand other physically important quantities such as the specific heat, the magnetic susceptibility, and the correlation length (in the ordered phase) at the transition point, which also diverge as $`q4_+`$, are not solved exactly.
Here we calculate the large-$`q`$ expansion series of the energy cumulants including the specific heat and the magnetization cumulants including the magnetic susceptibility in both the phases and the correlation length in the ordered phase at the transition point using the finite lattice method.$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ Obtained long series for the energy and magnetization cumulants give the estimates of the quantities that are more precise by a factor of $`10^210^4`$ than the Monte Carlo simulations. Especially its estimates are within an accuracy of $`0.1\%`$ at $`q=5`$, where the correlation length is as large as a few thousands of the lattice spacing. Bhattacharya et al.$`^\mathrm{?}`$ made a stimulating conjecture on the asymptotic behavior of the energy cumulants at the first order transition point; the relation between the energy cumulants and the correlation length in their asymptotic behavior at the first order transition point for $`q4_+`$ will be the same as their relation in the second order phase transition for $`q=4`$ and $`\beta \beta _t`$, which is well known from their critical exponents. The obtained series enables us to confirm the correctness of the conjecture.
As for the correlation length at the first order phase transition point, the results of the Monte Carlo simulation$`^\mathrm{?}`$ and the density matrix renormalization group$`^\mathrm{?}`$ indicate that the correlation lengths are very close to each other in the ordered and disordered phases for $`q>10`$. On the other hand, at the second order phase transition point($`q4`$) their ratio is known to be $`1/2`$. It is interesting whether the ratio is exactly equal to unity, remains close to unity, or approaches $`1/2`$ when $`q4_+`$. To investigate it we calculate the first few terms of the large-$`q`$ expansion for the eigenvalues of the transfer matrix and find that from the second largest to the $`N`$-th largest eigenvalues with $`N`$ the one-dimensional size of the lattice make a continuum spectrum in the thermodynamic limit both in the ordered and disordered phases. We also calculate the long series of the second moment correlation length in both the phases, which serve to investigate the behavior of the spectrum of the eigenvalues of the transfer matrix in the region of $`q`$ close to $`4`$.
## 2 Finite lattice method
Here we use the finite lattice method,$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ to generate the large-$`q`$ series for the Potts model, instead of the standard diagrammatic method used by Bhattacharya et al.$`^\mathrm{?}`$ The finite lattice method can in general give longer series than those generated by the diagrammatic method especially in lower spatial dimensions. In the diagrammatic method, one has to list up all the relevant diagrams and count the number they appear. In the finite lattice method we can skip this job and reduce the main work to the calculation of the expansion of the partition function for a series of finite size lattices, which can be done using the straightforward site-by-site integration$`^{\mathrm{?},\mathrm{?}}`$ without the diagrammatic technique. This method has been used mainly to generate the low- and high-temperature series in statistical systems and the strong coupling series in lattice gauge theory. We note that this method is applicable to the series expansion with respect to any parameter other than the temperature or the coupling constant. Using this method we calculated$`^{\mathrm{?},\mathrm{?}}`$ the series for the $`n`$-th energy cumulants ($`n=06`$) to $`z^{23}`$, $`n`$-th magnetization cumulants ($`n=13`$) to $`z^{21}`$, and second moment correlation length to $`z^{19}`$ with $`z1/\sqrt{q}`$.
## 3 Energy cumulants
The latent heat $``$ at the transition point are known to vanish at $`q4_+`$ as
$$3\pi x^{1/2},$$
with $`x=\mathrm{exp}(\pi ^2/2\theta )`$ and $`2\mathrm{cosh}\theta =\sqrt{q}`$. Bhattacharya et al.’s conjecture says that the $`n`$-th energy cumulants $`F_{d,o}^{(n)}`$ at the first order transition point $`\beta =\beta _t`$ will diverge for $`q4_+`$ as
$$F_d^{(n)},(1)^nF_o^{(n)}\alpha B^{n2}\frac{\mathrm{\Gamma }\left(n\frac{4}{3}\right)}{\mathrm{\Gamma }\left(\frac{2}{3}\right)}x^{3n/22}.$$
(1)
The constants $`\alpha `$ and $`B`$ in Eq.(1) should be common to the ordered and disordered phases from the duality relation for each $`n`$-th cumulants.
If this conjecture is true, the product $`F^{(n)}^{3n4}`$ is a smooth function of $`\theta `$, so we can expect that the Padé approximation of $`F^{(n)}^p`$ will give convergent result at $`p=3n4`$. It has been examined for the large-$`q`$ series obtained by the finite lattice methods for $`n=2,\mathrm{},6`$ both in the ordered and disordered phases, which in fact give quite convergent Padé approximants for $`p=3n4`$ and as $`p`$ leaves from this value the convergence of the approximants becomes bad rapidly. An example can be seen in Fig. 1 for n=2 in the disordered phase.
We give in Table 1 the values of the specific heat $`C=\beta _t^2F^{(2)}`$ evaluated from these Padé approximants for some values of $`q`$. These estimates are three or four orders of magnitude more precise than (and consistent with) the previous result for $`q7`$ from the large-$`q`$ expansion to order $`z^{10}`$ by Bhattacharya et al.$`^\mathrm{?}`$ and the result of the Monte Carlo simulations for $`q=10,15,20`$ carefully done by Janke and Kappler$`^\mathrm{?}`$ as in Table 2. What should be emphasized is that we obtained the values of the specific heat in the accuracy of about 0.1 percent at $`q=5`$ where the correlation length is as large as 2500. As for the asymptotic behavior of $`F^{(n)}`$ at $`q4_+`$, the Padé data of $`F_d^{(2)}/x`$ and $`F_o^{(2)}/x`$ have the errors of a few percent around $`q=4`$ and their behaviors are enough to convince us that the conjecture (1) is true for $`n=2`$ with $`\alpha =0.073\pm 0.002.`$ Furthermore from the conjecture (1) the combination $`\left\{\mathrm{\Gamma }\left(n\frac{4}{3}\right)|F^{(n)}|/\mathrm{\Gamma }\left(\frac{2}{3}\right)F^{(2)}\right\}^{\frac{1}{n2}}x^{\frac{3}{2}}`$ is expected to approach the constant $`B`$ for each $`n(3)`$, and in fact the Padé data for every $`n(=3,\mathrm{},6)`$ gives $`B=0.38\pm 0.05,`$ which also gives strong support to the conjecture for $`n3`$.
## 4 Magnetization cumulants
The behavior of the $`n`$-th magnetization cumulants $`M^{(n)}`$ for $`\beta \beta _t`$ at $`q=4`$ is well known as $`M_{d,o}^{(n)}A_{d,o}^{(n)}(\xi )^{\frac{15}{8}n2}`$ and parallel to the conjecture for the energy cumulants by Bhattacharya et al. we can make a conjecture that
$$M_{d,o}^{(n)}\mu _{d,o}^{(n)}x^{\frac{15}{8}n2}.$$
(2)
in the limit $`q4_+`$ with $`\beta =\beta _t`$. We have examined the Padé approximation of $`M^{(n)}^p`$ for the large-$`q`$ series generated by the finite lattice method for $`n=2`$ and $`3`$ both in the ordered and disordered phases, which in fact gives quite convergent Padé approximants for $`p=15n/44`$ and as $`p`$ leaves from this value the convergence of the approximants becomes bad rapidly again. In Table 3 we present the resulting estimates of the magnetic susceptibility $`\chi _{d,o}=M_{d,o}^{(2)}`$. Our result is much more precise than the Monte Carlo simulation$`^\mathrm{?}`$ at least by a factor of 100 as in Table 4. From the behavior of $`M_{d,o}^{(n)}/x^{(15n/82)}`$ we obtain the coefficients in Eq.(2) as $`\mu _d^{(2)}=0.0020(2)`$, $`\mu _o^{(2)}=0.0016(1)`$ and $`\mu _d^{(3)}=7.4(5)\times 10^5`$, $`\mu _o^{(3)}=7.9(2)\times 10^5`$. These convince us that the conjecture made for the magnetization cumulants is also true.
## 5 Exponential correlation length
Next we investigate the exponential correlation length $`\xi _{1,o}`$ in the ordered phase. Here the exponential correlation length is defined by $`\xi _1=\mathrm{log}(\mathrm{\Lambda }_1/\mathrm{\Lambda }_0)`$ with $`\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_1`$ the largest and the second largest eigenvalues of the transfer matrix, respectively. (We have omitted the subscript ’$`1`$’ in the previous sections) There is an obstacle to extract the correction to the leading term of the large-$`q`$ expansion for the correlation length at the phase transition point from the correlation function $`<𝒪(t)𝒪(0)>_c`$, since we know from the graphical expansion that it behaves like
$`<𝒪(t)𝒪(0)>_c`$ $``$ $`z^t(1+2zt^2+\mathrm{})`$
$``$ $`\mathrm{exp}(mt).`$
for a large distance $`t`$. Thus we will diagonalize the transfer matrix $`T`$ directly for large-$`q`$. The eigenfunction for the largest eigenvalue $`\mathrm{\Lambda }_0`$ in the leading order of $`z`$ is
$$|0>|\underset{N}{\underset{}{\mathrm{0\; 0\; 0}\mathrm{}\mathrm{..0\; 0\; 0}}}>$$
where all of the $`N`$ spin variables (each of which can take the value of $`s=0,1,\mathrm{},q1`$) are fixed to be zero, with the element of the transfer matrix $`<0|T|0>`$ $`=1+O(z^2)`$ and the corresponding eigenvalue is $`\mathrm{\Lambda }_0=1+O(z)`$. The eigenfunctions for the second largest eigenvalue are
$$|I>\frac{1}{\sqrt{NI+1}}|0\mathrm{}0\underset{I}{\underset{}{eee\mathrm{}ee}}0\mathrm{}0>$$
$$(I=1,\mathrm{},N),$$
$$|e>\frac{1}{\sqrt{q1}}\underset{s=1}{\overset{q1}{}}|s>$$
with the diagonal matrix elements
$$<I|T|I>=z+O(z^2),$$
and the off-diagonal matrix elements starting from higher orders in $`z`$. The second largest eigenvalues of $`T`$ degenerate in the leading order with $`\mathrm{\Lambda }_i=z+O(z^{3/2})`$. The degeneracy of the eigenvalues of the first $`N`$ ’excited states’ is the reason why the expansion series (5) of the correlation function cannot be exponentiated into a single exponential term. The off-diagonal matrix elements resolve the degeneracy with
$`\mathrm{\Lambda }_1/\mathrm{\Lambda }_0`$ $`=`$ $`z+4z^{3/2}+6z^2+O(z^3),`$
$`\mathrm{}`$
$`\mathrm{\Lambda }_N/\mathrm{\Lambda }_0`$ $`=`$ $`z4z^{3/2}+6z^2+O(z^3).`$
for $`N\mathrm{}`$. These eigenvalues constitute a continuum spectrum. It appears to be kept in any higher order of $`z`$. From Eq.(5) we obtain
$$1/\xi _{1,o}=\mathrm{log}z4z^{1/2}+2z+8/3z^{3/2}+O(z^2).$$
This is the same as the large-$`q`$ expansion of $`1/\xi _{1,d}`$ given by Buffenoir and Wallon$`^\mathrm{?}`$ to this order.
In the disordered phase the situation is quite similar. The eigenvalues of the transfer matrix for the first $`N`$ excited states constitute a continuum spectrum with their values exactly the same as in the ordered phase at least to the order of $`z^{5/2}`$.
The eigenvalues form the continuum spectrum just on the first order phase transition point. Off the transition point, we can see that the spectrum is discrete with $`\mathrm{\Lambda }_i\mathrm{\Lambda }_{i+1}O(ϵ)`$ for $`\sqrt{z}ϵ1`$ and $`\mathrm{\Lambda }_i\mathrm{\Lambda }_{i+1}O(ϵ^{2/3})`$ for $`ϵ\sqrt{z}1`$ where $`ϵ\beta /\beta _t1`$.
## 6 Second moment correlation length
Here we give the results of the large-$`q`$ expansion of the second moment correlation lengths $`\xi _{2nd,o}`$ in the ordered phase and $`\xi _{2nd,d}`$ in the disordered phase defined by
$$\xi _{2nd}^2=\frac{\mu _2}{2d\mu _0},$$
where $`\mu _2`$ and $`\mu _0`$ are the second moment of the correlation function and the magnetic susceptibility, respectively. The obtained expansion coefficients$`^\mathrm{?}`$ for the ordered and disordered phases coincide with each other to order $`z^3`$ and differ from each other in higher orders. The ratio of the second moment correlation length in the ordered phase to that in the disordered phase is estimated by the Padé analysis to be not far from unity even in the limit of $`q4`$ ($`\xi _{2nd,o}/\xi _{2nd,d}=0.930(3)`$). Another point is that the ratio $`\xi _{2nd,d}/\xi _{1,d}`$ of the second moment correlation length to the exponential correlation length is much less than unity in the region of $`q`$ where the correlation length is large enough. It approaches $`0.51(2)`$ for $`q4`$. It is known that in the limit of the large correlation length,
$$\frac{\xi _{2nd}^2}{\xi _1^2}\frac{_{i=1}^{\mathrm{}}c_i^2(\xi _i/\xi _1)^3}{_{i=1}^{\mathrm{}}c_i^2(\xi _i/\xi _1)}<1,$$
with $`\xi _i\mathrm{log}(\mathrm{\Lambda }_i/\mathrm{\Lambda }_0)`$. If the ’higher excited states’ ($`i=2,3,\mathrm{}`$) did not contribute so much, this ratio would be close to unity, as in the case of the Ising model on the simple cubic lattice, where $`\xi _{2nd}/\xi _1=0.970(5)`$ at the critical point.$`^{\mathrm{?},\mathrm{?}}`$ Our result implies that the contribution of the ’higher excited states’ is large in the disordered phase of the Potts model in two dimensions even when $`q`$ is close to $`4`$. This strongly suggests that the eigenvalues of the transfer matrix for the first $`N`$ excited states in the disordered phase form the continuum spectrum not only in the large-$`q`$ region but also when $`q`$ approaches $`4`$.
As for the exponential correlation length $`\xi _{1,o}`$ in the ordered phase for $`q4`$, it is difficult to calculate the eigenvalues of the transfer matrix in much higher orders. It is quite natural, however, to expect that the ratio $`\xi _{1,o}/\xi _{1,d}`$ would be close to unity even in the limit of $`q4`$. The reason is the following. If the ratio $`\xi _{1,o}/\xi _{1,d}`$ would be $`1/2`$ in the limit of $`q4`$, which is the known ratio in the second order phase transition point($`q4`$), then the ratio $`\xi _{2nd,o}/\xi _{1,o}`$ should be close to unity, which would imply that the higher excited states would not contribute so much to $`\xi _{2nd,o}`$ and the eigenvalue of the transfer matrix for the first excited state would be separated from the higher excited states. This scenario is not plausible, since as already mentioned in section 5 the continuum spectrum of the eigenvalues of the transfer matrix appears to be maintained in any high order in $`z`$. In this case, we can expect that the ratio $`\xi _{2nd,o}/\xi _{1,o}`$ would be around $`1/2`$ as is the case in the disordered phase, resulting that $`\xi _{1,o}/\xi _{1,d}`$ is close to unity.
## 7 Summary
We generated the large-$`q`$ series for the energy and magnetization cumulants at the first order phase transition point of the two-dimensional $`q`$-state Potts model in high orders using the finite lattice method. They gave very precise estimates of the cumulants for $`q>4`$ and confirmed the correctness of the Bhattacharya et al.’s conjecture that the relation between the cumulants and the correlation length for $`q=4`$ and $`\beta \beta _t`$ (the second order phase transition) is kept in their asymptotic behavior for $`q4_+`$ at $`\beta =\beta _t`$ (the first order transition point). If this kind of relation is satisfied as the asymptotic behavior for the quantities at the first order phase transition point in more general systems when the parameter of the system is varied to make the system close to the second order phase transition point, it would serve as a good guide in investigating the system.
Further the large-$`q`$ expansion of the eigenvalues of the transfer matrix was calculated in the first 4 terms. We found that they have the same spectra of the eigenvalues of the transfer matrix in the ordered and disordered phases giving the same exponential correlation length ($`\xi _{1,o}=\xi _{1,d}`$) to the order of $`z^{3/2}`$ and that the spectra are continuous in the thermodynamic limit. We also calculated the large-$`q`$ expansion of the second moment correlation length in the ordered and disordered phases in high orders and found that they differ from each other in higher orders than $`z^3`$, but that the ratio $`\xi _{2nd,d}/\xi _{1,d}`$ is not far from unity for all region of $`q>4`$. We also found that $`\xi _{2nd,d}/\xi _{d,1}`$ is far from unity even in the limit of $`q4`$. It receives significant contributions not only from the ’first excited state’ but also ’higher excited states’ and this suggest strongly that the continuum spectrum would be maintained (i.e. there would be no particle state) in the disordered phase. From these results it is quite natural to expect that the exponential correlation length $`\xi _{1,o}`$ in the ordered phase is not far from that in the disordered phase even in the limit of $`q4`$ and it is not plausible that their ratio approaches $`1/2`$ that is their ratio in the second order phase transition point($`q4`$).
## References
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# Self-Similar Hot Accretion Flow onto a Neutron Star
## 1 Introduction
At mass accretion rates less than a few per cent of the Eddington rate, black holes (BHs) are believed to accrete via a hot, two-temperature, radiatively inefficient, quasi-spherical, advection-dominated accretion flow, or ADAF (Ichimaru 1977; Narayan & Yi 1994, 1995a, 1995b; Abramowicz et al. 1995; Chen et al. 1995). The physical properties of ADAFs around BHs have been investigated by a number of authors, and detailed spectral models have been applied to observations of BH candidates see Narayan, Mahadevan & Quataert 1998 and Kato, Fukue & Mineshige 1998 for reviews.
At low mass accretion rates, accretion onto a neutron star (NS) is also expected to occur via a hot, two-temperature flow (Narayan & Yi 1995b; Yi et al. 1996), but the properties of such flows have not been investigated. NS flows are expected to differ from BH ADAFs in several respects. (i) Whereas in the case of a BH the accreting material flows freely and supersonically through the absorbing boundary at the event horizon, in the case of a NS the radial velocity of the material must decelerate to zero. (ii) The accreting material is expected to apply a spin-up or spin-down torque on the NS. Popham & Narayan (1991) and Paczyński (1991) investigated the nature of the torque for a cold thin disk, but the case of a hot flow has not been studied. (iii) For similar mass accretion rates (in Eddington units), the luminosity of a NS accreting via an ADAF is likely to be much higher than that of a BH because a NS has a surface while a BH has an event horizon (Narayan & Yi 1995b; Narayan, Garcia & McClintock 1997; Menou et al. 1999). (iv) The spectra are expected to be different.
We discuss in this paper the structure of a hot accretion flow around a NS, or any other relativistic star with a surface. The flow under consideration is global and extends radially to very large distances (at least thousands of NS radii or more) where it matches onto appropriate outer boundary conditions. We do not attempt a detailed analysis of the boundary layer region near the NS, where the accretion flow meets the star. We present only an approximate analysis of this region, which extends at most a few NS radii above the stellar surface for a more detailed discussion of the physics of the boundary layer see Popham & Narayan 1991; Paczyński 1991; Titarchuk, Lapidus, & Muslimov 1998; Titarchuk & Osherovich 1999.
The paper is organized as follows. We show in §2 that there is a radially extended region around the NS where the flow “settles” with a radial velocity much less than the local free-fall velocity. We obtain a self-similar solution for this settling region and show that, surprisingly, the density and temperature of this zone are independent of the mass accretion rate. In §3, we discuss physical properties of the accretion flow. In §4 we compare the analytical results with numerical computations and in §5 we discuss the relationship between the settling flow and an ADAF. We conclude with a discussion in §6.
## 2 Self-Similar Settling Solution
We consider a steady, rotating, axisymmetric, quasi-spherical, two-temperature accretion flow onto a star with a surface. We use the height-integrated form of the viscous hydrodynamic equations (Ichimaru 1977; Abramowicz et al. 1988; Paczyński 1991; Narayan & Yi 1994):
$`\dot{M}=4\pi R^2\rho v,`$ (1)
$`v{\displaystyle \frac{dv}{dR}}=\left(\mathrm{\Omega }^2\mathrm{\Omega }_K^2\right)R{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{d}{dR}}\left(\rho c_s^2\right),`$ (2)
$`4\pi \alpha {\displaystyle \frac{\rho c_s^2R^4}{\mathrm{\Omega }_K}}{\displaystyle \frac{d\mathrm{\Omega }}{dR}}=\dot{J}\dot{M}\mathrm{\Omega }R^2,`$ (3)
$`\rho vT_p{\displaystyle \frac{ds_p}{dR}}={\displaystyle \frac{\rho vc^2}{(\gamma _p1)}}{\displaystyle \frac{d\theta _p}{dR}}vc^2\theta _p{\displaystyle \frac{d\rho }{dR}}=(1\delta )q^+q_{\mathrm{Coul}},`$ (4)
$`\rho _evT_e{\displaystyle \frac{ds_e}{dR}}={\displaystyle \frac{\rho _evc^2}{(\gamma _e1)}}{\displaystyle \frac{d\theta _e}{dR}}vc^2\theta _e{\displaystyle \frac{d\rho _e}{dR}}=\delta q^++q_{\mathrm{Coul}}q^{},`$ (5)
where $`\dot{M}`$ is the mass accretion rate, $`R`$ is the spherical radius, $`\rho `$ is the mass density of the accreting gas, $`v`$ is the radial infall velocity, $`\mathrm{\Omega }`$ is the angular velocity, $`\mathrm{\Omega }_K(R)=(GM/R^3)^{1/2}`$ is the Keplerian angular velocity, $`c_s^2=c^2(\theta _p+\theta _em_e/m_p)`$ is the square of the isothermal sound speed, $`T_{p,e}`$ are the temperatures of protons and electrons, $`\theta _{p,e}=k_BT_{p,e}/m_{p,e}c^2`$ are the corresponding dimensionless temperatures, $`\alpha `$ is the Shakura-Sunyaev viscosity parameter, $`\dot{J}`$ is the rate of accretion of angular momentum, $`s_p`$ and $`s_e`$ are the specific entropies of the proton and electron fluids, $`\rho _e(m_e/m_p)\rho `$ is the mass density of the electron fluid, $`\gamma _p`$ and $`\gamma _e`$ are the adiabatic indices of protons and electrons (which, in general, may be functions of $`T_p`$ and $`T_e`$), and $`q^+`$, $`q^{}`$, and $`q_{\mathrm{Coul}}`$ are the viscous heating rate, radiative cooling rate, and energy transfer rate from protons to electrons via Coulomb collisions, per unit mass. We have assumed that a fraction $`\delta `$ of the viscous heat goes into electrons and a fraction $`1\delta `$ into protons; it is usually assumed in ADAF models that $`\delta 1`$, but our analysis is general for any value of $`\delta `$ between 0 and 1. Equations (1)–(5) describe the conservation of mass, radial momentum, angular momentum, proton energy and electron energy, respectively.
For simplicity, we have assumed in equation (1) that the flow is spherical. A more accurate treatment would replace $`R^2`$ with $`RH`$, where the scale height $`H=c_s/\mathrm{\Omega }_k`$. This would introduce minor differences in some of the quantitative results.
In the case of accretion onto a NS we expect the flow to slow down as it settles on the stellar surface, and we expect the density in this settling zone to be significantly higher than for a BH. The increased density would cause more efficient transfer of energy from protons to electrons via Coulomb collisions and more efficient radiation from the electrons. As we show below, this leads to a flow in which $`q^+`$, $`q^{}`$ and $`q_{\mathrm{Coul}}`$ are all of the same order, which is very different from the case of a BH ADAF, where $`q^+q^{},q_{\mathrm{Coul}}`$. Another feature of the settling zone, again the result of the large density, is that optically thin bremsstrahlung cooling (which is sensitive to $`\rho `$) dominates over self-absorbed synchrotron cooling. We therefore neglect synchrotron emission in our analysis. For simplicity, we neglect also thermal conduction. Comptonization is important over part of the settling zone. However, we have not been able to derive useful analytical results with both bremsstrahlung and Comptonization included. Therefore, we neglect Comptonization in this section, and discuss its effects in §3.3.
The set of equations (1)–(5) must satisfy certain boundary conditions at the neutron star. First, as the flow approaches the surface of the star at $`R=R_{NS}`$, the radial velocity must become very much smaller than the local free-fall velocity. Second, the angular velocity must approach the angular velocity of the star $`\mathrm{\Omega }_{NS}`$. We use the dimensionless parameter
$$s\frac{\mathrm{\Omega }_{NS}}{\mathrm{\Omega }_K(R_{NS})}$$
(6)
to represent the spin of the NS.
The radius of the star, and its spin, are the two principal boundary conditions applied at the inner edge of the accretion flow. We assume that the star is unmagnetized, so there are no magnetospheric effects to consider. Two outer boundary conditions, namely the temperature and angular velocity of the gas, are determined by the properties of the gas as it is introduced into the accretion flow on the outside (e.g. from an ambient medium or from a different type of accretion flow such as a thin disk). These outer boundary conditions have little effect on the interior of the flow see §5 of this paper, and Narayan, Kato & Honma 1997; but see also Yuan 1999. An additional important boundary condition is the mass accretion rate $`\dot{M}`$, which is determined by external conditions and which we take to be constant.
### 2.1 Inner Settling Solution
We consider first the inner region of the flow, $`R_{NS}R10^{2.5}R_S`$, where $`R_S=2GM/c^2`$ is the Schwarzchild radius. In this region we expect a two-temperature plasma, with $`T_p>T_e`$, in which the electrons are relativistic and the protons are non-relativistic: $`\theta _e1,\theta _p1`$. The viscous heating rate of the gas, the energy transfer rate from the protons to the electrons via Coulomb collisions, and the cooling rate of the electrons via bremsstrahlung emission are given by
$`q^+`$ $`=`$ $`\alpha {\displaystyle \frac{\rho c_s^2R^2}{\mathrm{\Omega }_K}}\left({\displaystyle \frac{d\mathrm{\Omega }}{dR}}\right)^2,`$ (7)
$`q_{\mathrm{Coul}}`$ $`=`$ $`Q_{\mathrm{Coul}}\rho ^2{\displaystyle \frac{\theta _p}{\theta _e}},Q_{\mathrm{Coul}}=4\pi r_e^2\mathrm{ln}\mathrm{\Lambda }{\displaystyle \frac{m_ec^3}{m_p^2}},`$ (8)
$`q^{}`$ $`=`$ $`Q_{\mathrm{ff},\mathrm{R}}\rho ^2\theta _e,Q_{\mathrm{ff},\mathrm{R}}=48\alpha _fr_e^2{\displaystyle \frac{m_ec^3}{m_p^2}},`$ (9)
where $`\alpha _f`$ is the fine structure constant, $`r_e`$ is the classical electron radius, $`\mathrm{ln}\mathrm{\Lambda }20`$ is the Coulomb logarithm, $`c_s^2c^2\theta _p`$, and we have neglected logarithmic corrections to the relativistic free-free emissivity. The subscript “R” in $`Q_{\mathrm{ff},\mathrm{R}}`$ denotes relativistic bremsstrahlung.
We now make a number of simplifications in equations (1)–(5). First, we neglect the radial velocity term $`vdv/dR`$ in equation (2). Second, we assume that $`\dot{J}`$ dominates over $`\dot{M}\mathrm{\Omega }R^2`$ on the right-hand-side of equation (3) and we neglect the latter term. Third, we neglect the entropy terms in equations (4), (5); that is, we assume that the heating, cooling and energy transfer terms in these equations dominate over the entropy gradient terms. All of these assumptions are justified a posteriori below. The equations then read
$`\dot{M}=4\pi R^2\rho v,`$ (10)
$`\left(\mathrm{\Omega }^2\mathrm{\Omega }_K^2\right)R={\displaystyle \frac{1}{\rho }}{\displaystyle \frac{d}{dR}}\left(\rho c_s^2\right),`$ (11)
$`4\pi \alpha {\displaystyle \frac{\rho c_s^2R^4}{\mathrm{\Omega }_K}}{\displaystyle \frac{d\mathrm{\Omega }}{dR}}=\dot{J},`$ (12)
$`(1\delta )q^+=q_{\mathrm{Coul}}=q^{}q^+\delta ,`$ (13)
where $`q^+,q_{\mathrm{Coul}},\text{ and }q^{}`$ are given by equations (7)–(9).
It is straightforward to show that the simplified equations (10)–(13) have the following self-similar solution,
$`\rho =\rho _0r^2,\theta _p=\theta _{p0}r^1,\theta _e=\theta _{e0}r^{1/2},`$
$`\mathrm{\Omega }=\mathrm{\Omega }_0r^{3/2},v=v_0r^0,`$ (14)
where $`r=R/R_S`$ is a dimensionless radius. The normalization coefficients are uniquely related to each other as follows
$`\theta _{p0}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left(1s^2\right),`$ (15a)
$`\theta _{e0}`$ $`=`$ $`\left({\displaystyle \frac{Q_{\mathrm{Coul}}}{Q_{\mathrm{ff},\mathrm{R}}}}{\displaystyle \frac{\theta _{p0}}{(1\delta )}}\right)^{1/2}=\left({\displaystyle \frac{\pi \mathrm{ln}\mathrm{\Lambda }}{12\alpha _f}}{\displaystyle \frac{\theta _{p0}}{(1\delta )}}\right)^{1/2}10.9\left({\displaystyle \frac{1s^2}{1\delta }}\right)^{1/2},`$ (15b)
$`\mathrm{\Omega }_0`$ $`=`$ $`\mathrm{\Omega }_{K0}s7.19\times 10^4m^1s\text{ rad/s},`$ (15c)
$`\rho _0`$ $`=`$ $`{\displaystyle \frac{9c^2\mathrm{\Omega }_{k0}}{4Q_{\mathrm{Coul}}}}\alpha (1\delta )\theta _{e0}s^28.10\times 10^4m^1\alpha (1\delta )\theta _{e0}s^2\text{ g/cm}^3,`$ (15d)
$`v_0`$ $`=`$ $`{\displaystyle \frac{\dot{M}}{4\pi R_S^2\rho _0}},`$ (15e)
$`\dot{J}`$ $`=`$ $`6\pi \alpha c^2R_S^3\rho _0\theta _{p0}s={\displaystyle \frac{9\pi }{2\sqrt{3}}}{\displaystyle \frac{(GM)^2c}{\sqrt{Q_{\mathrm{Coul}}Q_{\mathrm{ff},\mathrm{R}}}}}\alpha ^2(1\delta )^{1/2}s^3\left(1s^2\right)^{3/2}.`$ (15f)
Here, $`\mathrm{\Omega }_{K0}=\mathrm{\Omega }_K(R_S)`$ and the dimensionless spin parameter is $`s=\mathrm{\Omega }(R_S)/\mathrm{\Omega }_{K0}=\mathrm{\Omega }_{NS}/\mathrm{\Omega }_K(R_{NS})=\mathrm{\Omega }(R)/\mathrm{\Omega }_K(R)`$, as introduced in equation (6). Recall that $`s`$ is a boundary condition of the problem.
The range of $`r`$ over which the solution is valid is determined by the twin requirements that the protons be non-relativistic and that the electrons be relativistic. The former condition is satisfied for any $`r>1`$, while the latter condition requires $`r<\theta _{e0}^2120`$. A third condition is that Comptonization should be negligible (since we have assumed this). As we show in §3.3, this last condition requires $`r>\text{few tens}`$, with the exact limit depending on the NS spin, $`s`$, and viscosity, $`\alpha `$.
The radial velocity of the solution is independent of $`r`$, whereas the local free-fall velocity varies as $`v_{ff}=c/\sqrt{2r}`$. Thus, $`v/v_{ff}`$ decreases with decreasing radius. This shows that the solution corresponds to a settling flow and that it is quite different from self-similar ADAFs around BHs, where $`v/v_{ff}`$ either is constant (Narayan & Yi 1994, 1995a; Manmoto et al. 2000) or increases with decreasing $`r`$ (Narayan et al. 2000; Quataert & Gruzinov 2000). Although $`v/v_{ff}`$ is quite small as the flow approaches the NS surface, $`v`$ itself is still fairly large. At the NS surface, $`v`$ must reduce substantially from its self-similar value. As the numerical results of §4 show, this happens in a boundary layer where the accreting material cools catastrophically to a temperature that is orders of magnitude below virial. The boundary layer is distinct from the settling zone which is described by the above self-similar solution.
The angular velocity of the gas is a fixed fraction of the local Keplerian angular velocity, the ratio being determined by the dimensionless spin $`s`$ of the star. The gas in the settling solution radiates most of the energy dissipated through viscosity. In fact, the rates of viscous heating, Coulomb energy transfer and radiative emission are all equal, which is achieved by a suitable choice of the density, electron temperature and proton temperature in the gas. The two temperatures have universal forms, with only a weak dependence on $`s`$, while the density has a strong dependence on $`s`$.
The most outstanding feature of the self-similar solution is that, except for the radial velocity, none of the other gas parameters has any dependence on $`\dot{M}`$.
The fact that $`\dot{J}`$ is negative implies that the accretion flow removes angular momentum from the star and spins it down. This behavior is quite different from that seen in thin disks (Popham & Narayan 1991; Paczyński 1991), where for most choices of the stellar spin parameter $`s`$, the accretion disk spins up the star with a torque $`\dot{J}_{\mathrm{thin}}\dot{M}\mathrm{\Omega }_K(R_{NS})R_{NS}^2`$. Only when $`s`$ is very close to unity does the torque become negative. In contrast, for the self-similar solution derived here, the torque is negative for all values of $`s`$ (except extremely small values, see the discussion in §4). Moreover, $`\dot{J}`$ is independent of $`\dot{M}`$. Equivalently, the dimensionless torque, $`j=\dot{J}/\dot{M}\mathrm{\Omega }_K(R_{NS})R_{NS}^2`$, which is $`1`$ under most conditions for a thin disk, here takes on the value
$`j`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }}{8\sqrt{2}}}{\displaystyle \frac{m_p}{m_e}}{\displaystyle \frac{\eta }{\sqrt{\alpha _f\mathrm{ln}\mathrm{\Lambda }}}}r_{NS}^{1/2}(1\delta )^{1/2}\alpha ^2\dot{m}^1s^3\left(1s^2\right)^{3/2}`$ (16)
$``$ $`43\alpha ^2\dot{m}^1s^3\left(1s^2\right)^{3/2},`$
where $`\dot{m}=\dot{M}/\dot{M}_{\mathrm{Edd}}`$ is the mass accretion rate in Eddington units, with $`\dot{M}_{\mathrm{Edd}}=1.39\times 10^{18}m\text{ g/s}`$, (for a nominal $`\eta =0.1`$), $`r_{NS}3`$. Note that $`j`$ could be very large at low $`\dot{m}`$.
We now check under what conditions the approximations we made earlier are valid. First, we neglected the term $`vdv/dR`$ in equation (2). This is obviously valid since the self-similar solution has $`v=\mathrm{constant}`$.
Second, we assumed that $`|\dot{J}||\dot{M}\mathrm{\Omega }R^2|`$ in equation (3). Using (16) this condition may be cast into the form $`|j|sr^{1/2}/\sqrt{3}`$, or equivalently,
$$\dot{m}74\alpha ^2s^2r^{1/2}.$$
(17)
For $`\alpha 0.1,s0.3`$, and assuming a radial extent of $`r10^2`$ for the flow, we require $`\dot{m}7\times 10^3`$. A direct numerical simulation (§4) shows that the analytical solution is valid even for values of $`\dot{m}`$ that are a factor of a few larger than this limit.
Third, we neglected the entropy terms $`\rho _{p,e}vT_{p,e}ds_{p,e}/dR`$ in equations (4) and (5). For the protons, using the solution (14), we can show that the left-hand-side of equation (4) varies as $`r^4`$ while the right-hand-side varies $`r^{4.5}`$. Thus, with decreasing $`r`$, the entropy term becomes progressively less important than the other terms, thereby confirming the validity of the approximation. In the case of the electrons, the entropy is always small since $`s_e(m_eT_e/m_pT_p)s_ps_p`$.
### 2.2 Outer Settling Solution
For $`r>10^{2.5}`$, both protons and electrons are non-relativistic and the solution described in the previous subsection is not valid. Interestingly, another self-similar solution may be derived for this region of the flow. This solution has a nearly one-temperature plasma with $`T_pT_eT_p,T_e`$. The free-free cooling takes the form
$$q^{}Q_{\mathrm{ff},\mathrm{NR}}\rho ^2\theta _e^{1/2},Q_{\mathrm{ff},\mathrm{NR}}=5\sqrt{2}\pi ^{3/2}\alpha _f\sigma _T\frac{m_ec^3}{m_p^2},$$
()
where $`\sigma _T`$ is the Thompson cross-section, and the subscript $`\mathrm{NR}`$ stands for non-relativistic. Since the gas is effectively one-temperature, Equation (13) simplifies to
$$q^+q^{}.$$
()
We do not need to consider the Coulomb transfer rate $`q_{\mathrm{Coul}}`$, since this quantity is proportional to $`(T_pT_e)`$ and can be adjusted to have the right magnitude with small adjustments of the two temperatures. In the non-relativistic regime, $`\theta _e<1,\theta _p(m_e/m_p)\theta _e1`$, and the Coulomb transfer rate is
$$q_{\mathrm{Coul}}=\frac{3}{\sqrt{2\pi }}\frac{m_e}{m_p}\frac{\sigma _Tc}{m_p^2}\mathrm{ln}\mathrm{\Lambda }\rho ^2\frac{kT_pkT_e}{\theta _e^{3/2}e^{1/\theta _e}}.$$
(18)
¿From the condition $`q_{\mathrm{Coul}}q^{}`$ it follows that
$$kT_pkT_e\frac{10}{3\pi ^2}\frac{\alpha _f}{\mathrm{ln}\mathrm{\Lambda }}\sqrt{m_pm_e}c^2\theta _e^2e^{1/\theta _e},$$
(19)
which is exponentially small for $`\theta _e<1`$.
The flow is again described by the self-similar solution (14), with the following two exceptions:
$$\theta _e=\theta _{e0}r^1,\theta _{e0}=\frac{m_p}{m_e}\theta _{p0},$$
()
$$\rho =\rho _0r^2,\rho _0=\frac{9\alpha c^2\mathrm{\Omega }_{K0}}{2Q_{\mathrm{ff},\mathrm{NR}}}\frac{\theta _{p0}s^2}{\theta _{e0}^{1/2}}=0.12m^1\theta _{p0}^{1/2}\alpha s^2\text{ g/cm}^3,$$
()
where we have used the fact that the total pressure $`p=p_p+p_e=2p_p`$.
This solution is valid only if $`\dot{m}`$ is quite small, cf. equation (17):
$$\dot{m}<2.2\times 10^3\alpha _{0.1}^2s_{0.3}^2r_3^{1/2},$$
(20)
where $`r_3=r/10^3`$, $`\alpha _{0.3}=\alpha /0.3`$, and $`s_{0.1}=s/0.1`$. For greater $`\dot{m}`$, a self-similar power-law solution does not exist; numerically computed solutions exhibit non-power-law behavior, as discussed in §4.
## 3 Properties of the self-similar solution
### 3.1 Spin-Up/Spin-Down of the Neutron Star
The rate of spin-up of the accreting NS is given by
$`{\displaystyle \frac{d}{dt}}\left(I_{NS}\mathrm{\Omega }\right)`$ $`=`$ $`\dot{J}\dot{M}\mathrm{\Omega }(R_{NS})R_{NS}^2`$ (21)
$``$ $`43s^3\alpha ^2\dot{M}_{\mathrm{Edd}}\mathrm{\Omega }_K(R_{NS})R_{NS}^2,`$
where $`I_{NS}`$ is the moment of inertia of the NS. We have made use of the fact that $`|\dot{J}||\dot{M}\mathrm{\Omega }(R_{NS})R_{NS}^2|`$ for the self-similar solution, and used equation (16) for $`\dot{J}`$. The negative sign in the final expression implies that the accretion flow spins down the star. The above equation is for an unmagnetized NS. If the NS has a magnetosphere, the inner edge of the accretion flow is at the magnetospheric radius, $`R_m`$. In this case, let us define $`s`$ by $`\mathrm{\Omega }_{NS}=s\mathrm{\Omega }_K(R_m)=s\mathrm{\Omega }_K(R_{NS})(R_m/R_{NS})^{3/2}`$. Substituting this in equation (21) with $`I_{NS}=constant`$ and integrating, we obtain
$$s=\frac{s_0}{\sqrt{1+t/\tau }},\tau =\frac{I_{NS}}{86s_0^2\alpha ^2\dot{M}_{\mathrm{Edd}}R_{NS}^2}\left(\frac{R_m}{R_{NS}}\right)^{3/2},$$
(22)
where $`s_0=s(t=0)`$. The same result is valid for an unmagnetized NS by setting $`R_m=R_{NS}`$. The quantity $`\tau `$ is the characteristic spin-down time of the NS. For a spherical NS of constant density, $`I_{NS}=2MR_{NS}^2/5=(0.8\times 10^{33}\text{ g})mR_{NS}^2`$. Substituting this expression, we obtain the spin-down rate $`\dot{P}_{NS}/P_{NS}=\tau ^1`$ with
$$\tau 6.7\times 10^{12}s^2\alpha ^2\left(\frac{R_m}{R_{NS}}\right)^{3/2}\text{ s}=2\times 10^8s_{0.1}^2\alpha _{0.1}^2\left(\frac{R_m}{R_{NS}}\right)^{3/2}\text{ yr}.$$
(23)
Note the remarkable fact that the spin-down time scale is independent of the mass of the NS, and the mass accretion rate! For the magnetic case, the rate depends on the radius ratio $`R_m/R_{NS}`$.
It is customary to express the spin-down rate as $`\dot{P}_{NS}/P_{NS}^2`$. Writing
$$P_{NS}=\frac{2\pi }{s\mathrm{\Omega }_K(R_{NS})}\left(\frac{R_m}{R_{NS}}\right)^{3/2}$$
(24)
and $`\mathrm{\Omega }_K(R_{NS})10^4m_{1.4}^1\text{ rad/s}`$, where $`R_{NS}=3R_S`$ and $`m_{1.4}=M/(1.4M_{\mathrm{}})`$, we obtain
$$\frac{\dot{P}_{NS}}{P_{NS}^2}2.4\times 10^{10}m_{1.4}^1\alpha ^2s^3\text{ s}^2=2.7\times 10^{12}m_{1.4}^1\alpha _{0.3}^2s_{0.5}^3\text{ s}^2,$$
(25)
where $`s_{0.5}=s/0.5`$. This spin-down rate is in good agreement with observational data on the spin-down of X-ray pulsars for which Yi, Wheeler & Vishniac (1997) invoked ADAFs: 4U 1626-67 has $`\dot{P}/P^28\times 10^{13}\text{ s}^2`$ and $`P=7.7\text{ s}`$; OAO 1657-415 has $`\dot{P}/P^22\times 10^{12}\text{ s}^2`$ and $`P=38\text{ s}`$, and GX 1+4 has $`\dot{P}/P^23.7\times 10^{12}\text{ s}^2`$ and $`P=122\text{ s}`$. Since the spin-down rate is quite sensitive to $`\alpha `$ and $`s`$, the observed data in individual systems can be fitted by small adjustment of these parameters.
### 3.2 Luminosity
In computing the luminosity of the accretion flow, we must allow for the energy release in both the boundary layer and the self-similar settling zone. We calculate their luminosities separately.
Radiation from the self-similar settling flow may be calculated following the methods described by Popham & Narayan (1995) for a thin disk. This method assumes that the luminosity at a given radius is determined by the local viscous energy production. This is a legitimate approximation for the settling flow in which $`q^{}=q^+`$. Keeping only the dominant terms, we find
$$L_{SS}=\frac{GM_{NS}\dot{M}}{R_{in}}\left(1+\frac{1}{2}s^2js\right)+\dot{M}_{P_{\mathrm{in}}}^{P_{\mathrm{out}}}\frac{dP}{\rho },$$
(26)
where $`R_{in}=R_{NS}+\mathrm{\Delta }_{BL}`$ is the inner radius of the self-similar zone and $`\mathrm{\Delta }_{BL}R_{NS}`$ is the thickness of the boundary layer. Here the first two terms, $`(1+s^2/2)`$, represent the luminosity, associated with potential energy of the infalling gas, $`L_{\mathrm{pot}}`$, the third term $`js`$ is the luminosity, associated with the rotational energy extracted from the star, $`L_{\mathrm{rot}}`$ (note, $`j<0`$ in the self-similar solution), and the final integral is the “enthalpy correction”, $`L_{\mathrm{enth}}`$. Using the analytical solution (14)–(15) and assuming $`P_{\mathrm{out}}=0`$ for simplicity, we obtain
$`L_{\mathrm{pot}}`$ $`=`$ $`\dot{M}_{\mathrm{Edd}}c^2{\displaystyle \frac{\dot{m}}{2r_{NS}}}\left(1+{\displaystyle \frac{s^2}{2}}\right),`$ (27a)
$`L_{\mathrm{rot}}`$ $`=`$ $`43\dot{M}_{\mathrm{Edd}}c^2{\displaystyle \frac{\alpha ^2}{2r_{NS}}}(1\delta )^{1/2}s^4\left(1s^2\right),`$ (27b)
$`L_{\mathrm{enth}}`$ $`=`$ $`\dot{M}_{\mathrm{Edd}}c^2{\displaystyle \frac{\dot{m}}{2r_{NS}}}\left(1s^2\right).`$ (27c)
Note that the leading terms in $`L_{\mathrm{pot}}`$ and $`L_{\mathrm{enth}}`$ cancel each other exactly. The luminosity of the self-similar settling zone is thus
$$L_{SS}6.2\times 10^{34}mr_3^1\dot{m}_2s_{0.1}^2+8.9\times 10^{33}mr_3^1\alpha _{0.1}^2s_{0.1}^4\text{ erg\hspace{0.17em}s}^1,$$
(28)
where $`\dot{m}_2=\dot{m}/0.01`$, $`s_{0.1}=s/0.1`$, $`r_3=r_{NS}/3`$, and we have assumed $`s1`$. Note that luminosity not associated with with dissipation of rotational energy, represented by the first term in equation (28), is much less than the commonly assumed $`GM_{NS}\dot{M}/R_{NS}`$. This is because the negative enthalpy term has large magnitude, as a result of the fact that the settling flow is akin to a pressure supported, quasi-stationary atmosphere.
The second term in equation (28) is the luminosity of the settling zone. Since the self-similar solution for this zone is independent of $`\dot{m}`$, the luminosity too shows no $`\dot{m}`$ dependence. Indeed, the luminosity remains finite even as $`\dot{m}0`$. How is this possible, and where does the energy come from? The answer is that the luminosity of the settling zone is supplied by the central star. As the star spins down, it does work on the accretion flow and the energy released comes out as bremsstrahlung radiation.
The boundary layer luminosity requires a different method of calculation since viscous energy production is negligible in this zone: $`\mathrm{\Omega }`$constant, and so $`q^+\left(d\mathrm{\Omega }/dR\right)^20`$. As the accreting gas cools in the boundary layer, starting from a nearly virial temperature $`10^{12}`$ K on the outside down to the NS temperature $`10^7`$ K near the surface, the thermal energy in the gas is emitted as radiation. To estimate the luminosity, we use the energy balance equation, which is the sum of equations (4),(5):
$$q^{}=\frac{\rho v}{\gamma 1}\frac{dc_s^2}{dR}vc_s^2\frac{d\rho }{dR}=\frac{\gamma }{\gamma 1}\rho v\frac{dc_s^2}{dR}v\frac{dP}{dR}.$$
(29)
We can neglect the $`dP/dR`$ term because the pressure $`P`$ is essentially constant in the boundary layer. To obtain the luminosity we integrate over the boundary layer
$$L_{BL}=q^{}4\pi R^2𝑑R=\frac{\gamma }{\gamma 1}4\pi R^2\rho v\frac{dc_s^2}{dR}𝑑R=\frac{\gamma }{\gamma 1}\dot{M}\mathrm{\Delta }c_s^2.$$
(30)
Since $`c_s^2`$ starts from nearly virial value and reaches close to zero, $`\mathrm{\Delta }c_s^2GM_{NS}/R_{NS}`$. More precisely, $`\mathrm{\Delta }c_s^2=c^2\mathrm{\Delta }(\theta _p+\theta _e)c^2\theta _{p0}r_{NS}^1`$. Therefore, the boundary layer luminosity is
$$L_{BL}=\frac{\gamma }{\gamma 1}\dot{M}_{\mathrm{Edd}}c^2\frac{\dot{m}}{6r_{NS}}\left(1s^2\right)1.7\times 10^{36}m\dot{m}_2r_3^1,$$
(31)
where we have assumed $`\gamma =5/3`$. The total luminosity of the system is $`L=L_{SS}+L_{BL}`$.
### 3.3 Effect of Comptonization
Using the self-similar solution (14),(15), we may readily estimate the electron scattering optical depth and the $`y`$-parameter.<sup>1</sup><sup>1</sup>1Here we just estimate where the effect of Comptonization becomes significant. For better analytical approximations see, for instance, Dermer, Liang, & Canfield (1991); Titarchuk & Lyubarskij (1995). (In the latter paper, the expression for $`y`$ is not given, but it can be inferred using equation : $`y=\tau _{es}[(\alpha +3)\theta /(1+\theta )+4d_0^{1/\alpha }\theta ^2]`$.) Comptonization of free-free radiation has also been considered by Titarchuk (1989). The optical depth is
$`\tau _{\mathrm{es}}`$ $``$ $`\rho \kappa _{\mathrm{es}}R10^3\alpha (1\delta )^{1/2}(1s^2)^{1/2}s^2r^1`$ (32)
$``$ $`\alpha _{0.1}s_{0.1}^2r^1,`$
where $`\kappa _{\mathrm{es}}=\sigma _T/m_p`$ is the electron scattering opacity for ionized hydrogen. Since $`r3`$, we see that $`\tau _{\mathrm{es}}1/3`$ for reasonable parameters and the radiation is optically thin to electron scattering. The $`y`$-parameter is
$`y`$ $`=`$ $`16\theta _e^2\tau _{\mathrm{es}}2\times 10^6\alpha (1\delta )^{1/2}(1s^2)^{3/2}s^2r^2`$ (33)
$``$ $`2\times 10^3\alpha _{0.1}s_{0.1}^2r^2.`$
The radius at which $`y1`$ is
$$r_c45\alpha _{0.1}^{1/2}s_{0.1}.$$
(34)
Above this radius the inverse Compton scattering is small and the self-similar solution is valid. For $`r<r_c`$, however, Comptonization is important and the electron temperature profile will be modified from the self-similar form. Since the electron-proton collisions are relatively weak (the plasma is two-temperature), other quantities, e.g., the density, proton temperature, etc., are unaffected. Comptonization is unimportant for low-viscosity flows, $`\alpha 0.01`$ around slowly rotating NSs, $`s0.01`$, because then $`r_c<r_{NS}`$.
### 3.4 Spectrum
We now estimate the spectrum of radiation emitted from the settling accretion flow. Let us neglect inverse Compton scattering for the moment. The relativistic bremsstrahlung emissivity is approximated as $`ϵ_\nu \rho ^2\mathrm{exp}(h\nu /kT_e)\text{ erg cm}^3\text{ s}^1\text{ Hz}^1`$. Therefore the luminosity per unit frequency is
$`L_\nu `$ $``$ $`{\displaystyle _{R_{NS}}^{\mathrm{}}}\rho ^2e^{h\nu /kT_e}2\pi R^2𝑑R`$ (35)
$``$ $`{\displaystyle _{1/\nu _m}^{\mathrm{}}}t^3e^{\nu t}𝑑t\nu ^2\mathrm{\Gamma }(2,\nu /\nu _m),`$
where $`\mathrm{\Gamma }(a,z)=_z^{\mathrm{}}t^{a1}e^t𝑑t`$ is the incomplete gamma-function and $`\nu _m=kT_e(R_{NS})/h`$ is the maximum frequency. Above $`\nu _m`$ the spectrum falls exponentially and below $`\nu _m`$ it is nearly flat. We may, thus, replace the exponential in the integral with a square function which is equal to unity for $`\nu <\nu _m`$ and 0 for $`\nu >\nu _m`$. With this approximation
$$L_\nu \frac{3}{2}\frac{L_{SS}}{\nu _m}\left(1\frac{\nu ^2}{\nu _m^2}\right),$$
(36)
where $`L_{SS}=L_\nu 𝑑\nu `$ is the total luminosity of the self-similar flow, represented by equation (28). The break frequency, $`\nu _m`$, is roughly given by $`h\nu _m2.7\text{ MeV}`$ for a typical electron temperature $`T_{e,\mathrm{max}}10^{10.5}{}_{}{}^{}\text{K}`$ \[cf., equation (15)\]. At a typical x-ray energy, $`h\nu 3\text{ keV}`$, the observed luminosity per decade is
$$\nu L_\nu 1.7\times 10^{31}m\alpha _{0.1}^2s_{0.1}^4\left(\frac{h\nu }{3\text{ keV}}\right)\text{ erg s}^1,$$
(37)
i.e., $`\nu L_\nu 1.5\times 10^{32}`$ for a 300 Hz neutron star ($`s_{0.1}1.6`$). The luminosity per decade is much greater at higher photon energies and may be as high as $`\text{few}\times 10^{34}10^{35}\text{ erg/s}`$ at $`h\nu \text{ MeV}`$.
As shown in the previous section, Comptonization becomes important below the radius $`r_c`$. At $`r_c`$, $`y1`$ and the electron temperature is
$$T_e(r_c)2.7\text{ MeV}/\sqrt{r_c}400\alpha _{0.1}^{1/4}s_{0.1}^{1/2}\text{ keV}.$$
(38)
For $`r<r_c`$, the electron temperature will be determined self-consistently by Compton cooling rather than by bremsstrahlung emission. Computing the spectrum from this region is beyond the scope of the paper. We also do not attempt to calculate the spectrum of the radiation from the boundary layer.
### 3.5 Bernoulli parameter
It is known that the Bernoulli parameter of the accreting gas in BH ADAFs is positive for a wide range of $`r`$ (Narayan & Yi 1994, 1995a; Narayan, Kato & Honma 1997), and it has been suggested that the positive Bernoulli parameter may trigger strong winds or jets in these systems Narayan & Yi 1994, 1995a; Blandford & Begelman 1999, but see Abramowicz et al. 2000. Igumenshchev & Abramowicz (1999, 2000) confirmed with numerical simulations that strong outflows are produced from BH ADAFs when $`\alpha 1`$.
Normalizing the Bernoulli parameter, $`Be`$, by $`(\mathrm{\Omega }_KR)^2`$, and using equations (14),(15), we find that the self-similar settling flow has
$`b{\displaystyle \frac{Be}{\mathrm{\Omega }_K^2R^2}}`$ $`=`$ $`{\displaystyle \frac{1}{v_k^2}}\left({\displaystyle \frac{1}{2}}v^2+{\displaystyle \frac{1}{2}}\mathrm{\Omega }^2R^2\mathrm{\Omega }_K^2R^2+{\displaystyle \frac{\gamma }{\gamma 1}}c_s^2\right)`$ (39)
$`=`$ $`{\displaystyle \frac{v_0^2}{c^2}}r+{\displaystyle \frac{s^2}{2}}1+{\displaystyle \frac{\gamma }{\gamma 1}}{\displaystyle \frac{1s^2}{3}}`$
$``$ $`{\displaystyle \frac{2\gamma 3}{3(\gamma 1)}}{\displaystyle \frac{s^2}{2}}{\displaystyle \frac{3\gamma }{3(\gamma 1)}},`$
where $`\gamma `$ is the mean adiabatic index of the flow and in the last expression we have neglected the term in $`v_0^2`$ since it is negligible deep inside the settling region.
We see that the second term in the final expression is always negative. This term is proportional to $`s^2`$, which means that a more rapidly spinning NS stabilizes the accretion flow against outflows more effectively than a slower spinning star. The physical explanation is as follows. The centrifugal force increases with increasing $`s`$, which has the effect of making it easier for the gas to escape. At the same time, however, the centrifugal force causes the radial infall velocity to decrease, which increases the time available for cooling. The temperature and the gas pressure go down, making the gas more gravitationally bound. The net contribution of the two effects turns out to be negative.
The first term in the last line of equation (39) can be either positive or negative, depending on the value of $`\gamma `$. Combining the two terms, we find that the gas is gravitationally bound and unable to flow out in a wind (i.e. $`b<0`$) if the adiabatic index satisfies
$$\gamma >\frac{3\left(1\frac{1}{2}s^2\right)}{2\frac{1}{2}s^2}.$$
(40)
For $`s^21`$, the condition is $`\gamma >1.5`$; that is, the the accretion flow can produce a wind and/or a collimated outflow only if $`\gamma <1.5`$ and is stable to such outflows if $`\gamma >1.5`$. Normally, we expect $`\gamma `$ to be close to 5/3 for the accreting gas.
### 3.6 Stability to Convection
It is well known that if the entropy increases inwards in a gravitationally-bound non-rotating system, the gas is convectively unstable; otherwise the flow is stable. The specific entropy profile in the settling accretion flow around a NS can be readily calculated from equations (4),(5) using (14),(15). This gives
$$\frac{ds}{dR}=\frac{k}{m_p}\frac{1}{\gamma 1}\frac{d}{dR}\mathrm{ln}\left(\frac{c_s^2}{\rho ^{\gamma 1}}\right)=\frac{k}{m_p}\frac{2\gamma 3}{\gamma 1}\frac{1}{R}.$$
(41)
We see that the entropy increases outwards for $`\gamma >1.5`$ and inwards for $`\gamma <1.5`$. Hence if $`\gamma >1.5`$ the flow is stable against convection, while if $`\gamma <1.5`$ the flow is convectively unstable.
In the presence of rotation, the analysis is a little more complicated. Narayan et al. (2000) and Quataert & Gruzinov (2000) discuss the generalization of the Schwarzchild criterion for accretion flows with rotation. If the gas motions are restricted to the equatorial plane of a height-integrated flow, convective stability requires the following effective frequency to be positive:
$$N_{\mathrm{eff}}^2=N^2+\kappa ^2>0,$$
(42)
where $`N`$ is the Brunt-Väisälä frequency and $`\kappa `$ is the epicyclic frequency, $`\kappa =\mathrm{\Omega }`$ for $`\mathrm{\Omega }R^{3/2}`$. For a power-law flow with $`\rho R^a`$ and $`\mathrm{\Omega }(R)=s\mathrm{\Omega }_KR^{3/2}`$ with $`s^2=1(1+a)c_0^2`$, this criterion may be written as follows see Narayan et al. 2000 for more discussion
$$N_{\mathrm{eff}}^2=\mathrm{\Omega }_K^2\left([(\gamma +1)a(\gamma 1)]\frac{(1+a)c_0^2}{\gamma }+1\right)>0.$$
(43)
Since for the self-similar settling solution $`a=2`$, the stability criterion (43) becomes
$$N_{\mathrm{eff}}^2=\frac{\mathrm{\Omega }_K^2}{\gamma }\left[(2\gamma 3)+s^2(3\gamma )\right]>0$$
(44)
which yields that the flow is convectively stable if
$$\gamma >\frac{3(1s^2)}{2s^2}.$$
(45)
This condition is different from the stability criterion against outflows, given in equation (40).
Following the techniques developed by Quataert & Gruzinov (2000), Narayan et al. (2000) have also presented a more general analysis of convection in a self-similar accretion flow. This analysis, which does not restrict motions to lie in the equatorial plane, assumes that $`v_\varphi `$ and $`c_s`$ are independent of the polar angle $`\theta `$ as is valid for a marginally convectively stable system, cf Quataert & Gruzinov 2000. Narayan et al. (2000) find that the most unstable region of the flow is near the rotation axis, $`\theta =0,\pi `$. They show that the marginal stability criterion for this polar fluid coincides with the condition for the positivity of the Bernoulli parameter. That is, a flow which is convectively stable at all $`\theta `$ has a negative Bernoulli parameter, while a flow which is convectively unstable for at least some values of $`\theta `$ has a positive Bernoulli parameter. (The Bernoulli parameter itself is independent of $`\theta `$.)
We have verified this result for the solutions presented in this paper. Specifically, when we apply to our solution the more general convective stability criterion given by equation (A9) of Narayan et al. (2000), we recover the condition (40) above, namely that the self-similar flow is convectively stable if and only if
$$\gamma >\frac{3\left(1\frac{1}{2}s^2\right)}{2\frac{1}{2}s^2}.$$
(46)
## 4 Comparison with Numerical Results
We have numerically solved the system of height-integrated two-temperature fluid equations (1)–(5) with boundary conditions. In the energy equations we assume that viscous dissipation only heats the protons (i.e. $`\delta =0`$). We include energy transfer from protons to electrons via Coulomb collisions, and we take the cooling of electrons to be purely by free-free emission, as discussed in §2. For these processes we use the expressions given in Narayan & Yi (1995b), which smoothly interpolate between the regimes of non-relativistic and relativistic electrons. We take into account the variation of the electron adiabatic index $`\gamma _e`$ with temperature (Chandrasekhar 1939; Esin et al. 1997), using a simple interpolation formula from Gammie & Popham (1998). We assume that the protons have $`\gamma _p=5/3`$.
We employ the gravitational potential of Paczyński & Wiita (1980) to mimic the effect of strong gravity near the NS surface. In this potential the Keplerian angular velocity takes the form
$$\mathrm{\Omega }_K^2=\frac{GM}{(RR_S)^2R}.$$
(47)
Note that the analytical work presented in the previous sections is based on a Newtonian potential.
We specify the boundary conditions as follows. We take the outer boundary of the flow to be at $`r_{\mathrm{out}}=10^6`$. At this radius we specify that the angular velocity is equal to its value in the self-similar ADAF solution of Narayan & Yi (1994), and that the proton and electron temperatures are both equal to the self-similar ADAF temperature. We assume that the accreting star is a $`1M_{\mathrm{}}`$ neutron star with a radius $`R_{NS}=3R_S=8.85`$ km, unless stated otherwise. At $`R=R_{NS}`$, we specify the value of the NS spin parameter, $`s=\mathrm{\Omega }_{NS}/\mathrm{\Omega }_K(R_{NS})`$, and we require the proton temperature of the flow to be $`T=\mathrm{few}\times 10^7\text{ K}T_{virial}`$. (We do not assume that the electron and proton temperatures are equal, but in fact they are equal.) We do not constrain the density of the gas in any way at either boundary.
The numerical problem as posed here has a family of solutions characterized by three dimensionless parameters: the mass accretion rate $`\dot{m}`$ (in Eddington units), the NS spin $`s`$ (in units of the Keplerian angular velocity at the NS surface), and the viscosity parameter $`\alpha `$. The angular momentum flux $`\dot{J}`$, or equivalently the dimensionless flux $`j=\dot{J}/\dot{M}\mathrm{\Omega }_K(R_{NS})R_{NS}^2`$, is an eigenvalue of the problem.
Figure Self-Similar Hot Accretion Flow onto a Neutron Star shows representative solutions for $`\alpha =0.1`$ and a range of values of $`\dot{m}`$ and $`s`$. The solutions clearly have three radial zones. For $`r>10^{2.5}`$, there is a one-temperature zone in which the gas properties vary roughly as power-laws of the radius. For $`r<10^{2.5}`$, there is a second power-law zone with a two-temperature structure. Finally, close to the NS, the flow has a boundary layer region. In this final region, the gas experiences run-away cooling, the velocity falls precipitously, and the density increases very rapidly. This region of the flow does not have power-law behavior.
The numerical solutions are unreliable in the boundary layer; thermal conduction (not included in the calculations) is probably very important here, and optical depth effects (also not included) will modify the radiation properties significantly. The solutions are suspect also in the inner region of the two-temperature power-law zone, where Comptonization is likely to be important. Outside these regions, however, the numerical solution is expected to be accurate.
The numerical results in the two-temperature power-law zone below $`r10^{2.5}`$ agree quite well with the analytical solution presented in §2.1. Curves corresponding to a given value of $`s`$ and different values of $`\dot{m}`$ coincide with one other to very good accuracy, as predicted by the analytical solution. This is best seen in the profiles of $`\rho `$ and $`\mathrm{\Omega }`$. Changing $`s`$ causes an up/down shift of the curves but does not affect the slopes of the curves. The temperature profiles are sensitive to the spin $`s`$, especially for large values of $`s`$. The radial velocity varies approximately as $`v\dot{m}`$ and is roughly consistent with $`vr^0`$ for $`s>0.1`$.
For $`r>10^{2.5}`$, the numerical solutions are in reasonable agreement with the one-temperature self-similar solution described in §2.2. The agreement is less perfect than in the previous zone. This is primarily because the analytical solution requires a very low value of $`\dot{m}`$ in order to be valid as far out as the outer radius $`r10^6`$ \[cf., equation (20)\]. The numerical models shown have larger values of $`\dot{m}`$ than this limit.
As we discussed in §2 and §3, the transport of angular momentum in a hot settling flow differs dramatically from the well-known behavior of a thin disk. The solid line in Figure Self-Similar Hot Accretion Flow onto a Neutron Star indicates the dependence of the dimensionless angular momentum eigenvalue $`j`$ as a function of the dimensionless NS spin $`s`$. The long-dashed line indicates the corresponding results for a thin disk (Popham & Narayan 1991), where $`j+1`$ for most values of $`s`$, and goes negative only for stars nearly at break-up. In contrast, in the settling flow, $`j`$ is negative for almost all values of $`s`$. For the particular choice of parameters, namely $`\alpha =0.1,\dot{m}=0.01`$, $`s\text{few}\times 10^1`$, we find that $`j`$ few.
The short-dashed line in Figure Self-Similar Hot Accretion Flow onto a Neutron Star shows the analytical formula for $`j`$, as given in equation (16). The agreement with the numerical results is good for a wide range of $`s`$ below about 0.5. For $`s>0.5`$, the numerically determined $`j`$ levels off at a constant negative value, whereas the analytical result shows $`|j|`$ decreasing rapidly. The main reason for the discrepancy is the neglect of the ram pressure term in the radial momentum equation in the analytical work. Note the interesting fact that super-Keplerian accretion ($`s>1`$) is, in principle, possible (provided one can arrange to have a star with super-Keplerian rotation). In a super-Keplerian flow, the ram pressure of the infalling gas supplies the radial momentum needed to push the gas onto the NS. For extremely small $`s0.1`$, the numerical solutions show $`j`$ to be slightly positive; we find $`j10^3`$, as indicated in the lower panel in Figure Self-Similar Hot Accretion Flow onto a Neutron Star. Here again the self-similar solution, which predicts a small negative value for $`j`$, breaks down. For very small $`s`$, the $`\dot{J}`$ term in equation (3) is comparable to or smaller than the $`\dot{M}\mathrm{\Omega }R^2`$ term which was omitted in deriving the analytical solution. This is the reason for the discrepancy. The precise value of $`s`$ at which the analytical solution breaks down depends on the choice of parameters ($`\alpha `$, $`\dot{m}`$), but is essentially independent of the outer radius. Thus, the transport of angular momentum through the flow and the spin-down of the star are determined by the boundary conditions at the stellar surface, and not by the outer boundary of the flow.
One of the most surprising features of the self-similar solution is that the angular momentum flux $`\dot{J}`$ is independent of $`\dot{m}`$; equivalently, the dimensionless eigenvalue $`j`$ is $`\dot{m}^1`$. The solid line in Figure Self-Similar Hot Accretion Flow onto a Neutron Star is a plot of $`j`$ as a function of $`\dot{m}`$ for a flow with $`\alpha =0.1`$ and $`s=0.3`$, as determined from the numerical solutions. The dashed line indicates for comparison the analytical scaling, $`j\dot{m}^1`$, with the proportionality constant given in equation (16). Note the very good agreement between the numerical results and the analytical self-similar solution.
The highest value of $`\dot{m}`$ up to which we could obtain a numerical solution is $`\dot{m}_{\mathrm{crit}}=0.0313`$. Beyond this critical value, there is no hot solution. (The value of $`\dot{m}_{\mathrm{crit}}`$ depends on $`s`$, $`\alpha `$ and $`r_{out}`$.) For larger $`\dot{m}`$, the density in the flow is so high that there is runaway free-free cooling and the gas is unable to remain hot. We presume that the accretion then occurs via a thin accretion disk.
Figure Self-Similar Hot Accretion Flow onto a Neutron Star plots, for selected values of $`\dot{m}`$, the luminosity per logarithmic interval of $`D`$, where $`D`$ is the fractional distance from the NS surface:
$$D=\frac{(RR_{NS})}{R_{NS}}.$$
(48)
We see that the luminosity at the peak of the curve is very insensitive to $`\dot{m}`$. The emission in the peak corresponds to radiation from the settling flow. This emission represents energy released by the spin-down of the NS, and its luminosity is independent of $`\dot{m}`$ \[cf. equation (28)\]. For radii below the peak, the curves do show a dependence on $`\dot{m}`$. The radiation here corresponds to boundary layer emission, which is proportional to $`\dot{m}`$ according to equation (31).
All the models described above have $`R_{NS}=3R_S`$. However, different equations of state predict slightly different NS radii, $`R_{NS}=24R_S`$. We have computed numerical models for this range of $`R_{NS}`$ and we find that the subsonic settling solution exists for the whole range. Qualitatively, the solutions with different $`R_{NS}`$ are very similar. As $`R_{NS}`$ decreases, the peak temperature is higher, as expected for the deeper potential.
## 5 Relationship of the Settling Flow to an ADAF
The accretion solution we have discussed so far radiates all the energy dissipated by viscosity, and is therefore “cooling-dominated.” On the other hand, it is known that a hot flow around a black hole is an “advection-dominated” accretion flow (ADAF). Both solutions have a two-temperature structure for $`r10^{2.5}`$ and both are very hot (nearly virial) for all $`r`$. How are these two types of accretion flows related to each other?
By solving the equations numerically for different boundary conditions, we have found that the two solutions are part of a single sequence of solutions in which the spin of the star, $`s`$, plays a pivotal role as a control parameter. For relatively rapidly rotating stars, with $`s0.1`$, we obtain the settling solution in our numerical experiments. However, as $`s`$ is decreased, we find that the settling solution smoothly transforms to an ADAF-type solution, which becomes well-established for $`s0.01`$. The transition is not sharp, so it is difficult to identify a specific transition point $`s=s_t`$ at which the transformation occurs. Numerical experiments indicate that the value of $`s_t`$ (however it is defined) is not very sensitive to $`R_{out},\gamma `$, and $`\dot{m}`$ and is, roughly, $`s_t0.040.06`$.
The change of the nature of the flow as $`s`$ is varied is illustrated in Fig. Self-Similar Hot Accretion Flow onto a Neutron Star. The solid and dotted curves correspond to two solutions with $`s=0.3`$ and $`s=0.01`$, respectively, with all other boundary conditions being the same. We see that the solutions are markedly different from each other. This is most clearly seen in the profiles of density, where the $`s=0.3`$ model has a logarithmic slope of -2, as appropriate for the cooling-dominated settling solution described in this paper, and the $`s=0.01`$ model has a slope of -3/2, as expected for a standard self-similar ADAF (Narayan & Yi 1994, 1995a). There is a similar difference also in the profiles of the radial velocity, where the two solutions have logarithmic slopes of -1/2 and 0, respectively.
An interesting feature of the $`s=0.01`$ ADAF-type solution is that it consists of two distinct segments. For large radii (in Fig. Self-Similar Hot Accretion Flow onto a Neutron Star, for radii outside $`r20`$), the flow corresponds to the standard ADAF discussed in the literature, with the scalings
$$\rho r^{3/2},c_s^2r^1,\mathrm{\Omega }r^{3/2},vr^{1/2}.$$
(49)
However, at smaller radii, the numerical solution indicates the presence of a second advection-dominated zone, a “settling ADAF,” which was first seen in numerical calculations described in Narayan & Yi (1994). This settling ADAF is seen in Fig. Self-Similar Hot Accretion Flow onto a Neutron Star as a zone that lies between the boundary layer region and the outer standard ADAF, with different slopes for $`\rho `$ and $`v`$. The radial extent of the settling ADAF zone may be quite large and, in general, depends on the flow parameters and boundary conditions.
A self-similar model of the settling ADAF may be readily obtained as follows. In an ADAF, energy is not radiated, therefore $`q^{}=0`$. Close to the star $`\mathrm{\Omega }\text{ constant}`$, therefore $`q^+=0`$. Equations (4), (5) then simplify to the condition of entropy conservation, $`ds/dR=0`$, which yields $`c_s^2\rho ^{\gamma 1}`$. As the material settles on the star, its radial velocity decreases and we have $`vv_{ff},\mathrm{\Omega }\mathrm{\Omega }_K`$. Then, from equation (2), it follows that the temperature of the gas is nearly virial. Other quantities are determined straightforwardly, so that we have
$$\rho r^{\frac{1}{\gamma 1}},c_s^2r^1,\mathrm{\Omega }\text{const.},vr^{\frac{2\gamma 3}{\gamma 1}}.$$
(50)
The infall velocity decreases with radius if $`\gamma <1.5`$, and increases if $`\gamma >1.5`$. To highlight the difference between the standard ADAF and the settling ADAF, we have chosen $`\gamma =4/3`$ in the solutions shown in Fig. Self-Similar Hot Accretion Flow onto a Neutron Star.
Finally, the long-dashed curves in Fig. Self-Similar Hot Accretion Flow onto a Neutron Star correspond to a solution with $`s=0.3`$ for which we have increased the outer boundary value of $`T`$ by a factor of 10. We see that the solution in the interior is not sensitive to the outer boundary conditions (within a reasonable range, of course).
## 6 Summary and Discussion
In this paper, we have presented analytical and numerical solutions that describe a hot, viscous, two-temperature accretion flow onto a neutron star (NS). The results are relevant also for accretion onto other compact stars with a surface, e.g. white dwarfs. To our knowledge, this is the first study of viscous fluid dynamics for a hot flow around a NS.
The presence of a surface modifies the nature of the flow relative to the case of a black hole. We show that the accretion flow has an extended settling region in which the radial velocity $`v`$ is constant; $`v`$ is also small relative to the local free-fall velocity. The density in the settling region varies as $`\rho r^2`$, and the angular velocity has a Keplerian scaling, $`\mathrm{\Omega }=s\mathrm{\Omega }_Kr^{3/2}`$, with $`s`$ being a constant. Here, $`r`$ is the radius in Schwarzchild units, and the value of $`s`$ is set by the spin of the NS: $`s=\mathrm{\Omega }_{NS}/\mathrm{\Omega }_K(r_{NS})`$. At the inner edge of the settling region, there is a narrow boundary layer in which the velocity falls extremely rapidly and the density increases sharply to match the surface density of the NS.
The settling region consists of two distinct zones. In the inner zone, $`r10^{2.5}`$, the gas is two-temperature, with the proton temperature varying as $`T_pr^1`$ and the electron temperature varying as $`T_er^{1/2}`$. These scalings are derived assuming that electrons radiate primarily by free-free emission and that energy transfer from protons to electrons occurs via Coulomb collisions. We have derived a completely general analytical self-similar solution for this region which agrees well with numerical results.
In the outer zone of the settling region, $`r10^{2.5}`$, the gas is one-temperature, $`T_pT_er^1`$; here again, we derive an analytical self-similar solution which agrees reasonably well with numerical results.
The most surprising feature of the settling region is that nearly all the gas properties are independent of the mass accretion rate $`\dot{m}`$; only the radial velocity shows a dependence: $`v\dot{m}`$. Since the density and temperature are independent of $`\dot{m}`$, the luminosity is also independent of $`\dot{m}`$. Indeed, the settling solution is valid—with a finite luminosity—even in the limit when $`\dot{m}0`$. Clearly, the luminosity does not originate from the gravitational release of energy as mass accretes onto the NS.
The magnitude of the luminosity is very sensitive to the spin parameter $`s`$ of the NS, varying as the fourth power of this quantity, equation (28). For $`s0.1`$, as appropriate for the millisecond X-ray pulsar, SAX J1808.4-3658, the model predicts an X-ray (few keV) luminosity $`\nu L_\nu 10^{32}\mathrm{erg}\mathrm{s}^1`$3.2). This estimate does not include the contribution from Comptonization in the inner regions of the flow, which might increase the X-ray luminosity by an undetermined amount. We note that quiescent X-ray luminosities of soft X-ray transients, including SAX J1808.4-3658, are generally in the range $`\nu L_\nu 10^{33}\mathrm{erg}\mathrm{s}^1`$ (Narayan, Garcia & McClintock 1997; Asai et al. 1999; Menou et al. 1999; Stella et al. 2000).
The angular momentum flux $`\dot{J}`$ in the settling solution is dominated by the viscous transport term rather than the advection term $`\dot{M}\mathrm{\Omega }_KR^2`$. Consequently, $`\dot{J}`$ is negative, i.e. the angular momentum flux is oriented outward, and the accretion flow spins down the star. The analytical solution predicts that spin-down occurs for all values of the spin parameter $`s`$ of the central star. The numerical solutions by and large confirm this; for plausible parameters, spin-up is seen only for extremely small values of the spin parameter, $`s<0.005`$ (cf Fig. Self-Similar Hot Accretion Flow onto a Neutron Star). (The exact value of $`s`$ at which $`\dot{J}`$ changes sign depends on $`\alpha `$ and $`\dot{m}`$, but is relatively independent of the position of the outer edge of the flow.) This behavior is very different from the case of a thin accretion disk (Popham & Narayan 1991; Paczyński 1991), where one finds that the star is spun-up for nearly all values of $`s`$, and spin-down occurs only for $`s`$ close to unity (break-up limit).
Another surprising feature of the settling solution is that $`\dot{J}`$, like nearly all other quantities, is independent of $`\dot{m}`$. Indeed, the settling zone behaves like a stationary zone (since $`v`$ is very small), and essentially acts like a conventional “brake,” slowing down the star by viscosity. The brake can operates even if $`\dot{m}0`$, so long as there is a static atmosphere of the self-similar form and there is a sink for the angular momentum, say an external medium, at large $`r`$. Furthermore, the luminosity of the settling flow is almost entirely from the energy released by the viscous braking action. That is, the luminosity ultimately is fed by the loss of rotational kinetic energy of the star, and not by gravity.
This result has an interesting consequence. In accretion flows around black holes, one defines an efficiency factor $`\eta `$ by comparing the accretion luminosity $`L_{acc}`$ to the rest mass energy of the accreting gas, $`\eta L_{acc}/\dot{M}c^2`$. It is well-known that $`\eta =0.06`$ for a thin accretion disk around a Schwarzchild black hole and that the value increases to $`\eta =0.42`$ for a maximally rotating Kerr hole. A number of interesting ideas have been discussed in the literature for increasing the efficiency of an accretion flow around a black hole; these involve tapping the rotational energy of the black hole using magnetic fields or viscosity (Blandford & Znajek 1977; Krolik 1999; Gammie 1999). The general relativistic dragging of inertial frames by the spinning hole plays an important role in the mechanism of (Blandford & Znajek 1977).
For the hot settling solution described in this paper, the luminosity is almost entirely from the rotational energy of the star. Since $`L_{acc}`$ is independent of $`\dot{m}`$, the efficiency scales as $`\eta \dot{m}^1`$ and $`\eta \mathrm{}`$ as $`\dot{m}0`$. Thus, it would appear that accretion flows can tap the rotation energy of a star with a surface more easily than the energy of a spinning black hole. Sibgatullin & Sunyaev (2000) showed that the extraction of rotational energy of a NS may result in very high boundary layer efficiencies, up to $`\eta 0.67`$, in the counter-rotating NS–thin disk systems, as well. As in the case of the black hole, the energy extraction works best when the star is spinning rapidly: $`\eta s^4`$. Interestingly, the energy extraction is not a general relativistic effect — our analytical solution is based entirely on Newtonian physics.
Yi, Wheeler & Vishniac (1997) and Yi & Wheeler (1998) recently suggested that the sudden torque-reversal events seen in some accretion-powered pulsars may be due to the accretion flow switching between a Keplerian thin disk and a hot a sub-Keplerian state akin to an ADAF. Our work lends support to this suggestion. We find that the torque does reverse in sign between a thin accretion disk and a hot settling flow for almost any reasonable stellar spin parameter. We also find that the magnitude of the spin-down torque exerted by the settling flow is comparable to the measured value of the torque in torque-reversing pulsars (§3.1).
It is worth emphasizing that the settling solution described here is quite distinct from the self-similar ADAF solutions derived for black hole accretion (Narayan & Yi 1994, 1995a; Honma 1996; Kato & Nakamura 1998; Blandford & Begelman 1999; Manmoto et al. 2000). All the black hole solutions described in the literature have density varying relatively mildly with radius: $`\rho r^{3/2}r^{1/2}`$. Our settling solution has $`\rho r^2`$. Also, the black hole solutions are advection-dominated, whereas our solution radiates essentially all the energy it generates through viscous dissipation. There are also, as we now discuss, significant differences in the sign of the Bernoulli parameter and in the stability to convection.
Narayan & Yi (1994, 1995a) showed that their self-similar ADAF solution has a positive Bernoulli parameter so long as the adiabatic index $`\gamma `$ of the gas is less than $`5/3`$. They argued on the basis of this that ADAFs are likely to have strong outflows and winds but see Abramowicz et al. 2000. Such strong outflows were confirmed with numerical simulations by Igumenshchev & Abramowicz (1999, 2000); they found outflows for large values of the viscosity parameter: $`\alpha 1`$. Blandford & Begelman (1999) developed a self-similar model with inflow and outflow (the ADIOS model). For the settling solution described in this paper, the Bernoulli parameter is positive only if $`\gamma `$ is less than $`\gamma _{crit}`$, where $`\gamma _{crit}=1.5`$ for a slowly-spinning star and is smaller than 1.5 for a rapidly-spinning star \[cf equation (40)\]. Since the hot ionized two-temperature gas in the flow is likely to have $`\gamma `$ close to $`5/3`$ at most radii, we expect the Bernoulli parameter to be generally negative. Therefore, we do not expect a strong outflow. Of course, this conclusion assumes that we do not have dynamically important magnetic fields in the flow.
Narayan & Yi (1994, 1995a) also showed that their ADAFs are convectively unstable for a wide range of conditions. The convective instability has been seen in numerical simulations in which the viscosity parameter is assigned a low value: $`\alpha 0.1`$ (Igumenshchev, Chen, & Abramowicz 1996; Igumenshchev & Abramowicz 1999, 2000; Stone, Pringle, & Begelman 1999; Igumenshchev, Abramowicz, & Narayan 2000). Self-similar solutions for convection-dominated accretion flows (CDAFs) have been derived by Narayan et al. (2000) and Quataert & Gruzinov (2000). For the settling solution described in this paper, we find that the gas is convectively unstable only for the same low values of $`\gamma `$ for which the Bernoulli parameter is positive. (Narayan et al. 2000 showed that, in general, for self-similar flows the criterion for the Bernoulli parameter to be positive is the same as the criterion for the flow to be convectively unstable, cf. §3.6.) Thus, we do not expect hot settling flows around NSs to be convectively unstable, or to have a distinct CDAF mode of settling.
From numerical experiments we have discovered the interesting property that the settling flow can continue well inside the last stable orbit (down to at least $`r_{NS}2`$) and yet remain subsonic at all radii. This is despite the fact that the numerical models employ a pseudo-Newtonian potential which mimic the last stable orbit for test particles. We find that the structure of the flow is qualitatively similar for flows with $`r_{NS}<3`$ and $`r_{NS}>3`$. This is very different from black hole ADAFs which become supersonic close to the central object.
We have also shown (§5) that the settling solution and the ADAF are not two physically distinct solutions, but are related to each other. As the neutron star spin is decreased, we find that the settling flow smoothly transforms to an ADAF-type solution. The transformation proceeds over the spin range $`s0.010.1`$.
The settling solution is hot but cooling-dominated. It is thus most closely related to the two-temperature solution discovered by Shapiro, Lightman & Eardley (1976). The SLE solution is known to be thermally unstable (Piran 1978; Wandel & Liang 199?; Narayan & Yi 1995b), so one may wonder about the thermal stability of our settling solution. We defer discussion of this important topic to a future paper.
Could there be non-settling solutions around NSs, and could such flows be more analogous to the black hole ADAF, ADIOS, and CDAF solutions? Indeed this is possible if the NS radius is small enough. We could imagine, for instance, a standard black hole-like flow around a NS, which makes a sonic transition to a supersonic state, and then stops suddenly at a standing shock at the surface of the NS. Such a solution would be dynamically consistent, and, except for the shock, would be very similar to a black hole flow. However, for such a solution to exist, one requires the radius of the NS to be smaller than the sonic radius of the flow. The latter radius is estimated to be in the range $`r_{sonic}25`$, depending on the value of $`\alpha `$ (Narayan, Kato & Honma 1997; Popham & Gammie 1998), whereas NS radii are in the range $`r_{NS}24`$ for typical NS equations of state. Thus, for some choices of $`\alpha `$ and some equations of state, $`r_{NS}`$ could be smaller than $`r_{sonic}`$. In such cases, we could have four different hot solutions around a NS: (i) we could have the self-similar settling solution of the kind presented in this paper (as we discussed in §4, we find subsonic settling solutions for any choice of $`R_{NS}`$ in the range $`24R_S`$); or (ii) we could have a Narayan, Kato & Honma (1997)–like ADAF solution with a shock at the NS surface; or (iii) we could have a Blandford & Begelman (1999)–like ADIOS solution with a shock at the NS surface; or (iv) we could have a Narayan et al. (2000) and Quataert & Gruzinov (2000)–like CDAF, again with a shock at the NS surface. In the latter three cases, we expect the NS to be spun-up rather than spun-down by the accretion flow, and we also expect the radiative efficiency to be close to the standard value for a NS, namely $`\eta 0.10.5`$. The spectrum of the radiation is also likely to be very different in the four models. This may provide a way to distinguish which if any of these possibilities is found in nature.
In addition to the above possibilities, yet other flow configurations may be possible when we allow for the multi-dimensional nature of the flow. These could be explored with numerical hydrodynamics simulations.
This work was supported in part by grants PHY 9507695 and AST 9820686 from the National Science Foundation.
Fig. Self-Similar Hot Accretion Flow onto a Neutron Star
Fig.Self-Similar Hot Accretion Flow onto a Neutron Star
Fig. Self-Similar Hot Accretion Flow onto a Neutron Star
Fig.Self-Similar Hot Accretion Flow onto a Neutron Star
Fig. Self-Similar Hot Accretion Flow onto a Neutron Star
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# Dynamical Cosmological Constant From A Very Recent Phase Transition
## Abstract
Observation indicates that the expansion of the Universe is accelerating and favours a dynamical cosmological constant, $`\mathrm{\Lambda }(t)`$. We consider the possibility that this is due to a scalar field which has undergone a very recent phase transition. We study a simple class of model, corresponding to a $`\varphi ^4`$ potential with a time-dependent mass squared term. For the models considered the phase transition occurs at a red shift $`z1.2`$. The evolution of the equation of state $`\omega _\varphi `$ and energy density $`\rho _\varphi `$ with time is distinct from existing dynamical $`\mathrm{\Lambda }`$ models based on slowly rolling fields, with $`\omega _\varphi `$ and $`\rho _\varphi `$ rapidly changing in a characteristic way following the transition. The $`\varphi `$ energy density is composed of a time-dependent vacuum energy and coherently oscillating condensate component with a negative pressure. The condensate component will typically collapse to form non-topological soliton lumps, ’$`\varphi `$-axitons’, which smoothly populate the Universe.
<sup>1</sup>mcdonald@physics.gla.ac.uk
Perhaps the most remarkable cosmological observation of recent times is the accelerating expansion of the Universe . This requires the existence of an energy density which has a negative pressure and which is smooth on the 10Mpc scales relevant to dynamical estimates of the density of conventional dark matter. Other evidence for this comes from estimates of the age of the Universe, the Hubble constant, the baryon fraction in clusters, the galactic power spectrum and CMB measurements . The simplest explaination for a smooth energy density with a negative pressure is a time-independent cosmological constant, $`\mathrm{\Lambda }`$. However, there are a number of problems with a fixed cosmological constant. The first is that it is difficult to understand why the cosmological constant should just be dominating the energy density at the present epoch . This has led to two different views. One is that the cosmological constant is due to the evolution of a scalar field (’tracking solution’ ; see also ) whose energy density becomes significant recently due to dynamical effects. In general this approach has difficulties with nucleosynthesis and the present equation of state . (A particularly promising version which may overcome these problems is given in .) The other view is that the cosmological constant has become dominant in the present epoch due to anthropic selection (AS) . However, even if we assume that AS is responsible for the dominance of the cosmological constant, there are still problems. It is difficult to understand why a fixed $`\mathrm{\Lambda }`$ is so small compared to the mass scales of particle physics ($`\mathrm{\Lambda }10^{120}\mathrm{M}_{\mathrm{Pl}}^4`$). This has led to the suggestion that the absolute minimum of the vacuum energy should be exactly zero and that the smooth energy density is due to the evolution of a scalar field towards the minimum, resulting in a dynamical cosmological constant, $`\mathrm{\Lambda }(t)`$ . (An interesting alternatve based on quantum spinodial fluctuations has been suggested in .) There is also observational evidence in support of a specifically dynamical cosmological constant . The amplitude of the COBE-normalized galaxy clustering power spectrum is too large in the case of a fixed $`\mathrm{\Lambda }`$. However, if the effective cosmological constant is decreasing with time, then the amplitude of the power spectrum is reduced, alleviating the problem .
There have been a number of suggestions regarding the nature of the dynamical cosmological constant. Most popular are models based on pseudo-Nambu-Goldstone bosons (PNGBs) , exponential potentials and inverse-power law potentials . All of these produce characteristic time-dependent energy densities $`\rho _\varphi `$ and pressures $`p_\varphi (\omega _\varphi \rho _\varphi )`$ which come to dominate at recent red-shifts. For the PNGB models, the pressure slowly tends from negative values towards zero. For tracking models based on exponential and inverse-power law potentials, on the other hand, the pressure is evolving from a value typically positive and close to zero (corresponding to matter tracking) towards a negative value. This should allow the models to be distinguished by precision CMB angular spectra measurements . All of these models are based on very light scalar fields ($`m_\varphi ^<{}_{}{}^{}H_{o}^{}`$, where $`H_o`$ is the present value of the expansion rate) which are slowly rolling at the present time, $`|\dot{\varphi }/\varphi |_{}^<H_o`$.
In the present paper we wish to introduce an alternative model for a dynamical cosmological constant. This is based on the idea that the negative pressure energy density is associated with a conventional metastable false vacuum. This is perhaps the simplest form of negative pressure energy density in particle physics models. However, in order to have a dynamical cosmological constant we require that a phase transition from the metastable phase has occured recently. To achieve this we will consider the scalar field $`\varphi `$ to have a time dependent mass squared term which has recently become negative. We will refer to this scheme as the ’very recent phase transition’ (VRPT) scenario for a dynamical cosmological constant. The VRPT scenario is quite distinct from previous models based on slowly rolling scalar fields. In particular, there will be no need for the mass scale of the scalar to be extremely small, and the pressure and energy density will be rapidly evolving in a characteristic way at recent times following the phase transition. This should allow the VRPT scenario to be distinguished from the others by precision CMB measurements.
We will consider the usual spontaneous symmetry breaking potential for a real scalar field with $`\varphi \varphi `$ symmetry<sup>1</sup><sup>1</sup>1This can be generalized to a complex scalar. For a real scalar field domain walls may form during the VRPT, whilst for a complex field global strings may form.,
$$V(\varphi )=\frac{\mu ^2(t)}{2}\varphi ^2+\frac{\lambda }{4}\varphi ^4+\mathrm{\Lambda };\mathrm{\Lambda }=\frac{\mu _o^4}{4\lambda },$$
(1)
where
$$\mu ^2(t)=\mu _o^2\left(1\left(\frac{a_c}{a}\right)^n\right)$$
(2)
and $`a_c`$ is the scale factor at the time of the transition. In general the VRPT scenario requires an additional time-dependent $`\varphi ^2`$ term, which we will refer to as a ’stabilizing interaction’. We will discuss some possible sources for the stabilizing interaction later, but for now we simply model it phenomenologically. We will see that such time-dependent mass squared terms with integer $`n`$ can arise naturally in plausible models.
Typically the mass scale of the potential will be very large compared with $`H_o`$. Thus following the phase transition $`\varphi `$ will be coherently oscillating about the relatively slowly evolving time-dependent minimum of its potential. In order to discuss the time evolution of the $`\varphi `$ energy density and equation of state, we follow the discussion of , based on averaging over the rapid oscillations of the scalar about the minimum of its potential.
The $`\varphi `$ equation of motion is given by
$$\ddot{\varphi }+3H\dot{\varphi }=\frac{V}{\varphi }.$$
(3)
This may be rewritten as
$$\frac{}{t}\left(\frac{\dot{\varphi }^2}{2}+V\right)=3H\dot{\varphi }^2,$$
(4)
where the partial derivative is with respect to constant $`\mu ^2(t)`$.<sup>2</sup><sup>2</sup>2In practice we evolve the energy density and pressure of $`\varphi `$ by incrementing the energy density with $`\mu (t)`$ held constant and then incrementing $`\mu (t)`$ and calculating the pressure. Therefore we do not use the total derivative with respect to $`t`$ given in for an explicitly time-dependent potential. Taking the time average over an oscillation cycle, we obtain
$$\frac{\rho _\varphi }{t}=3H\gamma \rho _\varphi ,$$
(5)
where
$$\gamma =\frac{2_c(1V/V_{max})^{1/2}𝑑\varphi }{_c(1V/V_{max})^{1/2}𝑑\varphi }.$$
(6)
$`_c`$ denotes integration over one oscillation cycle and $`V_{max}\rho _\varphi `$ is the maximum $`\varphi `$ energy density during an oscillation cycle. $`\gamma `$ corresponds to the time average of $`\dot{\varphi }^2/\rho _\varphi `$ over an oscillation cycle. The equation of state of the $`\varphi `$ energy density is then given by
$$\omega _\varphi p_\varphi /\rho _\varphi =\gamma 1.$$
(7)
We have calculated the time evolution of the energy density and equation of state for different values of $`n`$ as a function of red-shift, $`z`$. The parameters of the potential are chosen such that we have at present $`\mathrm{\Omega }_\varphi 0.7`$ and $`\mathrm{\Omega }_m0.3`$ (where $`\mathrm{\Omega }_m`$ is the density of conventional clustered dark matter).
The numerical solutions for $`\rho _\varphi `$ and $`\omega _\varphi `$ as a function of $`z`$ depend on $`\mu _o^2`$ and $`\lambda `$ only through the ratio $`\mu _o^4/\lambda `$; we have calculated the evolution for the case $`\lambda =1`$. Because of this scaling property, for a given value of $`\omega _\varphi `$ today, $`\omega _{\varphi o}`$, the evolution of the $`\varphi `$ equation of state and energy density is completely fixed by the value of $`n`$ in the stabilizing interaction.
In Figure 1 we show the evolution of the equation of state $`\omega _\varphi `$ as a function of red-shift for $`n=2,3,4`$. In Table 1 we give the parameters of the VRPT ($`\omega _{\varphi o}`$, $`\mu _o`$ (for $`\lambda =1`$) and $`z_c`$, the red-shift at which the transition occurs) for the cases $`n=2`$ and $`n=4`$.
For $`n2`$ and $`\omega _{\varphi o}0.6`$ we find that the phase transition occurs at $`z_c1.2`$. (We consider $`\omega _{\varphi o}0.6`$, in keeping with observational limits for the case of a fixed $`\omega _\varphi `$ . Since in our case $`\omega _\varphi `$ is decreasing with $`z`$, this limit should be conservative.) From Table 1 we see that the mass of the scalar is $`{}_{}{}^{<}\mathrm{\hspace{0.33em}10}_{}^{3}\mathrm{eV}`$ for $`\lambda _{}^<\mathrm{\hspace{0.33em}1}`$. In Figure 2 we show the evolution of the energy density together with the matter energy density $`\mathrm{\Omega }_m`$ for the cases $`n=2`$ and $`n=4`$, where we have normalized the energy density by taking the ratio to the present critical density $`\rho _c`$. Both $`\omega _\varphi `$ and $`\rho _\varphi `$ rapidly change at recent red-shifts, the more so for larger values of $`n`$.
In addition, we have considered the effect of the VRPT on the age of the Universe, given by
$$t_U=\frac{2}{3}H_o^1f_U;f_U=\frac{3}{2}\frac{da}{a}(\mathrm{\Omega }_\varphi (t)+\mathrm{\Omega }_m(t))^{1/2}.$$
(8)
The age of globular clusters requires that $`f_U=1.5\pm 0.3`$ . For a fixed cosmological constant and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`f_U=1.45`$. For the VRPT $`f_U`$ is generally smaller, but not significantly so. The largest deviation in the examples considered corresponds to $`n=2`$ and $`\omega _{\varphi o}=0.6`$, for which $`f_U=1.38`$.
Table 1. VRPT parameters.
| n | $`\omega _{\varphi o}`$ | $`\mu _o`$ | $`z_c`$ |
| --- | --- | --- | --- |
| 2 | $`0.6`$ | $`3.45\times 10^3\mathrm{eV}`$ | $`1.20`$ |
| | $`0.8`$ | $`3.15\times 10^3\mathrm{eV}`$ | $`0.47`$ |
| 3 | $`0.6`$ | $`3.27\times 10^3\mathrm{eV}`$ | $`0.57`$ |
| | $`0.8`$ | $`3.13\times 10^3\mathrm{eV}`$ | $`0.28`$ |
| 4 | $`0.6`$ | $`3.21\times 10^3\mathrm{eV}`$ | $`0.38`$ |
| | $`0.8`$ | $`3.10\times 10^3\mathrm{eV}`$ | $`0.20`$ |
After the phase transition has occured<sup>3</sup><sup>3</sup>3We are assuming that $`\varphi `$ is out of thermal equilibrium, so that the transition is purely dynamical in nature., the $`\varphi `$ energy density will be composed of a time dependent vacuum energy $`\rho _{vac}`$ and a $`\varphi `$ condensate component $`\rho _{osc}`$ corresponding to coherent oscillations about the time dependent minimum of the potential. In Figure 3 we show the equation of state and energy density of the $`\varphi `$ oscillations as a function of $`z`$. We see that the equation of state rapidly tends towards the value $`\omega _{osc}=0`$ corresponding to effectively $`\varphi ^2`$ oscillations about the minimum. Nevertheless, the equation of state and so the pressure is significantly negative throughout. In Figure 4 we show the equation of state and energy density associated with the time dependent vacuum energy, $`\rho _{vac}V(\varphi _{min}(t))`$. We see that the equation of state does not significantly deviate from the value expected for a constant vacuum energy density, $`\omega _\mathrm{\Lambda }=1`$.
The negative pressure of the $`\varphi `$ oscillations implies that the $`\varphi `$ condensate is unstable with respect to spatial perturbations of $`\varphi `$ . Such perturbations may be expected to exist, coming, for example, from thermal or inflationary quantum fluctuations. In the linear approximation the spatial perturbations evolve as
$$\delta \ddot{\varphi }_𝐤=|\omega _\varphi |𝐤^2\delta \varphi _𝐤$$
(9)
(where $`\omega _\varphi <0`$ is assumed). Therefore,
$$\delta \varphi _𝐤=\mathrm{exp}(|\omega _\varphi |^{1/2}|𝐤|t)\delta \varphi _{𝐤o}.$$
(10)
This is true so long as the wavenumber of the perturbation satisfies $`|𝐤|^2{}_{}{}^{<}\mathrm{\hspace{0.33em}4}|\omega _\varphi |m_\varphi ^2`$ , otherwise the positive pressure associated with gradient energy in the perturbation will overcome the negative pressure responsible for the growth of the perturbation. (For this reason, negative pressure effects play no role in the dynamics of $`\mathrm{\Lambda }(t)`$ models based on slowly rolling fields, since their small mass, $`m_\varphi <H_o`$, implies that perturbations on sub-horizon scales cannot grow.) The perturbations in the condensate will grow exponentially until they become non-linear. The condensate will then fragment into non-topological soliton lumps which we will refer to as ’$`\varphi `$-axitons’ (following the existence of similar objects in axion cosmology and Affleck-Dine baryogenesis ). The radius of the $`\varphi `$-axitons will be determined by the first perturbation mode to go non-linear, $`r_\varphi (|\omega _\varphi |^{1/2}m_\varphi )^1`$ (assuming the exponential factor to be dominant in determining non-linearity). This will occur in a time $`\delta t(|\omega _\varphi |m_\varphi )^1\mathrm{log}(\varphi /\delta \varphi _𝐤)H_o^1`$.
Thus shortly after the VRPT the Universe will typically be filled with $`\varphi `$-axitons. So long as $`r_\varphi 10`$Mpc, which is true if $`m_\varphi 10^{39}\mathrm{GeV}`$, the $`\varphi `$-axiton density will initially act as a smooth component of pressureless matter as far as determinations of the dark matter density are concerned.
We need to check that subsequent infall into galactic halos does not result in the $`\varphi `$-axiton density clustering on 10Mpc scales and so no longer being effectively smooth. We do this via a simple Newtonian argument. We consider the mean distance between galaxies to be $`R_{gal}10`$Mpc, with $`\varphi `$-axitons being smoothly distributed initially. The time scale for the $`\varphi `$-axitons to fall a distance $`R_{gal}`$ due to the attraction of a galaxy of mass $`M_{gal}`$ is then
$$t_{infall}\frac{R_{gal}^{3/2}}{\sqrt{GM_{gal}}}.$$
(11)
Using $`M_{gal}10^{11}M_{}`$ for the mean galactic mass, this gives $`t_{infall}1.6\times 10^{12}`$yr, which is much longer than the age of the Universe, $`t_U10^{10}`$yr. Thus infall will not significantly alter the smoothness of the $`\varphi `$-axiton distribution. (The same argument holds for the $`\varphi `$ particles in a homogeneous condensate.) Although we have estimated this at the present time, it holds for earlier times also, since $`t_U(H^1)`$ and $`t_{infall}(R_{gal}^{3/2})`$ are both proportional to $`a^{3/2}`$.
Therefore in the presence of spatial perturbations, the energy being fed into the coherent $`\varphi `$ oscillations will be converted into a smooth density of pressureless $`\varphi `$-axitons. This will slightly alter the evolution of the total $`\varphi `$ equation of state and energy density from the case where there are coherent $`\varphi `$ oscillations with negative pressure. In Figure 5 we show the effect of replacing the $`\varphi `$ condensate with pressureless $`\varphi `$-axiton matter. The effect is a small alteration of the total $`\varphi `$ equation of state and energy density as a function of $`z`$.
After the phase transition, the $`\varphi `$ energy density consists of a time-dependent vacuum energy
$$\rho _{vac}=\frac{\mu _o^2\mu ^2(t)}{4\lambda },$$
(12)
with $`\omega _{vac}=1`$ to a very good approximation, and an energy density either in the form of a coherently oscillating $`\varphi `$ field with a negative pressure or, more likely, in the form of pressureless $`\varphi `$-axitons. However, the total $`\omega _\varphi `$ is insensitive to the pressure in the condensate component, as seen from Figure 5. This means that $`\omega _\varphi >1`$ following the VRPT is simply due to the dilution of $`\omega _\varphi `$ by the low pressure component of smooth $`\varphi `$ condensate or $`\varphi `$-axiton matter,
$$\omega _\varphi =\frac{\omega _{osc}\rho _{osc}+\omega _{vac}\rho _{vac}}{\rho _{osc}+\rho _{vac}}\frac{\rho _{vac}}{\rho _{osc}+\rho _{vac}}.$$
(13)
So far we have introduced a simple time-dependent mass squared term (’stabilizing interaction’) in order to trigger the phase transition and so produce a dynamical cosmological constant. We now consider some ways in which this term could be generated in particle physics models. One possibility is that the $`\varphi `$ field might couple to a light field which has a thermal distribution, such as the photons or neutrinos. For example, one could consider a coupling of the form $`\frac{\varphi ^2}{M^2}_{ke}`$, where $`_{ke}=\frac{1}{4}F^{\mu \nu }F_{\mu \nu },\overline{\psi }_\mu \gamma ^\mu \psi `$. On taking the average over thermal fluctuations of the fields in $`_{ke}`$ we obtain an effective $`\varphi ^2`$ term $`\frac{T^4}{M^2}\varphi ^2`$. With $`Ta^1`$ this results in a stabilizing interaction with $`n=4`$. Alternatively we might consider a very light additional scalar field $`\chi `$ with a thermal distribution (and a small enough energy density so as to avoid nucleosynthesis constraints), coupling to $`\varphi `$ via a term $`\chi ^2\varphi ^2`$. On averaging over thermal fluctuations of the $`\chi `$ field, this will give an effective $`\varphi ^2`$ term $`T^2\varphi ^2`$, corresponding to a stabilizing interaction with $`n=2`$. In both of these interactions we have effectively massless fields coupling to $`\varphi `$, so that the time dependence of the stablilzing interaction is due to the red-shifting of the energy of the thermal fluctuations of the light fields. An alternative is to consider the case where $`\chi `$ is massive and is coherently oscillating in an effectively $`\chi ^2`$ potential about $`\chi =0`$. We assume that the $`\chi `$-matter density is negligible compared with the conventional dark matter density. In this case the $`\chi `$ oscillation amplitude will be proportional to $`a^{3/2}`$, resulting in a stabilizing interaction with $`n=3`$. These are just a few examples; one could also consider, for example, $`\chi `$ to be a slowly rolling field which triggers the phase transition, analogous to what happens in hybrid inflation models .
There has recently been some discussion of the possibility of distinguishing between models with different time-dependent equations of state $`\omega (z)`$. In it was suggested that in order to distinguish between models with $`\omega (z)`$ and $`\omega =constant`$ we would require an accuracy of less than 1$`\%`$ in determining the luminosity distance as a function of red-shift, $`d_L(z)`$, whereas 1$`\%`$ is regarded as an optimistic estimate of the accuracy of future experiements . However, in , it is suggested that with an appropriate fit to $`\omega (z)`$, $`\omega (z)=\omega _0+\omega _1z+\mathrm{}`$, it is possible to distinguish between different models using the datasets expected from the SNAP satellite . In particular, they show that a toy model with $`\omega _0=0.6`$ and $`\omega _1=0.8`$ can be clearly distinguished from constant $`\omega `$ models. Comparing with Figure 1, we find that for $`n=2,3,4`$ the expansion parameters $`(\omega _o,\omega _1)`$ are given by (-0.6,-0.33), (-0.6,-0.70) and (-0.6,-1.05) respectively. So at least for $`n=3`$ and $`n=4`$ the VRPT scenario should be clearly distinguishible by SNAP, and possibly for $`n=2`$ also, although this is not directly apparent from the results of and requires further analysis.
In conclusion, we have introduced an alternative model for a dynamical cosmological constant, based on the idea that a scalar field underwent a phase transition from a metastable phase at a very recent epoch, $`z1.2`$. This VRPT scenario results in a characteristic evolution of the equation of state and energy density which is quite distinct from the case of models based on slowly rolling fields, with the pressure rapidly rising from $`\omega _\varphi =1`$ and $`\rho _\varphi `$ rapidly decreasing at very recent times. The solutions for $`\omega _\varphi (z)`$ and $`\rho _\varphi (z)`$ are uniquely determined by the present value of $`\omega _\varphi `$ and the form of the time dependent mass squared term in the potential. Following the phase transition, the $`\varphi `$ energy density will consist of a time dependent vacuum energy and a negative pressure $`\varphi `$ condensate which typically fragments to a smooth pressureless density of non-topological solitons, ’$`\varphi `$-axitons’.
As with most other dynamical $`\mathrm{\Lambda }`$ models, it is implicitly assumed that the reason for the recent dominance of the $`\varphi `$ energy density and the recent occurance of the phase transition is connected with anthropic selection. In general, dynamical $`\mathrm{\Lambda }`$ models require two conditions; that the energy density in the scalar field has recently become dominant and that the scalar field energy density is varying significantly on time scales of the order of $`H_o^1`$. Therefore two tunings are generally required. In this regard the VRPT scenario is on the same footing as other dynamical $`\mathrm{\Lambda }`$ models. Ultimately the question of which dynamical $`\mathrm{\Lambda }`$ model is correct is a matter to be decided by observations. A future goal will therefore be to understand the detailed predictions of this class of dynamical $`\mathrm{\Lambda }`$ model for observable quantities; high-z supernova, quasar lensing statistics, the galaxy clustering power spectrum and in particular the angular CMB spectrum. We expect that these will be clearly distinguishable from other dynamical $`\mathrm{\Lambda }`$ models, and indeed we have shown by comparison with recent analyses that at least some VRPT models are likely to be distinguishable by SNAP. In addition, the idea of a very recent phase transition leads to other interesting issues. One is whether the $`\varphi `$-axiton density could be observationally or experimentally distinguished from conventional cold dark matter. This will depend on the $`\varphi `$ mass and on how strongly the $`\varphi `$ field interacts with ordinary matter. In addition, it is possible that the VRPT could result in the very recent formation of domain walls or global strings, depending on the symmetry associated with the potential and whether $`\varphi `$ is real or complex. It would be interesting to consider whether there are any observable effects associated with the very recent creation of topological defects. With the possibility of recently formed topological and non-topological solitons, the phenomenology of the VRPT scenario for the dynamical cosmological constant may be surprisingly rich.
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# Angle resolved photoemission spectroscopy of Sr2CuO2Cl2 – a revisit
## I Introduction
Extensive effort continues to be expended to understand the electronic properties of the cuprate high $`T_c`$ superconducting compounds. Apart from many differences, the cuprate superconducting materials have important properties in common. Namely they are all layered compounds possessing CuO<sub>2</sub>-planes, where the origin of the superconducting behavior is to be found. The electronic structure of the states with the lowest binding energy is satisfactorily understandable only in models which take into account correlation effects. ARPES has proven to be a powerful tool to investigate the low binding-energy electronic structure of the cuprates. Important examples for the success of this method are the mapping of the Fermi surface and the determination of an anisotropy of the superconducting gap consistent with a d-wave order parameter .
In the undoped regime the cuprates are two-dimensional antiferromagnetic insulators with CuO<sub>2</sub>-planes resulting in a correlation-gapped half-filled Cu3$`d_{x^2y^2}`$-O2$`p_x`$/O2$`p_y`$ antibonding band of lowest binding-energy. Upon doping with holes, the states with lowest binding energy evolve towards the chemical potential, and the cuprate becomes metallic or superconducting. To clarify the main principles of the evolution of the lowest binding-energy states with doping by applying model Hamiltonians, one has to go step by step from simple situations or limiting cases to more complicated scenarios. A natural starting point, then, is to investigate the propagation of a single hole in a CuO<sub>2</sub>-plane by studying the undoped parent compounds of cuprate superconductors by photoemission, where one electron is removed from the CuO<sub>2</sub>-plane. The states thus investigated can be termed the first electron-removal states.
A special class of parent compounds are the oxyhalides such as Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>, Ca<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>, Sr<sub>2</sub>CuO<sub>2</sub>F<sub>2</sub> (one layer compounds), and Ca<sub>3</sub>Cu<sub>2</sub>O<sub>4</sub>Cl<sub>2</sub> (double layered compound) , all of which contain an apical halogen atom, rather than an oxygen. It is now well established that the apical oxygen which is contained in most of the cuprate superconductors is not necessary for high $`T_c`$ superconductivity as Hiroi et al. showed that Ca<sub>2-x</sub>Na<sub>x</sub>CuO<sub>2</sub>Cl<sub>2</sub> is superconducting ($`T_c=26`$K). Out of the class of cuprate parent compounds without apical oxygen, the layer compound Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> is remarkable due to its extremely stable stoichiometry. Despite considerable efforts so far it has not been possible to dope this substance chemically. Furthermore, the CuO<sub>2</sub>-planes are unbuckled and the crystal structure shows neither orthorhombic distortion nor a superstructure at least down to 10 K (Ref. ). Therefore, Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> can be seen as the best realization of a two-dimensional antiferromagnet at half-filling and it is in this limit an ideal test case for any Hamiltonian in the low-doping regime.
In theoretical studies, the lowest binding-energy ionization states of the cuprates were predicted to be essentially described by a singlet antibonding combination of a Cu3$`d_{x^2y^2}`$ orbital (containing the intrinsic Cu3$`d^9`$ hole) and a coherent combination of the four neighboring O2$`p_x`$/O2$`p_y`$ orbitals which can be thought of as containing the hole created in the photoemission process . This two-hole state is generally referred to as the Zhang-Rice singlet state . The singlet character of the first electron-removal states, at least in CuO, has been experimentally verified using spin-resolved resonant photoemission.
There are various possibilities to theoretically model the dynamics of this lowest binding-energy excitation. For instance, in the framework of a three-band Hubbard Hamiltonian ($`H_{3b}`$), the Zhang-Rice singlet state is a two-particle state and belongs to the $`A_{1g}`$ (totally symmetric) irreducible representation of the eigenstates of $`H_{3b}`$. In the one-band Hubbard Hamiltonian ($`H_{1b}`$), the $`t`$-$`J`$-model and its extensions, the Zhang-Rice singlet state is an one-particle state belonging to the $`B_{1g}`$ irreducible representation and has the same symmetry as a Cu3$`d_{x^2y^2}`$ orbital. Figure 1 shows a schematic representation of a Zhang-Rice singlet state, including the two mirror planes M<sub>1</sub> and M<sub>2</sub> which are perpendicular to the CuO<sub>2</sub>-plane and which are relevant for photoemission along along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) and $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$), respectively.
It is now generally accepted that the $`H_{3b}`$ is a good starting point to describe the low binding-energy electronic structure in the strongly correlated cuprate systems and consequently the $`H_{3b}`$ has been investigated in detail . Additionally, it has been shown that the $`H_{1b}`$ still carries enough information to describe the low-energy dynamics in the cuprates . In the next level of simplifying these models, the Hubbard Hamiltonian can be expressed in terms of new Fermionic operators. In the large-$`U`$-limit, this new Hamiltonian can be written in a $`t^2/U`$ series and gives an effective one-band Hamiltonian, the $`t`$-$`J`$-model. This has been shown to give the same results for both the $`H_{3b}`$ (Ref. ) and the $`H_{1b}`$ (Ref. ) approaches. Extensions of the $`t`$-$`J`$-model have been carried out mainly in two different ways. Firstly, two further parameters - $`t^{}`$ describing the diagonal hopping and $`t^{\prime \prime }`$ for the next nearest neighbor hopping - can be introduced in addition to the direct hopping $`t`$.
The fact that a $`t`$-mediated hopping creates strings of spin defects leads to a predicted bandwidth for the Zhang-Rice singlet state which is governed by $`t^2/UJ`$. Hopping events mediated by $`t^{}`$ and $`t^{\prime \prime }`$, however, take place on the same spin sub-lattice and therefore even taking small values of these parameters has a significant effect on the results the $`t`$-$`J`$-model gives for the $`𝐤`$-dependent spectral function of the Zhang-Rice singlet state. Secondly, the $`t`$-$`J`$-model has also been extended to take three-site hopping terms (proportional to $`J/4`$) into account , which appear in a natural way if one employs a more detailed derivation of an effective Hamiltonian for the low binding-energy excitations .
A different ansatz to understand the dynamics of holes in CuO<sub>2</sub>-planes was proposed by Laughlin . He argued that the hole created by photoemission decays into spin and charge degrees of freedom and that the photoemission experiment measures the dispersion of the spinon . However, in a gauge-field treatment of this model this decay of the quasiparticles is suppressed below the Néel-temperature, $`T_N`$, due to confinement . Another fundamental approach to calculate the properties of the Hubbard Hamiltonian and the $`t`$-$`J`$-models uses the fact that they share an approximate SO(5) symmetry .
The preceding discussion highlights the intense theoretical interest in the lowest binding-energy electron-removal states of the undoped cuprates. Although originally trailing a few years behind the theoretical work, experimental investigations of the first electron-removal state in undoped cuprates have been carried out using ARPES . Before going on to mention the status in the field to date, we first ask the question: what kind of information does an ARPES experiment provide us with?
Neglecting “extrinsic” effects such as scattering of the photoelectrons on their way to the surface and assuming the applicability of the sudden approximation, the photoemission intensity $`I(𝐤,E)`$, reads
$$I(𝐤,E)\underset{i,f}{}M_{i,f}^2S(𝐤,E)f(E)$$
(1)
where the sum runs over all final states and initial states, $`f(E)`$ denotes the Fermi function and $`S(𝐤,E)`$ the spectral function:
$$S(𝐤,E)=\frac{1}{\pi }\text{Im}ic_𝐤\frac{1}{EH+E_i}c_𝐤^{}i.$$
(2)
$`M_{i,f}=f𝐀𝐩i`$ is the matrix element to be taken between the initial state $`i`$ with energy $`E_i`$ and final states $`f`$ with energy $`E`$, $`𝐀𝐩`$ is the photoemission interaction operator and $`H`$ the Hamiltonian of the system. From an ARPES experiment one can gain mainly two kinds of information. Firstly, the spectral function $`S(𝐤,E)`$ gives direct information about the dispersion and the quasiparticle character of the states and can be compared directly to predictions of model Hamiltonians. All information carried by $`S(𝐤,E)`$ is expected to be independent of the excitation energy and geometry of the experiment. Secondly, the matrix element $`M_{i,f}`$ includes all information concerning the photoemission interaction. As such it is sensitive to the experimental geometry via symmetry selection rules as well as to the photon-energy.
The first ARPES experiments on Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> (Ref. ) were able to make the important observation that the Zhang-Rice singlet state bandwidth along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) is of the order of 2$`J`$, and that the Zhang-Rice singlet state approaches closest to the chemical potential at the ($`\pi /2`$,$`\pi /2`$) point. These observations fit the predictions of the $`t`$-$`J`$-model. However, these early experiments found very little dispersion along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$), which is in conflict with the predictions of the same model. Subsequent studies have confirmed the behavior along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$), but have also observed dispersion along the $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$)-direction in reciprocal space ($`𝐤`$-space).
Nevertheless, the situation as regards the exact dispersion relation along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$), and in particular concerning the minimum energy difference between the states along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) and those at the ($`\pi /2`$,$`\pi /2`$) point is unclear. This is more than a mere ARPES detail, as it is along the $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) direction that the predictions of the different models vary the most - thus making this direction in $`𝐤`$-space important for the quantitative comparison between theory and experiment.
Furthermore, the recent controversy surrounding the correct Fermi surface topology in the doped high temperature superconductor Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> (Ref. ) has illustrated that data sets recorded from the same system with different experimental conditions can be remarkably dissimilar. Consequently, both the photon-energy and exact polarization geometry used in an ARPES experiment are important parameters which, if they cannot be treated at a quantitative, microscopic level, should at least be thoroughly investigated on the experimental side.
In this paper, we address the electronic structure and dynamics of the lowest binding-energy electron-removal states in the “standard” undoped model cuprate Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> using ARPES. Following the description of the experimental status given above, our “re-visit” of this system concentrates on the following points:
a) a thorough characterization of the photon-energy dependence of the first electron-removal states at ($`\pi /2`$,$`\pi /2`$) and close to ($`\pi /2`$,$`\mathrm{\hspace{0.17em}0}`$) (the two points along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) and $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) for which the first electron-removal states have maximum spectral weight); b) a thorough characterization of the polarization dependence of the first electron-removal states along these high-symmetry directions in $`𝐤`$-space; c) the determination of the dispersion relation of the first electron-removal states along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) and $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$), with subsequent comparison of the results with both existing ARPES data and the predictions of an extended $`t`$-$`J`$-model and d) the determination of the $`𝐤`$-dependent evolution of the coherent and incoherent parts of the spectral weight of the first electron-removal state along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) and $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$), and comparison of these data with theoretical predictions.
## II Experiment
Experiments were performed at the beamlines F$`2.2`$ and W$`3.2`$ at the Hamburg Synchrotron Laboratory (HASYLAB), at the undulator beamline U$`2`$-FSGM and the $`2`$m SEYA beamline at the Berliner Elektronenspeicherring Gesellschaft für Synchrotronstrahlung mbH (BESSY). At the storage rings and monochromators used in these studies, highly linearly polarized synchrotron radiation was available. In addition, the crossed undulator U$`2`$-FSGM beamline gave the possibility to use vertically oriented linearly polarized light. At this latter facility, polarization-dependent measurements could be performed without changing any other parameter in the experiment. A total energy resolution (beamline and analyzer) of better than $`70`$ meV and an angular acceptance of the analyzer of $`\pm 1`$ degree was used. Fresh samples of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> were cleaved in situ at a base pressure of $`1\times 10^{10}`$mbar, and spectra were taken within six hours after cleavage. Samples were either pre-oriented using X-ray diffraction measurements or aligned in-situ with the aid of low energy electron diffraction (LEED). In all cases, the fine angular adjustment was carried out using the $`𝐤`$-space symmetry of the sharp ARPES features related to non-bonding O2$`p`$ states around ($`\pi `$,$`\pi `$) and ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$), as discussed in Refs. and .
The high-quality single crystalline Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> was grown according to the method described in Ref. , where pre-dried high-purity SrCO<sub>3</sub>, SrCl<sub>2</sub> and CuO in a ratio 1:1:1 were melted at 1100$`^{}`$C. Although T<sub>N</sub> for Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> is $`251`$ K (Ref. ), all ARPES experiments were carried out at room temperature. This decision was based on two arguments. Firstly, and most importantly, as these experiments involve electron ejection from perfect single crystals of a compound with an energy gap of the order of 2 eV, we had to eliminate uncertainties due to charging effects, which meant measuring at room temperature. Secondly, although we measure at a temperature some 50 K above the Néel temperature, it has been shown that the antiferromagnetic spin correlation length of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> at $`350`$ K is still $`250`$ Å (Ref. ), meaning that the lowest binding-energy hole state created in the photoemission process is still embedded in an antiferromagnetic spin background.
To be absolutely sure that variations in the flux of the incident photons from the storage ring did not lead to charging-induced energy shifts, we also adjusted the beamline such that the photon flux impinging on the sample (monitored with a gold mesh upstream of the sample) was constant. This means that measuring time is the only parameter required for the normalization of the data. Bearing in mind that the data presented here are a selection of data from about 80 different cleavages, we made the qualitative observation that best and sharpest energy distribution curves (EDC) were observed for samples which behaved most sensitively to charging effects. This relation, of course, is reasonable as it correlates cleavage quality with the intrinsically highly insulating nature of perfect Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>.
Essentially two different experimental geometries were used. They can be characterized by the orientation of two planes: a) the plane of polarization and b) the emission plane. The plane of polarization is the plane spanned by the vector of the direction of the synchrotron light and the vector of its electric field (vector of polarization). The emission plane includes the vector of the direction of the photoelectron and the surface normal. Throughout this paper we call the experimental geometry parallel if these two planes are parallel, perpendicular if these two planes are perpendicular to each other (see Fig. 2 for a sketch of the two geometries).
## III Photon-energy and polarization dependence of the first electron-removal states of an undoped CuO<sub>2</sub>-plane
As mentioned above, both the inconsistencies in the literature regarding the dispersion relation of the first electron-removal states along the $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) direction (Ref. ) in Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>, as well as the ongoing controversy regarding the photon-energy dependence of the ARPES data from Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub>, mean that before comparison with theory is carried out, the photon-energy dependence of the first electron-removal states should be examined in detail.
In Fig. 3, photon-energy dependent ARPES data for the first electron-removal state in Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> at two points in the Brillouin zone: a) ($`\pi /2`$,$`\pi /2`$) (perpendicular geometry); and b) ($`0.7\pi `$,$`\mathrm{\hspace{0.17em}0}`$) (parallel geometry) are shown. In these regions of $`𝐤`$-space the spectral weight of the first electron-removal state is known to be at a local maximum along the respective $`𝐤`$-space directions . The spectra, which are normalized as described in the last section, illustrate clearly the strong variation of the first electron-removal state intensity with photon-energy. The data shown in Fig. 3 represent only a small portion of the photon-energy dependent data recorded, and are intended to give the reader a direct impression of the strength of the effects at play. Each of the shown spectra is part of a short k-series of three or five spectra. This was done to ensure we always captured the spectrum with the highest spectral intensity for the first electron-removal states. Within the errors given by the finite $`𝐤`$-resolution of the experimental setup ($`0.054`$ Å<sup>-1</sup> for 16 eV photon-energy up to $`0.152`$ Å<sup>-1</sup> for 80 eV photon-energy) we found the highest intensities always at the same $`𝐤`$-positions, namely at ($`\pi /2`$,$`\pi /2`$) and at ($`0.7\pi `$,$`\mathrm{\hspace{0.17em}0}`$).
This is in contrast to the results reported in Ref. , where series of EDCs on Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) taken with two different photon-energies show differences not only in intensity, but also in the EDC-derived dispersion relation of the first ionization states. This was discussed in Ref. as not being due to either a) the experimental setup; b) the excitation of different initial states or c) in the variation of k of the photoelectron. The effect was rather attributed to the strong impact of the matrix element, not only on the strength of the photoemission signal but also on the position of the quasiparticle peak in the photoemission spectra, leading to differences in the EDC-derived dispersion relations. We are forced to disagree with the last point, as in our extensive collection of ARPES data, there was never evidence for a photon-energy dependent shift of the $`𝐤`$-position for which the first electron-removal states have minimum binding energy along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$). The same, in fact, holds for other directions in the Brillouin zone as well as for the dispersion of other features with low binding energy. If matrix element effects change the EDC-derived dispersion relation, then this change is, at least for our data, smaller than the energy resolution and angular resolving power of our ARPES experiments, which is equal to $`\mathrm{\Delta }_k`$ or better than $`\mathrm{\Delta }_E`$ the values used in Ref. ).
Figure 4 shows an analysis of the spectral weight of the first electron-removal states at the same $`𝐤`$-points as shown in Fig. 3. These data are derived from a large set of photon-energy dependent ARPES data covering measurements from 10 (6) cleaves for the upper (lower) panels. From Fig. 4 it is clear that the ARPES spectral weight of the first electron-removal states oscillates with the final state kinetic energy, that is with k. We used a value of $`8.0`$ eV for the inner potential $`E_0`$ ($`E_0=V_0\mathrm{\Phi }`$, where $`\mathrm{\Phi }`$ is the work function) to calculate k, a value which in the range of $`6.98.9`$ eV used by other groups . We see clear maxima at k=$`0.82`$, $`1.63`$, $`2.40`$ and $`3.12`$ Å<sup>-1</sup>, corresponding to photon-energies of $`16`$, $`25`$, $`35`$ and $`48`$ eV.
The oscillatory nature of the photon-energy dependence, coupled to the absence of a classical resonance behavior at the Cu 3p threshold (around 74-76 eV photon-energy) indicates that the factor dominating the observed behavior is something other than the atomic photoionization cross-sections, and could be related, for example, to the extreme two-dimensionality of the electronic states concerned. This could lead to a matching of the final state k to the periodicity of the unit cell of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> in c-direction. Note that the differences between the four k is $`0.720.81`$ Å<sup>-1</sup> which represents in real space a distance of $`7.88.7`$ Å which is comparable to c-axis separation of two neighboring CuO<sub>2</sub>-planes ($`=7.805`$ Å) in Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>. The oscillation of the photoemission intensity with photon-energy can therefore be attributed to interference effects of the photoelectron wave diffracted from the c-axis periodicity of the layered crystal structure, similar to the explanation of the strong photon-energy dependence of the photoemission intensity from the molecular orbitals of C<sub>60</sub> (Ref. ).
These strong variations in intensity as a function of photon-energy are also reminiscent of ARPES data of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>, in which an out-of-phase behavior regarding intensity vs. k was observed for the CuO<sub>2</sub>-derived band which crosses the Fermi energy near the $`\overline{\text{X}}`$ point and the so-called “1 eV peak” . This fact was, at that time, used to argue against the surface-state origin of the 1 eV peak, a feature which is now believed to be due to non-mixing O2$`p`$ states of particular symmetry . Nevertheless, the data presented here, taken together with the theoretical and experimental investigations on Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> (Refs. and ), indicate clearly that care should be taken in the interpretation of absolute spectral weights observed in the ARPES data of the layered cuprates, as matrix element and diffraction effects do play an important role in these systems.
A further experimental variable in an ARPES experiment is the polarization of the incoming radiation. Figure 5 shows two series of ARPES measurements on Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> recorded along the high-symmetry directions $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) and $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$). The photon-energy was set to $`22.4`$ eV, which is near a maximum of intensity as shown in Fig. 3 above. In each case, the series are presented in pairs of data sets recorded with perpendicular and parallel polarization geometries as described in the experimental section.
Along the $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) direction in $`𝐤`$-space the first electron-removal state feature shows highest photoemission intensity with perpendicular geometry (Fig. 5a), whereas along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) maximal intensities are observed in the parallel geometry (Fig. 5d), a result which tallies with earlier measurements with 25 eV photon-energy . This remarkable dependence of the photoemission intensity on the polarization was first explained in the context of the strong polarization dependence in photoemission data from surface states .
The physical picture behind the polarization dependence can be described as follows. The interaction operator $`𝐀𝐩`$ has even (+) parity in a parallel and odd (-) parity in a perpendicular experimental geometry. Assuming the applicability of the Zhang-Rice singlet state construction to the first electron-removal final state, in a many-body picture this state belongs to the $`A_{1g}`$ representation, therefore being totally symmetric. Thus in this representation the Zhang-Rice singlet has even parity with respect to a mirror plane along the Cu3$`d_{x^2y^2}`$-O2$`p_x`$/O2$`p_y`$ orbital bonds ($`M_2`$) as well as at $`45`$ degrees to the bonds ($`M_1`$). In $`𝐤`$-space $`M_1`$ corresponds to the $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) direction and $`M_2`$ to $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$). The initial state (ground state) is a one-hole state with $`d_{x^2y^2}`$-symmetry, and therefore has even parity with respect to $`M_2`$ and odd parity with respect to $`M_1`$. The matrix element thus formally vanishes for the two cases
parallel for $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) $``$ $`M+|+|=0`$
perpendicular for $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) $``$ $`M+||+=0`$
This argumentation also holds in the one band picture. Here the initial state (ground state) is totally symmetric and the final state has $`d_{x^2y^2}`$-symmetry (as shown in Fig. 1), leading to the same result.
Thus, the observed polarization dependence of the first electron-removal states of an undoped CuO<sub>2</sub>-plane indicates that these states have a symmetry fully consistent with that of the Zhang-Rice singlet. These results amend earlier reports (Ref. ) regarding the polarization dependence of the first electron-removal states along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) in Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>.
## IV The dispersion relation and spectral function of the first electron-removal states
The common picture given by the $`t`$-$`J`$-model and its extensions is a strong dispersion of the Zhang-Rice singlet state along the $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) direction with the minimum binding energy at ($`\pi /2`$,$`\pi /2`$). At this $`𝐤`$-point, the spectral weight of the first electron-removal state also has its maximum and vanishes going away from ($`\pi /2`$,$`\pi /2`$). Recent ARPES experiments on Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> confirm this . Along the $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) direction the $`t`$-$`J`$-model predicts a rather low binding energy near the high-symmetry point ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) which is, in fact, almost energetically degenerate with that at ($`\pi /2`$,$`\pi /2`$). In the same model, the spectral weight of the quasiparticle increases and is maximal at ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$). This is in contrast to the result given in Ref. where no dispersion was observed and in Refs. where the binding energy of the first electron-removal state increases and the quasiparticle weight decreases after ($`\pi /2`$,$`\mathrm{\hspace{0.17em}0}`$). For the ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$)$``$($`0`$,$`\pi `$) direction, the $`t`$-$`J`$-model predicts only little dispersion, whereas experiment has shown a strong isotropic dispersion around the ($`\pi /2`$,$`\pi /2`$) point (Ref. ). Extensions of the $`t`$-$`J`$-model and the spin and charge separation ansatz exhibit new properties along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) and ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$)$``$($`0`$,$`\pi `$). The Zhang-Rice singlet state around ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) can now lower its energy by delocalization and is thus pushed to higher binding energies and the dispersion in this direction is essentially parabolic. In these models, both a quasiparticle dispersion along ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$)$``$($`0`$,$`\pi `$) which is isotropic around ($`\pi /2`$,$`\pi /2`$) as well as a reduced spectral weight near ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) can be achieved. In the spin and charge separation model the isotropy along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) and ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$)$``$($`0`$,$`\pi `$) is intrinsic and near ($`\pi /2`$,$`\pi /2`$) the dispersion relation is wedge-like rather than parabolic. Surprisingly, the data in Ref. agree well with the extended $`t`$-$`J`$-model along ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$)$``$($`0`$,$`\pi `$), but not along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$). In contradiction to that, other data along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) are consistent with the extended $`t`$-$`J`$-model and the spin and charge ansatz , a view which appears to be supported by further data sets . Thus, the experimental dispersion along the $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) remains controversial, and yet it is the dispersion and the evolution of the spectral weight along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) that is of the deepest theoretical interest for the following reasons. Firstly, the states near ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) evolve with increasing hole concentration to become the flat bands located near the chemical potential in the metallic systems. These flat band regions are believed by some to hold the key to high temperature superconductivity. Secondly, the dispersion and quasiparticle spectral weight along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) are fingerprints for the different models and parameters used in them. Thirdly, there are strong indications that the dispersion in the insulator is closely connected with the pseudogap in the underdoped region of the high $`T_c`$ superconductor phase diagram and thus may also be related to the superconducting gap in the doped superconducting systems .
Consequently, the importance of experimental data which allow the determination of the dispersion and spectral weight of the first electron-removal states along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) cannot be overestimated. Nevertheless, the experimentalist is faced with a number of challenges when collecting data along this direction: firstly, at several points in the Brillouin zone, there is no or only weak photointensity from the first electron-removal state. For example, near $`\mathrm{\Gamma }`$ the Zhang-Rice singlet has no spectral weight, because the Cu3$`d_{x^2y^2}`$ and O2$`p_x`$/O2$`p_y`$ orbitals cannot hybridize there , in addition the matrix elements for emission along the surface normal are formally zero for a perfectly two-dimensional electronic state located in the CuO<sub>2</sub>-plane. Furthermore, near ($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) the photoemission intensity of the first electron-removal states becomes weak as predicted by the extended $`t`$-$`J`$-models and the spin and charge separation ansatz. Secondly, only perfectly aligned samples with fresh and clean surfaces excited with UV-light in a well-adjusted measurement geometry will show usable photoemission data from the first electron-removal state. Lastly, as the clean, defect-free surface of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> is highly insulating, avoiding charging effects can be difficult.
In Fig. 6 we present series of ARPES spectra recorded along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) taken with parallel polarization and $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) with perpendicular polarization and with photon-energies of $`16`$, $`22`$ and $`35`$ eV - thus measuring near a peak in the photoemission intensity oscillation in each case. The position of the low energy part of the dispersive structure has been fitted using a Gaussian function while the high binding-energy background was modelled simply with the tail of another Gaussian. Although this fit is physically somewhat simplistic, it proved to be sufficient to reliably find the positions of the peak as indicated by the triangles in each plot. The absolute value on a binding-energy scale changes from one cleavage to the next and is rather arbitrary due to the insulating nature of the substance. However, the two series with $`16`$ eV were taken during one and the same cleavage, which is the prerequisite to compare the binding-energy scales of both series. The binding-energy scales of the $`22`$ and $`35`$ eV data have been adjusted to that of the $`16`$ eV data. We present a summary of the dispersion relations derived from the fits to the data in Fig. 7.
We repeat here that we observe no evidence for a photon-energy dependence of the dispersion relation of the first ionization states, as is consistent with the arguments regarding the spectral function given in the introduction. The dispersion of the first electron-removal states along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) agrees with the previous results having a parabolic shape with its minimum binding energy at ($`\pi /2`$,$`\pi /2`$). Along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) we also find a parabolic dispersion, in agreement with the dispersion reported in Refs. , but in contrast to the data given in Ref. . For $`16`$ eV photon-energy, we measured 72 meV difference in the lowest binding energies of the first electron-removal state along each of the two high-symmetry directions, which is half as small as in Ref. . In Ref. , a series of spectra representing a “maximum intensity cut” for a photon-energy $`22.4`$ eV along a line from ($`\pi /2`$,$`\pi /2`$) to ($`0.67\pi `$,$`\mathrm{\hspace{0.17em}0}`$) was shown, which displayed a dispersion of about $`300`$ meV. However, as we will show, the points of maximum intensity along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) do not coincide with those possessing minimum binding energy. Thus such a maximum cut will not trace the line of minimum binding energy. It turns out that this difference strongly influences the parameters in model Hamiltonians fitted to the dispersion curves.
We now go on to compare the experimental dispersion with theoretical results obtained within the extended $`t`$-$`J`$-model
$$H=\underset{i\sigma }{}\underset{l=1}{\overset{3}{}}t_lX_i^{\sigma 0}X_{i+l}^{0\sigma }+J\underset{i,m}{}𝐒_i𝐒_m,$$
(3)
with the additional hopping terms to second ($`t_2`$) and third ($`t_3`$) neighbors now added besides the dominating nearest neighbor hopping $`t=t_1`$. The Hamiltonian is written in terms of Hubbard operators $`X_i^{\sigma 0}=c_{i\sigma }^{}(1n_{i\sigma })`$ where $`\sigma =\pm 1`$ is the spin index. The values of the hopping terms are determined by mapping the more realistic Emery model with parameters derived from a constrained density-functional calculation to its low binding-energy part by means of the cell-perturbation method. This procedure gives $`t_2/t_1=0.08`$ and $`t_3/t_1=0.15`$. The theoretical dispersion relation shown in Fig. 7 is calculated using a variational method involving a spin polaron of small radius . There are also related works which lead to similar dispersions.
It has been shown in exact diagonalisation studies of small clusters that the dispersive bandwidth scales with the single parameter, $`J`$, whereas the form of the dispersion curve is fixed by the ratios $`t_2/t_1`$ and $`t_3/t_1`$. Therefore we use the variational ansatz for $`J=t_1`$ and then we scale the bandwidth in Fig. 7 with a factor $`J=0.22`$ eV. For reference, we also show the dispersion relation without additional hopping terms ($`t_2=t_3=0`$, scaling factor $`J=0.28`$ eV), which shows a much too small energy difference between the lowest binding energies of the first electron-removal states along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) and $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) to obtain good agreement with the experimental spectra.
Another important prediction from model Hamiltonians is the evolution of the spectral weight along certain directions in the Brillouin zone. Here we map this evolution and distinguish between the coherent and incoherent parts of the first electron-removal states following the procedure proposed in Ref. , in which the contribution from the main valence band tail was first subtracted from the spectra. The photoemission intensity was then defined as the coherent spectral weight and fitted using a Gaussian. The ARPES intensity which is neither part of the valence band tail nor in the Gaussian is then taken to be a measure of the incoherent part of the first electron-removal state spectral function \[see Ref. \]. In our case, we have used a Lorentzian broadened by the experimental resolution (i.e. a Voigt profile) to fit the lowest binding energy intensity (the coherent part). Of course, this procedure is not physically rigorous, but does offer a rough estimation of the possible split between the coherent and incoherent spectral weight.
In Fig. 8, the momentum distributions of the two parts of the ARPES spectra are shown for $`16`$ eV and $`22`$ eV excitation energy. Note that the overall shape of these distributions is independent of the photon-energy. Along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) both the coherent and incoherent parts of the first electron-removal states are symmetric around ($`\pi /2`$,$`\pi /2`$) at which point they both reach their maximum weight. In the $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) direction, we observe a steady increase in spectral weight of both components up to ($`0.7\pi `$,$`\mathrm{\hspace{0.17em}0}`$) after which point it drops fast. One can find signs of a qualitatively similar behavior in ARPES data from other groups and for measurements on the related substance Ca<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> (Ref. ), although the coherent and incoherent spectral weights were not analysed in these studies. Note that along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) the minimum binding energy of the low binding-energy feature and the maximum of the coherent intensity do not coincide at the same $`𝐤`$-point. In calculations using extensions of the $`t`$-$`J`$-model this behavior was predicted and can now be considered as experimentally verified.
Within the simple fit procedure we can give these quantities some numbers: the ratio of the coherent to incoherent spectral weight for the two locations in $`𝐤`$-space is 0.6 (($`\pi /2`$,$`\pi /2`$)) and 0.5 (($`0.7\pi `$,$`\mathrm{\hspace{0.17em}0}`$)) for 16 eV photon-energy and 2.4 (($`\pi /2`$,$`\pi /2`$)) and 3.8 (($`0.7\pi `$,$`\mathrm{\hspace{0.17em}0}`$)) for 22 eV photon-energy. It is important to realise that we observe here an extremely strong apparent dependence of the ratio between the coherent and incoherent spectral weight upon the photon-energy. To discuss the physical significance of this, we re-write equation 2 as
$$S(𝐤,E)=\frac{1}{\pi }\text{Im}ic_𝐤(\frac{1}{EE(𝐤)\mathrm{\Sigma }}+G_{\text{inc}})c_𝐤^{}i,$$
(4)
in order to separate the coherent from the incoherent spectral weight ($`\mathrm{\Sigma }`$ is the self-energy of the quasiparticle). Plugging this into equation 1, it is immediately clear that in taking the ratio of the coherent and the incoherent intensity, the matrix element cancels out. Thus, this ratio is then determined by the spectral function $`S(𝐤,E)`$ alone and should not depend on the excitation energy.
Our observation that the ratio does appear to depend on h$`\nu `$ leads then to the following possible conclusions: (i) the fit procedure used is too simple and therefore does not correctly quantify, however roughly, the coherent and incoherent parts of the spectral function. (ii) there are significant extrinsic contributions to the photoemission signal in this binding energy region. These could be the result of an energetic shift of spectral weight due to inelastic losses , which could be sensitively dependent on the photoelectron kinetic energy, as are the data from electron energy loss spectroscopy in reflection of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> (Ref. ). (iii) the photoemission intensity between the the main valence band edge and the chemical potential is not derived from the spectral weight (coherent + incoherent) resulting from a single electronic state - i.e. there is spectral weight from an additional, different electronic state in this energy region.
The fit procedure, as we have discussed above, is simplistic - however the coherent-incoherent intensity ratios vary by a factor of more than four between the two photon-energies. The third possible conclusion - that there has to be spectral weight from more than one state in this energy region could have important implications. We can rule out a significant contribution from secondary electrons for the relatively high kinetic energies dealt with here. Intensity from surface states is unlikely, as XPD data shows that the cleavage surface of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> is terminated within the SrCl units, which are essentially ionic and therefore do not support electronic states close to the chemical potential. In addition, our LEED investigations also gave no evidence for a reconstruction of this termination layer. Spin-resolved resonant photoemission of CuO (Ref. ) has shown intensity due to triplet states within 1eV of the Zhang-Rice singlet, although those states would be expected to have the same photon-energy dependence as the singlet in our experiment.
If there was an additional, different state in this energy region, this would lead to a complicated $`𝐤`$\- and h$`\nu `$-dependent overlap between the intensity of this state and the higher binding energy components attributed up till now to incoherent weight from the Zhang-Rice singlet. While this would not be expected to have a large impact on the observed dispersion relation for the Zhang-Rice singlet (which is, after all, a quantity derived from the spectral structure at lowest binding energies for each $`𝐤`$-point), it could, however lead to the observed photon-energy dependence of the ratio between the low binding-energy (“coherent”) and higher binding-energy (“incoherent”) parts of the photoemission spectra.
## V Summary
In conclusion, we have presented a detailed ARPES study of the low binding-energy occupied electronic structure of Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>, which corresponds to an investigation of the first electron-removal states of an undoped CuO<sub>2</sub>-plane. Our experiments, and the comparison of their results with theoretical models, have revealed the following main points:
1. The photoemission signal of the first electron-removal states at both ($`\pi /2`$,$`\pi /2`$) and ($`0.7\pi `$,$`\mathrm{\hspace{0.17em}0}`$) exhibits a marked photon-energy dependence. The intensity profile shows strong oscillations with maxima near $`16`$, $`25`$, $`35`$ and $`49`$ eV, corresponding to final state crystal momenta k=$`0.82`$, $`1.63`$, $`2.40`$ and $`3.12`$ Å<sup>-1</sup>. This strong photon-energy dependence has complicated comparisons between data in the literature from different groups as regards both the spectral weight and spectral form of the first electron-removal states in these systems. We attribute the oscillation of photoemission intensity (which has a period in k of ca. $`0.75`$ Å<sup>-1</sup>) to the diffraction of the photoelectron wave on the periodic c-axis separation of the CuO<sub>2</sub>-planes of $`8.2`$ Å.
2. Along the high-symmetry directions $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) and $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) the first electron-removal states shows a strong polarization dependence. This can be linked to the strongly polarization-dependent matrix element, which in turn allows the determination of the symmetry of the first electron-removal state itself. For both high-symmetry directions we observe a polarization dependence in keeping with that expected for a Zhang-Rice singlet state in the framework of either a three-band or one-band model Hamiltonian.
3. Our data show that the dispersion of the first electron-removal states along both high symmetry directions ($`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) and $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$)) is parabolic-like and independent of the excitation energy. This, and the rather large difference in lowest binding energy of the first electron-removal state along these directions, shows the validity of the extended $`t`$-$`J`$-model for describing the disperion relation of a single hole in an antiferromagnetic CuO<sub>2</sub> plane. Thus, the inclusion of second ($`t_2`$) and third ($`t_3`$) neighbor hopping terms with realistic values of $`t_2=0.08`$ and $`t_3=0.15`$ in units of the next neighbor hopping $`t=t_1`$ are required.
4. Upon application of a simple fit procedure, we infer the momentum distribution of the spectral weight of the coherent and incoherent part of the first electron-removal state to have its maximum along $`\mathrm{\Gamma }`$($`\pi `$,$`\pi `$) at ($`\pi /2`$,$`\pi /2`$) being symmetrically suppressed away from this point. Along $`\mathrm{\Gamma }`$($`\pi `$,$`\mathrm{\hspace{0.17em}0}`$) the spectral weights of both parts reach their maximum at ($`0.7\pi `$,$`\mathrm{\hspace{0.17em}0}`$) and then drop fast. The ratio between the coherent and incoherent spectral weight is strongly photon-energy dependent, which, at first sight would appear to violate the physics of the spectral function. There are different possible explanations for this including: (i) the necessity for a more sophisticated framework in which to analyse the weight of the coherent and incoherent contributions to the spectral weight; (ii) significant (h$`\nu `$-dependent) intensity due to extrinsic processes (iii) intensity in this energy region due to intrinsic electronic states other than the Zhang-Rice singlet.
We gratefully acknowledge the stimulating conversations with S. Haffner. This work was supported by the BMB+F under contract number 05 SB8 BDA6 and by the DFG under Fi439/7-1 and as part of the Graduiertenkolleg “Struktur und Korrelationseffekte im Festkörper” der TU Dresden.
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# The ATESP Radio Survey
## 1 Introduction
Recent deep radio surveys ($`S<<1`$ mJy) have shown that normalized radio source counts flatten below a few mJy. This has been interpreted as being due to a new population which does not show up at higher fluxes where counts are dominated by active galactic nuclei (AGN). To clarify its nature it is necessary to get detailed information on the radio properties of normal galaxies in the nearby universe (z $``$ 0.15), down to faint flux limits and to have at hand large samples of mJy and sub-mJy sources, for subsequent optical identification and spectroscopic work. To this end we have surveyed a large area ($``$ 26 sq. degr) with the ATCA at 1.4 GHz with a bandwidth of $`2\times 128`$ MHz. The properties of normal nearby galaxies can be easily derived, because this area coincides with the region in which the ESP redshift survey was conducted (Vettolani et al. Vettolani97 (1997)). Samples of faint galaxies over large areas are necessary to avoid bias due to local variations in their properties. Present samples of faint radio sources are confined to small regions with insufficient source statistics.
The present paper contains the radio source catalogue derived from the ATESP survey. The full outline of the radio survey, its motivation in comparison with other surveys, and the description of the mosaic observing technique which was used to obtain an optimal combination of uniform and high sensitivity over the whole area has been presented in Prandoni et al. (2000, Paper I). The source counts derived from the ATESP survey will be presented and discussed separately in a following paper. The paper is organized as follows: in Sect. 2 we describe the source detection and parameterization; the catalogue is described in Sect. 3 and the accuracy of the parameters (flux densities, sizes and positions) are discussed in Sect. 4.
## 2 Source Detection and Parameters
The ATESP survey consists of 16 radio mosaics with spatial resolution $`8\mathrm{}\times 14\mathrm{}`$. The survey was designed so as to provide uniform sensitivity over the whole region ($`26`$ sq. degrees) of the ESP redshift survey. To achieve this goal a larger area was observed, but we have excluded from the analysis the external regions (where the noise is not uniform and increases radially). In the region with uniform sensitivity the noise level varies from 69 $`\mu `$Jy to 88 $`\mu `$Jy, depending on the radio mosaic, with an average of 79 $`\mu `$Jy (see $`\sigma _{\mathrm{fit}}`$ values reported in Table 3 of paper I and repeated also in Table 2 of Appendix B, at the end of this paper). For consistency with paper I, in the following such sensitivity average values are denoted by the symbol $`\sigma _{\mathrm{fit}}`$. This means that sensitivity values have been defined as the *full width at half maximum* (FWHM) of the Gaussian that fits the pixel flux density distribution in each mosaic (see paper I for more details).
A number of source detection and parameterization algorithms are available, which were developed for deriving catalogues of components from radio surveys. We decided to use the algorithm *Image Search and Destroy* (IMSAD) available as part of the *Multichannel Image Reconstruction, Image Analysis and Display* (MIRIAD) package (Sault & Killeen Sault95 (1995)), as it is particularly suited to images obtained with the ATCA.
IMSAD selects all the connected regions of pixels (*islands*) above a given flux density threshold. The *islands* are the sources (or the source components) present in the image above a certain flux limit. Then IMSAD performs a two-dimensional Gaussian fit of the *island* flux distribution and provides the following parameters: position of the centroid (right ascension, $`\alpha `$, and declination, $`\delta `$), peak flux density ($`S_{\mathrm{peak}}`$), integrated flux density ($`S_{\mathrm{total}}`$), fitted angular size (major, $`\theta _{\mathrm{maj}}`$, and minor, $`\theta _{\mathrm{min}}`$, FWHM axes, not deconvolved for the beam) and position angle (P.A.).
IMSAD proved to have an average success rate of $`90\%`$ down to very faint flux levels (see below). Since IMSAD attempts to fit a single Gaussian to each *island*, it obviously tends to fail (or to provide very poor source parameters) when fitting complex (*i.e.* non-Gaussian) shapes.
### 2.1 Source extraction
We used IMSAD to extract and parameterize all the sources and/or components in the uniform sensitivity region of each mosaic<sup>1</sup><sup>1</sup>1To avoid interpolation and completeness problems the effective search area in each mosaic was slightly larger both in declination and in right ascension. The masked region in mosaic fld69to75 has been excluded from the search (see paper I for more details).. As a first step, a preliminary list containing all detections with $`S_{\mathrm{peak}}4.5\sigma `$ (where $`\sigma `$ is the average mosaic rms flux density) was extracted. Detection thresholds vary from 0.3 mJy to 0.4 mJy, depending on the radio mosaic.
### 2.2 Inspection
We visually inspected all *islands* ($`5000`$) detected, in order to check for obvious failures and/or possibly poor parameterization, that need further analysis. Problematic cases were classified as follows:
* *islands* for which the automated algorithm provides poor parameters and therefore are to be re-fitted with a single Gaussian, constraining in a different way the initial conditions ($`5\%`$ of the total).
* *islands* that could be better described by two or more split Gaussian components ($`3\%`$ of the total);
* *islands* which cannot be described by a single or multiple Gaussian fit ($`1\%`$ of the total);
* *islands* corresponding to obviously unreal sources, typically noise peaks and/or image artifacts in noisy regions of the images ($`2\%`$ of the total);
The goodness of Gaussian fit parameters was checked by comparing them with reference values, defined as follows. Positions and peak flux densities were compared to the corresponding values derived by a second-degree interpolation of the *island*. Such interpolation usually provides very accurate positions and peak fluxes. Gaussian integrated fluxes were compared to the ones derived directly by summing pixel per pixel the flux density in the source area, defined as the region enclosed by the $`3\sigma `$ flux density contour. Flux densities were considered good whenever the difference between the Gaussian and the reference value $`SS_{\mathrm{ref}}`$ was $`0.2S_{\mathrm{ref}}`$. Positions were considered good whenever they did fall within the $`0.9S_{\mathrm{peak}}`$ flux contour.
### 2.3 Re-fitting
For the first three groups listed above *ad hoc* procedures were attempted aimed at improving the fit.
Single-component fits were considered satisfactory whenever positions and flux densities satisfy the tolerance criteria defined above.
In a few cases Gaussian fits were able to provide good values for positions and peak flux densities, but did fail in determining the integrated flux densities. This happens typically at faint fluxes ($`<10\sigma `$). Gaussian sources with a poor $`S_{\mathrm{total}}`$ value are flagged in the catalogue (see Sect. 3.4).
The *islands* successfully split in two or three components are 67 in total (64 with two components and 3 with three components).
For non-Gaussian sources we adopted as parameters the reference positions and flux densities defined above. The source position angle was determined by the direction in which the source is most extended and the source axes were defined as *largest angular sizes* (*las*), i.e. the maximum distance between two opposite points belonging to the $`3\sigma `$ flux density contour along (major axis) and perpendicular to (minor axis) the same direction. All the non-Gaussian sources are flagged in the catalogue (see Sect. 3.4).
## 3 The Source Catalogue
The procedures described in the previous section yielded a preliminary list of sources (or source components) for further investigation. In order to minimize the incompleteness effects present at fluxes approaching the source extraction threshold (see Fig. 1) we decided to insert in the final catalogue only the sources with $`S_{\mathrm{peak}}6\sigma `$, where $`\sigma `$ is the mosaic rms flux density. This threshold has been chosen after inspection of the local noise distribution. The local noise ($`\sigma _{\mathrm{local}}`$) has been defined as the average noise value in a box of about $`8\mathrm{}\times 8\mathrm{}`$ size around a source. Usually the local noise does not show significant systematic departures from the mosaic average rms value: the $`\sigma _{\mathrm{local}}/\sigma _{\mathrm{fit}}`$ distribution can be described fairly well by a Gaussian with $`\mathrm{FWHM}=0.14`$ and peak value equal to 1.01 (see Fig. 2).
This can be seen also in Fig. 3, where we show, for each source, the signal-to-noise ratio defined using both $`\sigma _{\mathrm{local}}`$ and $`\sigma _{\mathrm{fit}}`$. The two signal-to-noise ratios mostly agree with each other, although a number of significant departures are evident for the faintest sources. This is due to the presence of some residual areas where the noise is not random due to systematic effects (noise peaks and stripes). This is caused by the limited dynamical range in presence of very strong sources (stronger than $`50100`$ mJy, see also paper I). It is worth noting that also the systematic departures from the expected behavior at the brightest end of the plot ($`S_{\mathrm{peak}}/\sigma _{\mathrm{local}}<S_{\mathrm{peak}}/\sigma _{\mathrm{fit}}`$) are a consequence of the same problem.
In mosaic regions where local noise is significantly larger, we applied a $`6\sigma _{\mathrm{local}}`$ cut-off if $`\sigma _{\mathrm{local}}/\sigma _{\mathrm{fit}}1.2`$ (assuming a normal distribution for the noise, the probability to get a local to average noise value $`1.2`$ is $`0.1\%`$). This resulted in the rejection of 32 sources (see the region in Fig. 3 defined by $`S_{\mathrm{peak}}6\sigma _{\mathrm{fit}}`$, $`S_{\mathrm{peak}}<6\sigma _{\mathrm{local}}`$ and $`\sigma _{\mathrm{local}}/\sigma _{\mathrm{fit}}1.2`$).
The criteria discussed above proved to be very effective in selecting out noise artifacts from the catalogue. Nevertheless a few (6) sources, which satisfy both the average and the local noise constraints are, however, evident noise artifacts at visual inspection. Such objects have been rejected from the final catalogue.
The adopted criteria for the final catalogue definition also allowed us to significantly reduce the number of poor Gaussian fits (see Sect. 2.3) since a large fraction of them ($`65\%`$) have fainter $`S_{\mathrm{peak}}`$: we are left with 50 poor Gaussian fits (flagged ‘S\*’, see Sect. 3.4) in the final catalogue.
A total number of 3172 source components entered the final catalogue. Some of them have to be considered different components of a unique source and, as discussed later in Sect. 3.2, the number of distinct sources in the ATESP catalogue is 2960.
### 3.1 Deconvolution
The ratio of the integrated flux to the peak flux is a direct measure of the extension of a radio source:
$$S_{\mathrm{total}}/S_{\mathrm{peak}}=\theta _{\mathrm{min}}\theta _{\mathrm{maj}}/b_{\mathrm{min}}b_{\mathrm{maj}}$$
(1)
where $`\theta _{\mathrm{min}}`$ and $`\theta _{\mathrm{maj}}`$ are the source FWHM axes and $`b_{\mathrm{min}}`$ and $`b_{\mathrm{maj}}`$ are the synthesized beam FWHM axes. The flux ratio can therefore be used to discriminate between extended (larger than the beam) and point-like sources.
In Fig. 4 we have plotted the flux ratio as a function of the signal-to-noise for all the sources (or source components) in the ATESP catalogue. The flux density ratio has a skew distribution, with a tail towards high flux ratios due to extended sources. Values for $`S_{\mathrm{total}}/S_{\mathrm{peak}}<1`$ are due to the influence of the image noise on the measure of source sizes (see Sect. 4). To establish a criterion for extension, such errors have to be taken into account. We have determined the lower envelope of the flux ratio distribution (the curve containing 90% of the $`S_{\mathrm{total}}<S_{\mathrm{peak}}`$ sources) and we have mirrored it on the $`S_{\mathrm{total}}>S_{\mathrm{peak}}`$ side (upper envelope in Fig. 4). We have considered as unresolved all sources laying below the upper envelope. The upper envelope can be characterized by the equation:
$$S_{\mathrm{total}}/S_{\mathrm{peak}}=1.05+\left[\frac{10}{(S_{\mathrm{peak}}/\sigma _{\mathrm{fit}})^{1.5}}\right].$$
(2)
From this analysis we found that 1864 of the 3172 sources (or source components) in the catalogue (i.e. $`60\%`$) have to be considered unresolved.
It is worth noting that the envelope does not converge to 1 going to large signal-to-noise values. This is due to the radial smearing effect. It systematically reduces the source peak fluxes, yielding larger $`S_{\mathrm{total}}/S_{\mathrm{peak}}`$ ratios (see discussion in Sect. 4.2). From Fig. 4 we can quantify the smearing effect in $`5\%`$ on average.
Deconvolved angular sizes are given in the catalogue only for sources above the upper curve (filled circles in Fig. 4). For unresolved sources (dots in Fig. 4) deconvolved angular sizes are set to zero.
Note that no bandwidth correction to deconvolved sizes has been applied. Correcting for such effect would be somewhat complicated by the fact that each source in the radio mosaics is a sum of contributions from several single pointings.
### 3.2 Multiple Sources
In Fig. 5 the (nearest neighbor) pair density distribution is shown as a function of distance (histogram). Also indicated is the expected distribution if all the sample sources (components) were randomly distributed in the sky. The expected distribution has been scaled so as to have the same area below the curve and the observed histogram. The excess at small distances is clearly due to physical associations and, because of the normalization chosen, is compensated by a deficiency at larger distances (between $`80\mathrm{}`$ and $`300\mathrm{}`$).
All the components closer than $`45\mathrm{}`$ (*i.e.* about three times the beam size) have been considered as possibly belonging to a unique double source. Triple sources are defined whenever one additional component is closer than $`45\mathrm{}`$ to (at least) one of the pair components. For multiple sources the same criterion is applied iteratively. Applying this distance constraint we expect that $`20\%`$ of the pairs are random superpositions.
The flux ratio distribution between the pair components has a large spread at all distances (see Fig. 6). To reduce the contamination we have discarded all the pairs with flux ratio larger than a factor 10. For triples and multiple sources the probable core is not considered when computing the flux ratios.
A few departures from the adopted criteria are present in the catalogue. For example the triple source ATESP J005620-394145 and the double source ATESP J011029-393253 have $`d45\mathrm{}`$ but do not satisfy the flux ratio constraint. All exceptions are based on source geometry considerations and/or the analysis of the source field.
In order to increase the multiple sources’ sub-sample completeness, we added 31 sources with distances $`45\mathrm{}<d<150\mathrm{}`$, which show clear signs of physical associations between their components (see Fig. 7 for some examples). No flux ratio constraints have been applied to such sources. In Fig. 6 are shown the flux ratios for all the pairs in the final sample of multiple sources (filled circles).
As a final result we have 189 multiple sources: 168 doubles, 19 triples and 2 sources with four components. As a consequence, the initial list of 3172 radio components results in a catalogue of 2960 distinct radio sources.
### 3.3 Non-Gaussian Sources
In the final catalogue we have 23 non-Gaussian sources whose parameters have been defined as discussed in Sect. 2.3. In particular we notice that positions refer to peak positions, which, for non-Gaussian sources does not necessarily correspond to the position of the core. We also notice that we can have non-Gaussian components in multiple sources. Some examples of single and multiple non Gaussian sources are shown in Fig. 7.
### 3.4 The Catalogue Format
The electronic version of the full radio catalog is available through the ATESP page at http://www.ira.bo.cnr.it. Its first page is shown as an example in Table 1. The source catalogue is sorted on right ascension. The format is the following:
Column (1) - Source IAU name. Different components of multiple sources are labeled ‘A’, ‘B’, etc.
Column (2) and (3) - Source position: Right Ascension and Declination (J2000).
Column (4) and (5) - Source peak ($`S_{\mathrm{peak}}`$) and integrated ($`S_{\mathrm{total}}`$) flux densities in mJy (Baars et al. Baarsetal77 (1977) scale). The flux densities are not corrected for the systematic effects discussed in Sect. 4.2.
Column (6) and (7) - Intrinsic (deconvolved from the beam) source angular size. Full width half maximum of the major ($`\mathrm{\Theta }_{\mathrm{maj}}`$) and minor ($`\mathrm{\Theta }_{\mathrm{min}}`$) axes in arcsec. Zero values refer to unresolved sources (see Sect. 3.1 for more details).
Column (8) - Source position angle (P.A., measured N through E) for the major axis in degrees.
Column (9) - Flag indicating the fitting procedure and parameterization adopted for the source or source component (see Sects. 2.3 and 3.2). ‘S’ refers to Gaussian fits. ‘S\*’ refers to poor Gaussian fits. ‘E’ refers to non-Gaussian sources. ‘M’ refers to multiple sources (see below).
The parameters listed for non-Gaussian sources are defined as discussed in Sect. 2.3.
For multiple sources we list all the components (labeled ‘A’, ‘B’, etc.) preceded by a line (flagged ‘M’) giving the position of the radio centroid, total flux density and overall angular size of the source. Source positions have been defined as the flux-weighted average position of all the components (source centroid). For sources with more than two components the centroid position has been replaced with the core position whenever the core is clearly recognizable.
Integrated total source flux densities are computed by summing all the component integrated fluxes.
The total source angular size is defined as *las* (see Sect. 2.3) and it is computed as the maximum distance between the source components.
## 4 Errors in the Source Parameters
Parameter uncertainties are the quadratic sum of two independent terms: the calibration errors, which dominate at high signal-to-noise ratios, and the internal errors, due to the presence of noise in the maps. The latter dominate at low signal-to-noise ratios.
In the following sections we discuss the parameter internal accuracy of our source catalogue. Master equations for total rms error derivation, with estimates of the calibration terms are reported in Appendix A.
### 4.1 Internal accuracy
In order to quantify the internal errors we produced a one square degree residual map by removing all the sources detected above $`5\sigma `$ in the radio mosaic fld20to25. We performed a set of Montecarlo simulations by injecting Gaussian sources in the residual map at random positions and re-extracting them using the same detection algorithm used for the survey (IMSAD). The Montecarlo simulations were performed by injecting samples of 30 sources at fixed flux and intrinsic angular size. We sampled peak fluxes between $`5\sigma `$ and $`50\sigma `$ and intrinsic angular sizes (FWHM major axis) between $`4\mathrm{}`$ and $`20\mathrm{}`$. Intrinsic sizes were convolved with the synthesized beam ($`8.5\mathrm{}\times 16.8\mathrm{}`$ for mosaic fld20to25) before injecting the source in the residual map.
The comparison between the input parameters and the ones provided by IMSAD permitted an estimate of the internal accuracy of the parameters as a function of source flux and intrinsic angular size. In particular we could test the accuracy of flux densities, positions and angular sizes and estimate both the efficiency and the accuracy of the deconvolution algorithm.
#### 4.1.1 Flux Densities and Source Sizes
The flux density and fitted angular size errors for point sources are shown in in Figs. 8 and 9 where we plot the ratio of the parameter value found by IMSAD (output) over the injected one (input), as a function of the signal-to-noise ratio.
We notice that mean values very far from 1 could indicate the presence of systematic effects in the parameter measure. The presence of such systematic effects is clearly present for peak flux densities in the faintest bins (see Fig. 8). This is the expected effect of the noise on the catalogue completeness at the extraction threshold. Due to its Gaussian distribution whenever an injected source falls on a noise dip, either the source flux is underestimated or the source goes undetected. This produces an incompleteness in the faintest bins. As a consequence, the measured fluxes are biased toward higher values in the incomplete bins, beacause only sources that fall on noise peaks have been detected and measured. We notice that the mean values found for $`S_{\mathrm{output}}/S_{\mathrm{input}}`$ are in good agreement with the ones expected taking into account such an effect (see dashed line). It is worth pointing out that our catalogue is only slightly affected by this effect because the detection threshold ($`4.5\sigma `$) is much lower than the $`6\sigma `$-threshold chosen for the catalogue (indicated by the vertical solid line in Fig. 8): at $`S_{\mathrm{peak}}6\sigma `$ we expect flux over-estimations $`5\%`$.
Some systematic effects appear to be present also for the source size at $`5\sigma `$: the major and minor axes tend to be respectively under- and over-estimated (see Fig. 9). Such effects disappear at $`S_{\mathrm{peak}}6\sigma `$ (ATESP cut-off). For both the flux densities and the source axes, the rms values measured are in very good agreement with the ones proposed by Condon (Condon97 (1997)) for elliptical Gaussian fitting procedures (for details see Appendix A):
$`\sigma (S_{\mathrm{peak}})/S_{\mathrm{peak}}`$ $`=`$ $`0.93\left({\displaystyle \frac{S_{\mathrm{peak}}}{\sigma }}\right)^1`$ (3)
$`\sigma (\theta _{\mathrm{maj}})/\theta _{\mathrm{maj}}`$ $`=`$ $`1.24\left({\displaystyle \frac{S_{\mathrm{peak}}}{\sigma }}\right)^1`$ (4)
$`\sigma (\theta _{\mathrm{min}})/\theta _{\mathrm{min}}`$ $`=`$ $`0.69\left({\displaystyle \frac{S_{\mathrm{peak}}}{\sigma }}\right)^1`$ (5)
where we have applied Eqs. (21) and (41) in Condon (Condon97 (1997)) to the case of ATESP point sources (see dotted lines in Figs. 8 and 9).
The fact that a source is extended does not affect the internal accuracy of the fitting algorithm for both the peak flux density and the source axes. In other words the errors quoted for point sources apply to extended sources as well.
However, this is not true for the deconvolution algorithm. The errors for the deconvolved source axes depend on both the source flux and intrinsic angular size. The higher the flux and the larger the source, the smaller the error. In particular, at 1 mJy ($`12\sigma `$) the errors are in the range $`35\%`$$`10\%`$ for angular sizes in the range $`6\mathrm{}`$$`20\mathrm{}`$. For fluxes $`>50\sigma `$ the errors are always $`<10\%`$. Deconvolved angular sizes are unreliable for very faint sources ($`56\sigma `$), where only a very small fraction of sources can be deconvolved. The deconvolution efficiency increases with the source flux. In particular, the fraction of deconvolved sources with intrinsic dimension $`4\mathrm{}`$ never reaches $`100\%`$: it goes from $`3\%`$ at the lowest fluxes, to $`15\%`$ at 1 mJy, to $`50\%`$ at the highest fluxes. We therefore can assume that $`4\mathrm{}`$ is a critical value for deconvolution at the ATESP resolution, and that ATESP sources with intrinsic sizes $`4\mathrm{}`$ are to be considered unresolved.
#### 4.1.2 Source Positions
The positional accuracy for point sources is shown in Fig. 10, where we plot the difference ($`\mathrm{\Delta }\alpha `$ and $`\mathrm{\Delta }\delta `$) between the position found by IMSAD (output) and the injected one (input), as a function of flux. No systematic effects are present and the rms values are in agreement with the ones expected for point sources (Condon Condon98 (1998), for details see Appendix A):
$`\sigma _\alpha `$ $``$ $`{\displaystyle \frac{b_{\mathrm{min}}}{2}}\left({\displaystyle \frac{S_{\mathrm{peak}}}{\sigma }}\right)^1`$ (6)
$`\sigma _\delta `$ $``$ $`{\displaystyle \frac{b_{\mathrm{maj}}}{2}}\left({\displaystyle \frac{S_{\mathrm{peak}}}{\sigma }}\right)^1`$ (7)
where we have assumed $`b_{\mathrm{min}}=8\mathrm{}`$ and $`b_{\mathrm{maj}}=14\mathrm{}`$, the average synthesized beam values of the ATESP survey (see dotted lines in Fig. 10). The positional accuracy of ATESP sources is therefore $`1\mathrm{}`$ at the limit of the survey ($`6\sigma `$), decreasing to $`0.5\mathrm{}`$ at $`12\sigma `$ ($`1`$ mJy) and to $`0.1\mathrm{}`$ at 50$`\sigma `$.
### 4.2 Systematic Effects
Two systematic effects are to be taken into account when dealing with ATESP flux densities, the clean bias and the bandwidth smearing effect. Clean bias has been extensively discussed in paper I of this series (see also Appendix B at the end of this paper). It is responsible for flux density under-estimations of the order of $`1020\%`$ at the lowest flux levels ($`S<10\sigma `$) and gradually disappears going to higher fluxes (no effect for $`S50100\sigma `$).
The effect of bandwidth smearing is well-known. It reduces the peak flux density of a source, correspondingly increasing the source size in radial direction. Integrated flux densities are therefore not affected.
The bandwidth smearing effect increases with the angular distance ($`d`$) from the the pointing center of phase and depends on the passband width ($`\mathrm{\Delta }\nu `$), the observing frequency ($`\nu `$) and the synthesized beam FWHM width ($`\theta _b`$). The particular functional form that describes the bandwidth smearing is determined by the beam and the passband shapes. It can be demonstrated, though, that the results obtained are not critically dependent on the particular functional form adopted (e.g. Bridle & Schwab Bridle89 (1989)).
In the simplest case of Gaussian beam and passband shapes, the bandwidth smearing effect can be described by the equation (see Eq. 12 in Condon et al. Condonetal98 (1998)):
$$\frac{S_{\mathrm{peak}}}{S_{\mathrm{peak}}^0}=\frac{1}{\sqrt{1+\frac{2\mathrm{ln}2}{3}\left(\frac{\mathrm{\Delta }\nu }{\nu }\frac{d}{\theta _b}\right)^2}}$$
(8)
where the $`\frac{S_{\mathrm{peak}}}{S_{\mathrm{peak}}^0}`$ ratio represents the attenuation on peak flux densities by respect to the unsmeared ($`d=0`$) source peak value.
We have estimated the actual smearing radial attenuation on ATESP single fields, by measuring $`S_{\mathrm{peak}}`$ for a strong source (ATESP J233758-401025) detected in eight contiguous ATESP fields (corrected for the primary beam attenuation) at increasing distance from the field center (full circles in Fig. 11, top panel). The data were then fitted using Eq. 8 by setting $`\nu =1.4`$ GHz and $`\theta _b11\mathrm{}`$ (from ($`b_{\mathrm{maj}}+b_{\mathrm{min}})/2=(14\mathrm{}+8\mathrm{})/2`$), as for the ATESP data. The best fit (dashed line in Fig. 11, top panel) gives $`S_{\mathrm{peak}}^0=131`$ mJy and an effective bandwidth $`\mathrm{\Delta }\nu =9`$ MHz (in very good agreement with the nominal channel width $`\mathrm{\Delta }\nu =8`$ MHz (see Sect. 5.2 in paper I). As expected, the measured integrated flux density (stars in the same plot) remains constant.
The mosaicing technique consists in a weighted linear combination of all the single fields in a larger mosaiced image (see Eq. 1 in paper I). This means that, given single fields of size $`1800\times 1800`$ pixels, source flux measures at distances as large as $`35\mathrm{}`$ from field centers are still used to produce the final mosaic (even if with small weights). As a consequence, the radial dependence of bandwidth smearing tends to cancel out.
For instance, since ATESP pointings are organized in a $`20\mathrm{}`$ spacing rectangular grid, a source located at the center of phase of one field ($`d=0`$) is measured also at $`d=20\mathrm{}`$ in the 4 contiguous E, W, S and N fields and at $`d=20\sqrt{2}28\mathrm{}`$ in the other 4 *diagonally* contiguous fields. Using Eq. (1) of paper I, we can estimate a 4% smearing attenuation for the mosaic peak flux of that source. In the same way we can estimate indicative values for mosaic smearing attenuations as a function of $`d_{\mathrm{min}}`$, defined as the distance to the closest field center (see dotted line in Fig. 12). We notice that actual attenuations vary from source to source depending on the actual position of the source in the mosaic.
From Fig. 12 we can see that at small $`d_{\mathrm{min}}`$ mosaic smearing is much worse than single field’s one (indicated by the solid line). The discrepancy becomes smaller going to larger distances and disappears at $`d_{\mathrm{min}}14\mathrm{}`$, which represents the maximum distance to the closest field center for ATESP sources. This maximum $`d_{\mathrm{min}}`$ value gives an upper limit of $`6\%`$ to mosaic smearing attenuations.
The expected mosaic attenuations have been compared to the ones obtained directly estimating the smearing from the source catalogue. As already noticed (Sect. 3.1), a ratio $`S_{\mathrm{peak}}/S_{\mathrm{total}}<1`$, is purely determined, in case of point sources and in absence of flux measurement errors, by the bandwidth smearing effect, which systematically attenuates the source peak flux, leaving the integrated flux unchanged.
We have then considered all the unresolved ($`\mathrm{\Theta }_{\mathrm{maj}}=0`$) ATESP sources with $`S_{\mathrm{peak}}>2`$ mJy and we have plotted the average values of the $`S_{\mathrm{peak}}/S_{\mathrm{total}}`$ ratio in different distance intervals (full dots in Fig. 12). The 2 mJy threshold ($`25\sigma `$) was chosen in order to find a compromise between statistics and flux measure accuracy. The average flux ratios plotted are in very good agreement with the expected ones, especially when considering that the most reliable measures are the intermediate distance ones, where a larger number of sources can be summed. In general we can conclude that on average smearing attenuations are $`5\%`$ and do not depend on the actual position of the source in the mosaics. This result also confirms the $`5\%`$ estimate drawn from Fig. 4. We finally point out that smearing will affect to some extent also source sizes and source coordinates.
### 4.3 Comparison with External Data
To determine the quality of ATESP source parameters we have made comparisons with external data. Unfortunately in the region covered by the ATESP, there are not many 1.4 GHz data available. The only existing 1.4 GHz radio survey is the NVSS (Condon et al. Condonetal98 (1998)), which covers about half of the region covered by the ATESP survey ($`\delta >40\mathrm{°}`$) with a point source detection limit $`S_{lim}2.5`$ mJy.
The NVSS has a poor spatial resolution ($`45\mathrm{}`$ FWHM beam width) compared to ATESP and this introduces large uncertainties in the comparison, especially for astrometry. To test the positional accuracy we have therefore used data at other frequencies as well. In particular we have used VLBI sources extracted from the list of the standard calibrators at the ATCA and the catalogue of PMN compact sources with measured ATCA positions (Wright et al. Wright97 (1997)).
#### 4.3.1 Flux Densities
In order to estimate the quality of the ATESP flux densities we have compared ATESP with NVSS. To minimize the uncertainties due to the much poorer NVSS resolution we should in principle consider only point-like ATESP sources. Nevertheless, in order to increase the statistics at high fluxes ($`S>10`$ mJy), we decided to include extended ATESP sources as well. Another source of uncertainty in the comparison is due to the fact that in the NVSS distinct sources closer together than $`50\mathrm{}`$ will be only marginally separated. To avoid this problem we have restricted the comparison to bright ($`S_{\mathrm{peak}}>1`$ mJy) *isolated* ATESP sources: we have discarded all multi-component sources (as defined in Sect. 3.2) and all single-component sources whose nearest neighbor is at a distance $`100\mathrm{}`$ (as in the comparison between the FIRST and the NVSS by White et al. White97 (1997)). We have not used isolated ATESP sources fainter than 1 mJy because we have noticed that there are several cases where NVSS point sources are resolved in two distinct objects in the ATESP images, only one being listed in the ATESP catalogue (i.e. the other one has $`S_{\mathrm{peak}}<6\sigma `$).
In Fig. 13 we have plotted the NVSS against the ATESP flux density for the 443 $`S_{\mathrm{peak}}>1`$ mJy ATESP sources identified (sources within the $`3\sigma `$ confidence circle in Fig. 14). We have used integrated fluxes for the sources which appear extended at the ATESP resolution (full circles) and peak fluxes for the unresolved ones (dots). Also indicated are the $`90\%`$ confidence limits (dashed curves), drawn taking into account both the NVSS and the ATESP errors in the flux measure. In drawing the upper line we have also taken into account an average correction for the systematic under-estimation of ATESP fluxes due to clean bias and bandwidth smearing (see Sect. 4.2). The provided NVSS fluxes are already corrected for any systematic effect.
The scatter plot shows that at high fluxes the ATESP flux scale agrees with the NVSS one within a few percents ($`3\%`$). This gives an upper limit to calibration errors and/or resolution effects at high fluxes. Going to fainter fluxes the discrepancies between the ATESP and the NVSS fluxes become larger, reaching deviations as high as a factor of 2 at the faintest levels. ATESP fluxes tend to be lower than NVSS fluxes. This could be, at least partly, due to resolution effects. Such effects can be estimated from the comparison between NVSS and ATESP sources itself and from theoretical considerations. Assuming the source integral angular size distribution provided by Windhorst et al. (Windhorst90 (1990)) we have that at the NVSS limit ($`S2.5`$ mJy) about 40% of the sources can be resolved by the ATESP synthesized beam (intrinsic angular sizes $`4\mathrm{}`$). This fraction goes up to $`50\%`$ at $`S10`$ mJy.
On the other hand, we point out that close to the NVSS catalogue cut-off, incompleteness could affect NVSS fluxes. For instance we have noticed that below 5 mJy, there are several cases where the flux given in the NVSS catalogue is overestimated with respect to the one measured in the NVSS image (even taking into account the applied corrections).
#### 4.3.2 Astrometry
The region covered by the ATESP survey contains two VLBI sources from the ATCA calibrator catalogue: 2227-399 and 0104-408. These sources were not used to calibrate our data and therefore provide an independent check of our ATCA positions. The offset between ATESP and VLBI positions (ATESP–VLBI) for the first and the second source respectively are: $`\mathrm{\Delta }\alpha =0.277\mathrm{}`$ and $`0.023\mathrm{}`$; $`\mathrm{\Delta }\delta =+0.239\mathrm{}`$ and $`0.172\mathrm{}`$. Such offsets indicate that the uncertainty in the astrometry should be within a fraction of arcsec. Obviously, we cannot exclude the presence of systematic effects, but an analysis of the ATESP–NVSS positional offsets gives $`\mathrm{\Delta }\alpha _{\mathrm{med}}=0.115\mathrm{}`$ and $`\mathrm{\Delta }\delta _{\mathrm{med}}=0.8\mathrm{}`$ (see Fig. 14).
A more precise comparison could be obtained by using the PMN sources with ATCA position measurements available. Unfortunately the number of such PMN sources in the region covered by the ATESP survey is very small: we found only 12 identifications. Using the 4.8 GHz positions for the PMN sources, we derived (ATESP–ATPMN) $`\mathrm{\Delta }\alpha _{\mathrm{med}}=0.115\mathrm{}`$ and $`\mathrm{\Delta }\delta _{\mathrm{med}}=0.3\mathrm{}`$.
We can conclude that all comparisons give consistent results and that the astrometry is accurate within a small fraction of an arcsec. Also systematic offsets, if present, should be very small.
###### Acknowledgements.
We acknowledge R. Fanti for reading and commenting on an earlier version of this manuscript.
The authors aknowledge the Italian Ministry for University and Scientific Research (MURST) for partial financial support (grant Cofin 98-02-32). This project was undertaken under the CSIRO/CNR Collaboration programme. The Australia Telescope is funded by the Commonwealth of Australia for operation as a National Facility managed by CSIRO.
## Appendix A Master Equations for Error Derivation
As discussed in Sect. 4, internal errors for the ATESP source parameters are well described by Condon (Condon97 (1997)) equations of error propagation derived for two-dimensional elliptical Gaussian fits in presence of Gaussian noise. In order to get the total rms error on each parameter, a calibration term should be quadratically added. Using Eqs. (21) and (41) in Condon (Condon97 (1997)), total percentage errors for flux densities ($`\frac{\sigma (S_{\mathrm{peak}})}{S_{\mathrm{peak}}}`$) and fitted axes ($`\frac{\sigma (\theta _{\mathrm{maj}})}{\theta _{\mathrm{maj}}}`$, $`\frac{\sigma (\theta _{\mathrm{min}})}{\theta _{\mathrm{min}}}`$) can be calculated from:
$$\sqrt{\frac{2}{\rho ^2}+ϵ^2}$$
(9)
where $`ϵ`$ is the calibration term and the effective signal-to-noise ratio, $`\rho `$, is given by:
$$\rho ^2=\frac{\theta _{\mathrm{maj}}\theta _{\mathrm{min}}}{4\theta _\mathrm{N}^2}\left[1+\left(\frac{\theta _\mathrm{N}}{\theta _{\mathrm{maj}}}\right)^2\right]^{\alpha _\mathrm{M}}\left[1+\left(\frac{\theta _\mathrm{N}}{\theta _{\mathrm{min}}}\right)^2\right]^{\alpha _\mathrm{m}}\frac{S_{\mathrm{peak}}^2}{\sigma ^2}$$
(10)
where $`\sigma `$ is the image noise ($`79`$ $`\mu `$Jy on average for ATESP images), $`\theta _\mathrm{N}`$ is the FWHM of the Gaussian correlation length of the image noise (assumed $``$ FWHM of the synthesized beam) and the exponents are:
$`\alpha _\mathrm{M}=3/2\mathrm{and}`$ $`\alpha _\mathrm{m}=3/2`$ $`\mathrm{for}\sigma (S_{\mathrm{peak}})`$ (11)
$`\alpha _\mathrm{M}=5/2\mathrm{and}`$ $`\alpha _\mathrm{m}=1/2`$ $`\mathrm{for}\sigma (\theta _{\mathrm{maj}})`$ (12)
$`\alpha _\mathrm{M}=1/2\mathrm{and}`$ $`\alpha _\mathrm{m}=5/2`$ $`\mathrm{for}\sigma (\theta _{\mathrm{min}}).`$ (13)
Similar equations hold for position rms errors (Condon Condon97 (1997), Condon et al. Condonetal98 (1998)):
$`\sigma ^2(\alpha )=`$ $`ϵ_\alpha ^2+\sigma ^2(x_0)\mathrm{sin}^2(\mathrm{P}.\mathrm{A}.)+\sigma ^2(y_0)\mathrm{cos}^2(\mathrm{P}.\mathrm{A}.)`$ (14)
$`\sigma ^2(\delta )=`$ $`ϵ_\delta ^2+\sigma ^2(x_0)\mathrm{cos}^2(\mathrm{P}.\mathrm{A}.)+\sigma ^2(y_0)\mathrm{sin}^2(\mathrm{P}.\mathrm{A}.)`$ (15)
where $`\sigma ^2(x_0)=\sigma ^2(\theta _{\mathrm{maj}})/8\mathrm{ln}2`$ and $`\sigma ^2(y_0)=\sigma ^2(\theta _{\mathrm{min}})/8\mathrm{ln}2`$ represent the rms lengths of the major and minor axes respectively.
Calibration terms are in general estimated from comparison with consistent external data of better accuracy than the one tested. Unfortunately there are no such data available in the region of sky covered by the ATESP survey. Nevertheless, from our typical flux and phase calibration errors, we estimate calibration terms of about $`510\%`$ for both flux densities and source sizes.
As a caveat we remind (see discussion in Paper I) that the 500 m baseline cutoff applied to our data makes the ATESP survey progressively insensitive to sources larger than $`30\mathrm{}`$: assuming a Gaussian shape, only $`50\%`$ of the flux for a $`30\mathrm{}`$ large source would appear in the ATESP images. It is important to have this in mind when dealing with flux densities and source sizes of the largest ATESP sources.
Right ascension and declination calibration terms have been estimated from the astrometry results reported in Sect. 4.3.2. As already discussed, the ATESP astrometry can be considered accurate within a small fraction of an arcsec, even though the scarcity of (accurate) external data available in the ATESP region makes it difficult to quantify this statement. Nevertheless from the rms dispersion of the median offsets found between ATESP and the external comparison samples (see Sect. 4.3.2) we can tentatively estimate $`ϵ_\alpha 0.1\mathrm{}`$ and $`ϵ_\delta 0.4\mathrm{}`$.
It can be easily demonstrated that the master equations (9), (14) and (15) reduce to Eqs. (3$``$ (7) in Sect. 4 (where the calibration term is neglected) in the case of ATESP point sources ($`\theta _{\mathrm{maj}}\times \theta _{\mathrm{min}}14\mathrm{}\times 8\mathrm{}`$, P.A.$`+2\mathrm{°}`$), assuming $`\theta _\mathrm{N}11\mathrm{}`$ (or $`\theta _\mathrm{N}^2\theta _{\mathrm{maj}}\theta _{\mathrm{min}}`$).
For a complete and detailed discussion of the error master equations of source parameters obtained through elliptical Gaussian fits we refer to Condon (Condon97 (1997)) and Condon et al. (Condonetal98 (1998)).
## Appendix B Flux Density Corrections for Systematic Effects
As already discussed in Sect. 4.2, two systematic effects are to be taken into account when dealing with ATESP flux densities, the clean bias and the bandwidth smearing effect.
The flux densities reported in the ATESP source catalogue are not corrected for such systematic effects. The corrected flux densities ($`S^{\mathrm{corr}}`$) can be computed as follows:
$$S^{\mathrm{corr}}=\frac{S^{\mathrm{meas}}}{k[a\mathrm{log}(S^{\mathrm{corr}}/\sigma )+b]}$$
(16)
where $`S^{\mathrm{meas}}`$ is the flux actually measured in the ATESP images (the one reported in the source catalogue). The parameter $`k`$ represents the smearing correction. It is set equal to 1 (i.e. no correction) when the equation is applied to integrated flux densities and $`<1`$ when dealing with peak flux densities. From the analysis reported in Sect. 4.2 we suggest to set $`k=0.95`$ ($`5\%`$ smearing effect).
The clean bias correction is taken into account by the term in the square brackets. As discussed in paper I, Sect. 5.3, the importance of the clean bias effect varies from mosaic to mosaic depending on the average number of clean components (cc’s). In particular we derived the values for the parameters $`a`$ and $`b`$ in three different mosaics representing the case of low (fld34to40, 1616 cc’s), intermediate (fld44to50, 2377 cc’s) and high (fld69to75, 3119 cc’s) average number of cc’s (see Table 4 of paper I).
In correcting the source fluxes for the clean bias, we suggest to set $`(a,b)=(0.09,0.85)`$ whenever the mosaic average number of cc’s is $`<2000`$ (low–cc’s case); $`(a,b)=(0.13,0.75)`$ whenever the mosaic cc’s average number is between 2000 and 3000 (intermediate–cc’s case); $`(a,b)=(0.16,0.67)`$ whenever the mosaic cc’s average number exceeds 3000 (high–cc’s case). The average number of clean components for each mosaic is reported in Table 2. In order to trace back the sources to the original mosaics, Table 2 lists also the right ascension range covered by each mosaic (indicated by the r.a. of the first and the last source in that mosaic).
The clean bias is a function of the source signal-to-noise ratio $`S^{\mathrm{corr}}/\sigma `$. Since the noise level is fairly uniform within each mosaic, it is possible to assume $`\sigma `$ equal to the mosaic average noise value ($`\sigma _{fit}`$ in Table 2, we refer to paper I for details on mosaic noise analysis and average noise value definition).
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# Large scale dynamos with ambipolar diffusion nonlinearity
## 1 Ambipolar diffusion as a toy nonlinearity
In this Letter we test and exploit the idea that the exact type of nonlinearity in the MHD equations is unessential as far as the nature of large scale field generation is concerned. At first glance this may seem rather surprising, especially if one pictures large scale field generation as the result of an inverse cascade process (Frisch et al. 1975, Pouquet et al. 1976). Like the direct cascade in Kolmogorov turbulence, the inverse cascade is accomplished by nonlinear interactions, suggesting that nonlinearity is important. However, a special type of inverse cascade is the strongly nonlocal inverse cascade process, which is usually referred to as the ‘alpha-effect’; see Moffatt (1978) and Krause & Rädler (1980). This effect exists already in linear (kinematic) theory.
Until recently it was unclear which, if any, of the two effects (inverse cascade in the local sense or the $`\alpha `$-effect) played the dominant role in large scale field generation as seen in simulations (e.g. Glatzmaier & Roberts 1995, Brandenburg et al. 1995, Ziegler & Rüdiger 2000) or in astrophysical bodies (stars, galaxies, accretion discs). A strong indication that it is actually the $`\alpha `$-effect (i.e. the strongly nonlocal inverse cascade) that is responsible for large scale field generation, comes from detailed analysis of recent three-dimensional simulations of forced isotropic non-mirror symmetric turbulence (Brandenburg 2000, hereafter B2000). In those simulations a strong and nearly force-free magnetic field was produced, and most of the energy supply to this field was found to come from the forcing scale of the turbulence.
In the absence of nonlinearity, however, the field seen in the simulations of B2000 became quickly swamped by magnetic fields at smaller scales. In that sense a purely kinematic large scale turbulent dynamo is impossible! Any hope for analytic progress is therefore slim. However, the model of Subramanian (1997, 1999) is an exception. Subramanian (1997; hereafter S97) extended the kinematic models of of Kazantsev (1968) and Vainshtein & Kitchatinov (1986) by including ambipolar diffusion (in the strong coupling approximation) as a nonlinearity. Under the common assumption that the velocity is delta-correlated in time, S97 derived a nonlinear equation for the evolution of the correlation functions of magnetic field and magnetic helicity. Although the models of Kazantsev (1968) and Novikov et al. (1983) are usually known to describe small-scale field generation, Subramanian (1999; hereafter S99) found that in the presence of fluid helicity there is the possibility of tunnelling of bound-states corresponding to small scales to unbounded states corresponding to large scale fields, which are force-free.
In this Letter we present numerical solutions to the closure model of S99. We stress that we do not advocate ambipolar diffusion (AD) as being dominant over the usual feedback from the Lorentz force in the momentum equation. Instead, our motivation is to establish a useful toy model to study effects of nonlinearity in dynamos. Our numerical solutions may provide guidance for further analytic treatment of these equations in parameter regimes otherwise inaccessible. We begin however by considering first solutions of the fully three-dimensional MHD equations in a periodic box using AD as the only nonlinearity.
## 2 Box simulations for a finite system
In this section we adopt the MHD equations for an isothermal compressible gas, driven by a given body force $`𝒇`$, in the presence of AD, but ignoring the Lorentz force
$$\frac{\mathrm{D}\mathrm{ln}\rho }{\mathrm{D}t}=\mathbf{}𝒖,$$
(1)
$$\frac{\mathrm{D}𝒖}{\mathrm{D}t}=c_\mathrm{s}^2\mathbf{}\mathrm{ln}\rho +\frac{\mu }{\rho }(^2𝒖+\frac{1}{3}\mathbf{}\mathbf{}𝒖)+𝒇,$$
(2)
$$\frac{𝑨}{t}=(𝒖+𝒖{}_{\mathrm{D}}{}^{})\times 𝑩\eta \mu _0𝑱,$$
(3)
where $`\mathrm{D}/\mathrm{D}t=/t+𝒖\mathbf{}`$ is the advective derivative, $`𝑩=\mathbf{}\times 𝑨`$ is the magnetic field, $`𝑱=\mathbf{}\times 𝑩/\mu _0`$ is the current density, and $`𝒇`$ is the random forcing function as specified in B2000. The nonlinear drift velocity $`𝒖_\mathrm{D}`$ due to AD can be written as $`𝒖{}_{\mathrm{D}}{}^{}=a𝑱\times 𝑩`$. We use nondimensional units where $`c_\mathrm{s}=k_1=\rho _0=\mu _0=1`$. Here, $`c_\mathrm{s}`$ is the sound speed, $`k_1`$ the smallest wavenumber of the box (so its size is $`2\pi `$), $`\rho _0`$ is the mean density, and $`\mu _0`$ is the vacuum permeability. Since AD is the only nonlinearity in Eq. (3) we can always normalize $`𝑩`$ such that $`a=1`$.
The model presented here is similar to Run 3 of B2000, where $`\mu =\eta =2\times 10^3`$. With a root-mean-square velocity of around 0.3 the magnetic Reynolds number based on the size of the box is around 1000. The forcing wavenumber $`k_\mathrm{f}`$ is chosen to be 5. In Fig. 1 we show a grey scale representation of a slice of the magnetic field and the current density at $`t=337`$. Note the presence of a large scale magnetic field that varies in the $`z`$-direction. In Fig. 2 we show the spectra of magnetic and kinetic energies. The peak of magnetic energy at $`k=1`$ shows the development of large scale magnetic fields. Further, the current density is concentrated into narrow filamentary structures, typical of AD (see Brandenburg & Zweibel 1994).
Unfortunately, the severity of the diffusive timestep limit, $`\delta t0.16\delta x^2/\eta _{\mathrm{AD}}`$, where $`\eta _{\mathrm{AD}}=a𝑩^2`$, prevented us from running much longer at high resolution ($`120^3`$ meshpoints). For $`60^3`$ meshpoints this limit is unimportant, and so we were able to run until $`t=900`$, a time when the large scale field was much more clearly defined. In the inset of Fig. 2, we show the evolution for such a case, but with a forcing at $`k_\mathrm{f}=10`$ (giving larger scale separation). Note again the peak of $`E_M`$ at $`k=1`$ and also the suppression of magnetic field at the next smaller scale, corresponding to $`k2`$. Both these features are very similar to the magnetic field evolution in the case with full Lorentz force and without AD (Figs 3 and 17 of B2000).
Our main conclusion from these results is first of all that large scale field generation works in spite of AD, contrary to earlier suggestions that AD might suppress the large scale dynamo process (Kulsrud & Anderson 1992). Secondly, AD provides a nonlinear saturation mechanism for the magnetic field at all scales, except for the scale of the box, where a force-free field develops for which $`𝒖_\mathrm{D}`$ vanishes. Like in the simulations of B2000 this provides a ‘self-cleaning’ mechanism, without which the field would be dominated by contributions from small scales.
Having established the close similarity between models with AD versus full Lorentz force as nonlinearity, we now move on to discuss the nonlinear closure model of S99 with AD as a ‘toy’ nonlinearity.
## 3 Closure model for an infinite system
Under the assumptions that the velocity is delta-correlated in time and the magnetic field is a gaussian random field S97 derived equations for the longitudinal correlation function $`M(r,t)`$ and the correlation function for magnetic helicity density, $`N(r,t)`$. The velocity is represented by a longitudinal correlation function $`T(r)`$ and a correlation function for the kinetic helicity density, $`C(r)`$. We change somewhat the notation of S99 and define the operators
$$\stackrel{~}{D}()=\frac{1}{r^4}\frac{}{r}(r^4),D()=\frac{}{r}(),$$
(4)
so the closure equations can be written as
$$\dot{M}=2\stackrel{~}{D}(\eta _\mathrm{T}DM)+2GM+4\alpha H,$$
(5)
$$\dot{N}=2\eta _\mathrm{T}H+\alpha M,$$
(6)
where $`H=\stackrel{~}{D}DN`$ is the correlation function of the current helicity, $`G=\stackrel{~}{D}DT`$ is the effective induction,
$$\alpha =\alpha _0(r)+4aH(0,t)$$
(7)
$$\eta _\mathrm{T}=\eta +\eta _0(r)+2aM(0,t)$$
(8)
are functions resembling the usual $`\alpha `$-effect and the total magnetic diffusivity. Here $`\alpha _0(r)=2[C(0)C(r)]`$ and $`\eta _0(r)=T(0)T(r)`$. Note that at large scales
$$\alpha _{\mathrm{}}\alpha (r\mathrm{})=\frac{1}{3}\tau 𝝎𝒖+\frac{1}{3}\tau _{\mathrm{AD}}𝑱𝑩/\rho _0$$
(9)
$$\eta _{\mathrm{}}\eta _\mathrm{T}(r\mathrm{})=\frac{1}{3}\tau 𝒖{}_{}{}^{2}+\frac{1}{3}\tau _{\mathrm{AD}}𝑩{}_{}{}^{2}/\mu _0\rho _0,$$
(10)
where $`\tau _{\mathrm{AD}}=2a\rho _0`$. Expression (9) is very similar to the $`\alpha `$-suppression formula first found by Pouquet et al. (1976). Here $`\alpha `$ and $`\eta _\mathrm{T}`$ are scale dependent (i.e. they are largest on large scales) and, in addition, both are affected by AD.
We construct $`T(r)`$ and $`C(r)`$ from an analytic approximation of the kinetic energy and helicity spectra, $`E_\mathrm{K}(k)`$ and $`H_\mathrm{K}(k)`$, respectively. Zero velocity at large scales means that $`E_K(k)k^4`$ for $`k0`$. At some wavenumber $`k=k_\mathrm{f}`$ the spectrum turns to a $`k^{5/3}`$ Kolmogorov spectrum, followed by an exponential cutoff, so we take
$$E_\mathrm{K}(k)=\frac{E_0(k/k_\mathrm{f})^4}{1+(k/k_\mathrm{f})^{17/3}}\mathrm{exp}(k/k_\mathrm{d}).$$
(11)
We use parameters representative of the simulations of B2000, so $`E_0=0.01`$, $`k_\mathrm{f}=5`$ and $`k_\mathrm{d}=25`$. Like in B2000 we assume the turbulence fully helical, so $`H_\mathrm{K}=2kE_\mathrm{K}`$ (e.g. Moffatt 1978). The correlation functions $`T(r)`$ and $`C(r)`$ are then obtained via
$$T(r)=\frac{2}{\tau }_0^{\mathrm{}}E(k)\frac{j_1(kr)}{kr}dk(E(k)),$$
(12)
and $`C(r)=(F(k))/4`$, where $`j_1(x)=(\mathrm{sin}xx\mathrm{cos}x)/x^2`$ and $`\tau `$ is the correlation time. (We use $`\tau =4`$, representative of the kinematic stage of Run 3 of B2000.)
We solve Eqs. (5) and (6) using second order finite differences and a third order time step on a uniform mesh in $`0<x<L`$ with up to 10,000 meshpoints and $`L=10\pi `$, which is large enough so that the outer boundary does not matter. In the absence of helicity, $`C=0`$, and without nonlinearity, $`a=0`$, we recover the model of Novikov et al. (1983). The critical magnetic Reynolds number based on the forcing scale is around 60. In the presence of helicity this critical Reynolds number decreases, confirming the general result that helicity promotes dynamo action (cf. Kim & Hughes 1997, S99). In the presence of nonlinearity the exponential growth of the magnetic field terminates when the magnetic energy becomes large. After that point the magnetic energy continues however to increase nearly linearly. Unlike the case of the periodic box (Sect. 2) the magnetic field can here extend to larger and larger scales; see Fig. 3. The corresponding magnetic energy spectra,
$$E_\mathrm{M}(k,t)=\frac{1}{\pi }_0^LM(r,t)(kr)^3j_1(kr)dk,$$
(13)
are shown in Fig. 4.
The resulting magnetic field is strongly helical and the magnetic helicity spectra (not shown) satisfy $`H_\mathrm{M}\begin{array}{c}<\\ \end{array}(2/k)E_\mathrm{M}`$. The development of a helicity wave travelling towards smaller and smaller $`k`$, as seen in Fig. 4, is in agreement with the closure model of Pouquet et al. (1976). In the following we shall address the question of whether or not the growth of this large scale field depends on the magnetic Reynolds number (as in B2000). We have checked that to a very good approximation the wavenumber of the peak is given by
$$k_{\mathrm{peak}}(t)\alpha _{\mathrm{}}(t)/\eta _{\mathrm{}}(t).$$
(14)
This result is familiar from mean-field dynamo theory (see also S99) and is consistent with simulations (B2000, section 3.5). Note that here $`k_{\mathrm{peak}}`$ decreases with time because $`\alpha _{\mathrm{}}`$ tends to a finite limit and $`\eta _{\mathrm{}}`$ increases. (This is not the case in the box calculations where $`k_{\mathrm{peak}}2\pi /L`$.)
## 4 Resistively limited growth on large scales
In an unbounded system the magnetic helicity, $`𝑨𝑩=6N(0,t)`$, can only change if there is microscopic magnetic diffusion and finite current helicity, $`𝑱𝑩=6H(0,t)`$,
$$\mathrm{d}𝑨𝑩/\mathrm{d}t=2\eta 𝑱𝑩.$$
(15)
The closure model of S97 and S99 also satisfies this constraint. (Note that ambipolar and/or turbulent diffusion do not enter!) As explained in B2000, this constraint limits the speed at which the large scale field can grow, but not its final amplitude. One way to relax this constraint is if there is a flux of helicity through open boundaries (Blackman & Field 2000, Kleeorin et al. 2000), which may be important in astrophysical bodies with boundaries. Here, however, we consider an infinite system.
In Fig. 5 we show that, after some time $`t=t_\mathrm{s}`$, $`𝑱𝑩`$ reaches a finite value. This value increases somewhat as $`\eta `$ is decreased. In all cases, however, $`𝑱𝑩`$ stays below $`𝝎𝒖(2\tau /a)`$, so that $`|\alpha _{\mathrm{}}|`$ remains finite; see (9). A constant $`𝑱𝑩`$ implies that $`𝑨𝑩`$ grows linearly at a rate proportional to $`\eta `$. However, since the large scale field is helical, and since most of the magnetic energy is by now (after $`t=t_\mathrm{s}`$) in the large scales, the magnetic energy is proportional to $`𝑩{}_{}{}^{2}k_{\mathrm{peak}}𝑨𝑩`$, and can therefore only continue to grow at a resistively limited rate, see Fig. 5.
## 5 Conclusions
Our results have shown that ambipolar diffusion (AD) provides a useful model for nonlinearity, enabling analytic (or semi-analytic) progress to be made in understanding nonlinear dynamos. There are two key features that are shared both by this model and by the full MHD equations: (i) large scale fields are the result of a nonlocal inverse cascade as described by the $`\alpha `$-effect, and (ii) after some initial saturation phase the large scale field continues to grow at a rate limited by magnetic diffusion. We reiterate that in astrophysical bodies the presence of open boundaries may relax the helicity constraint. Furthermore, the presence of large scale shear or differential rotation provides a means of amplifying toroidal magnetic fields quite independently of magnetic helicity, but this still requires poloidal fields for which the above conclusions hold.
###### Acknowledgements.
KS thanks Nordita for hospitality during the course of this work. Use of the PPARC supported supercomputers in St Andrews and Leicester is acknowledged.
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# Replica symmetry breaking in the minority game
## 1 Introduction
The minority game has drawn much attention recently as a toy model of a market . In the simplest possible case, when no public information is present, its definition is fairly simple. At each time step, $`N`$ players have to choose between two actions, such as buying a certain stock or selling it. Those who end up in the minority side win. This mechanism can be obtained by abstracting the well known law of supply-and-demand. When the majority of traders is buying a certain asset it is convenient to be a seller, for prices are likely to be high, and viceversa. The minority side has an advantage.
The full complexity of the model arises in the presence of public information, which is modeled by the occurrence of one of $`P`$ events representing, e.g., some political news or a price change. In the minority game agents resort to choice rules or information-processing devices – called strategies henceforth – which suggest them whether to “buy” or “sell” given the information they have received. Then each player acts according to the suggestion of his best performing strategy. This mechanism allows to tackle the central problem one faces when trying to understand the collective behaviour of systems of heterogeneous agents interacting under strategic interdependence (as in markets), that is, how agents react to public information (e.g., political news or price changes) and the feedback effects that these reactions have on public information.
Early numerical simulations have shown a remarkably rich behaviour where both cooperativity and crowd effects arise. Much emphasis was initially put on the emergence of a cooperative phase in the stationary state as compared to the reference situation of random agents – i.e., agents who toss a coin to decide which action to take. Agents in the minority game are able to coordinate their actions and reduce the global waste of resources below the level corresponding to random agents.
Later work has revealed that the agents’ adaptive dynamics minimizes a global function, related to market predictability, and that the system undergoes a phase transition between an asymmetric and a symmetric, unpredictable market, as the ratio $`\alpha =P/N`$ decreases. A full characterization of the model’s behaviour for $`N\mathrm{}`$ was derived studying the minima of the global function by the replica method <sup>3</sup><sup>3</sup>3 It should be mentioned at this point that the replica approach fails to describe the system’s behaviour in a certain range of the parameters, as discussed in Refs . This point will be made more precise later in the text.. It was also realized that agents can greatly improve their performance and global efficiency if they account for their own impact on the market . In this case the steady state is a *Nash equilibrium*, that is a configuration where no agent can improve his performance by changing his behavior if others stick to theirs. Also in this case the dynamics is related asymptotically to the minima of a global function, and hence statistical mechanics again allows one to describe in detail the stationary state. However, for Nash equilibria replica symmetry breaking (RSB) occurs. This makes the replica symmetric calculation of Refs only an approximation.
In this paper we move the first steps towards a complete characterisation of the set of NE of the minority game. Our analysis will be slightly more general, for we shall be able to embody the cooperative state as well. First, we shall briefly outline the replica approach, showing that the replica-symmetric solution has a limited validity due to entropy arguments. Then we compute the number of Nash equilibria, following Ref. , and show that there are exponentially many (in $`N`$). By analogy with the de Almeida-Thouless stability analysis of the SK spin glass model , we find the phase transition line (AT line) separating the replica symmetric from the RSB phase. Finally we shall break the replica symmetry and study the solution in the one-step RSB approximation (1RSB). Probably the exact solution requires infinitely many steps of RSB, but 1RSB provides already an extremely close agreement with numerical results. Finally, we will be able to draw complete phase diagram of the model within 1RSB.
Our discussion will focus on the statistical mechanical properties of the model. The economic and game-theoretic aspects of the model, which are discussed in detail elsewhere , will only be described briefly.
## 2 The model
### 2.1 Basic definitions
The essential ingredients of the minority game are:
* $`N`$ players, labeled by the index $`i`$;
* for each player $`i`$ a strategy variable $`s_i\{\pm 1\}`$, saying which strategy ($`+1`$ or $`1`$) player $`i`$ is adopting (we restrict ourselves to the case where players have two strategies each);
* $`P`$ different information patterns, labeled by the index $`\mu `$.
At each time step $`t`$ all players receive the same information $`\mu `$, drawn at random with equal probability in $`\{1,\mathrm{},P\}`$ . Strategies $`s`$ are the label of information processing devices, that suggest an action as a value of a binary variable (like “buy” or “sell”) upon receiving information $`\mu `$:
$$s:\{1,\mathrm{},P\}\mu a_{i,s}^\mu \{\pm 1\}.$$
(1)
Two such strategies are assigned to each agent and they are drawn at random and independently for each agent, from the set of all $`2^P`$ such functions. In practice, the $`a_{i,s}^\mu `$ play the role of quenched disorder, analogous to the random couplings $`\{J_{ij}\}`$ in spin glass models.
It is convenient to make the dependence of $`a_{i,s}^\mu `$ on $`s`$ explicit by introducing auxiliary random variables $`\omega _i^\mu `$ and $`\xi _i^\mu `$ such that
$$a_{i,s}^\mu =\omega _i^\mu +s\xi _i^\mu .$$
(2)
Clearly, both $`\xi _i^\mu `$ and $`\omega _i^\mu `$ take on values in $`\{0,\pm 1\}`$ but they are not independent.
The payoff to player $`i`$ under information $`\mu `$ is defined as
$$u_i^\mu (s_i,s_i)=a_{i,s_i}^\mu A^\mu ,A^\mu =\underset{j=1,N}{}a_{j,s_j}^\mu $$
(3)
where $`s_i=\{s_j\}_{ji}`$. It is positive whenever $`i`$ is in the minority group, whence the name of the game. Moreover, players interact with each other only through a global quantity (namely, $`A^\mu `$). This feature clarifies the mean-field character of the model. The total loss experienced by players under information $`\mu `$ simply reads
$$\underset{i=1,N}{}u_i^\mu (s_i,s_i)=(A^\mu )^2,$$
(4)
which is always positive.
### 2.2 Dynamics
A snapshot configuration of the system corresponds to a point $`\{s_i\}_{i=1}^N`$ in the (pure) *strategy space* $`\{\pm 1\}^N`$. The game is repeated and at each time step players face the problem of choosing the strategy to follow. By assumption, each player keeps a “score” $`U_{i,s}(t)`$ for each strategy $`s=\pm 1`$ and updates it as the game proceeds. In the beginning, players set $`U_{i,\pm 1}(0)=0`$. Then for $`t0`$ scores are updated according to the map
$$U_{i,s}(t+1)=U_{i,s}(t)\frac{1}{P}a_{i,s}^{\mu (t)}\left[A^{\mu (t)}\eta \left(a_{i,s_i(t)}^{\mu (t)}a_{i,s}^{\mu (t)}\right)\right]$$
(5)
where $`\eta `$ and $`s_i(t)`$ denotes the strategy that player $`i`$ actually uses at time $`t`$. The term proportional to $`\eta `$ is introduced to model agents who account for their market impact. We refer the reader to Ref. for a detailed discussion of this term<sup>4</sup><sup>4</sup>4In Refs. the term proportional to $`\eta `$ is $`\eta \delta _{s,s_i(t)}/P`$. This leads, however, to the same results of the last term in Eq. (5), in the statistical mechanics approach. The reason is that the approach of Refs. is based on the average of the evolution equation in the stationary state. Observing that $`a_{i,s_i(t)}^{\mu (t)}a_{i,s}^{\mu (t)}`$ is $`1`$ when $`s_i(t)=s`$ and a random sign, with zero average, otherwise, we find that the time average of the last term of Eq. (5) is the same as the average of $`\eta (\delta _{s,s_i(t)}1)/P`$. Hence the two equations are equivalent (apart from an irrelevant constant $`\eta /P`$)..
Let it suffice to say that with $`\eta =0`$ Eq. (5) reduces to the standard minority game dynamics: In this case agents reward (penalize) strategies which would have prescribed a sign opposite (equal) to that of $`A^\mu `$. In doing so, agents ignore the fact that if they actually had played those strategies, their contribution to $`A^\mu `$, and hence $`A^\mu `$ itself, could have changed. With $`\eta =1`$ instead, agents correctly accounting for their contribution to $`A^\mu `$. Hence the reward to strategy $`s`$ is really the payoff that agent $`i`$ would have received had he played that strategy. The parameter $`\eta `$ tunes the extent to which agents account for their market impact.
Following , we assume that the probability with which player $`i`$ chooses the strategy to adopt at time $`t`$ depends on the strategy’s score as follows:
$$\mathrm{Prob}\{s_i(t)=\pm 1\}=C\mathrm{exp}[\mathrm{\Gamma }U_{i,\pm 1}(t)],$$
(6)
where $`C`$ is a normalization constant and $`\mathrm{\Gamma }>0`$ is the learning rate. With this rule, the most successful strategy is more likely to be chosen<sup>5</sup><sup>5</sup>5For $`\eta =0`$, our approach gives correct results only for sufficiently large $`\mathrm{\Gamma }`$, namely for $`\mathrm{\Gamma }>\mathrm{\Gamma }_c(\alpha )`$ of Ref. .
Note that $`A^\mu `$ is the contribution of $`N`$ terms whereas the term proportional to $`\eta `$ in Eq. (5) is of order one. One may naively argue that the $`\eta `$ term is irrelevant, as $`N\mathrm{}`$. This is not so for exactly the same reason for which the Onsager reaction term – or cavity field – is relevant in mean-field spin glass theory <sup>6</sup><sup>6</sup>6While the contribution of other spins to the effective field acting on spin $`i`$ have fluctuating signs, the contribution of spin $`i`$ has always the sign of spin $`i`$. Mean-field theory needs to be corrected with the subtraction of the self-interaction from the effective-field, which becomes a cavity field. Likewise the contribution of agent $`ji`$ to $`A^\mu `$ is uncorrelated with the action of agent $`i`$, whereas the contribution of agent $`i`$ itself is totally correlated with his action.. Indeed Refs. have shown that for $`\eta =1`$ the dynamics converges to a NE. This means that, in a sense, players have become fully sophisticated. For $`\eta =0`$, instead, agents to converge to a sub-optimal state.
We introduce the continuous variables (“soft spins”) $`\varphi _i(t)[1,1]`$ as
$$\varphi _i=s_i$$
(7)
where $`\mathrm{}`$ stands for the average over the distribution of $`s_i`$ in the stationary state. The system is then described by a point $`\{\varphi _i\}_{i=1}^N`$ in the hypercube $`[1,1]^N`$. Hence the $`\varphi _i`$’s are the relevant dynamical variables and the phase space is $`[1,1]^N`$.
The analytic study of the dynamics Eq. (5) has been carried out in Ref. . We shall therefore omit the details and limit ourselves to a brief outline of the results. One can show that the stationary states of the dynamics correspond to the minima of the function
$$H_\eta =\underset{ij}{\overset{1,N}{}}\overline{\xi _i^\mu \xi _j^\mu }\varphi _i\varphi _j+2\underset{i=1}{\overset{N}{}}\overline{\mathrm{\Omega }^\mu \xi _i^\mu }\varphi _i+\eta \underset{i=1}{\overset{N}{}}\overline{(\xi _i^\mu )^2}(1\varphi _i^2)+\overline{(\mathrm{\Omega }^\mu )^2},$$
(8)
where $`\overline{\mathrm{}}=(1/P)_{\mu =1,P}\mathrm{}`$ and $`\mathrm{\Omega }^\mu =_{i=1,N}\omega _i^\mu `$.
Analyzing the stationary states of the dynamics is then equivalent to minimizing $`H_\eta `$. The limiting cases $`\eta =0`$ and $`\eta =1`$ correspond to
$$H_0=\overline{A^\mu ^2}\text{and}H_1=\overline{(A^\mu )^2}$$
(9)
respectively. $`H_0`$ (whose minima describe the standard minority game) is related to the market’s predictability or available information, as explained at length in Refs. . In fact, if $`A^\mu 0`$, then one can predict that the action $`a^\mu =\text{sign}A^\mu `$ is more likely to be successful than the other whenever pattern $`\mu `$ arises. $`H_1`$ is the long time total loss of players averaged over $`\mu `$, as is clear from Eq. (4). In previous works this quantity is usually denoted as $`\sigma ^2`$. We stress the fact that the minima of $`H_1`$ are the game’s NE.
It is easy to understand that :
1. $`H_0`$ is a positive definite quadratic form. Hence for $`\eta =0`$ there is a unique stationary state, corresponding to the cooperative state observed in early numerical simulations;
2. for any $`\eta >0`$ both the global efficiency and the individual payoffs are sensibly improved with respect to the $`\eta =0`$ case;
3. for $`\eta =1`$ there is a large number of stationary states, i.e., of NE. These states have $`\varphi _i^2=1`$, i.e., agents play pure strategies.
Point 1. has been treated extensively in previous works. Here we shall focus on points 2. and 3., namely on the $`\eta >0`$ case.
## 3 Replica approach
### 3.1 Replica-symmetric theory
In order to minimize the function in Eq. (8) we can resort to statistical mechanics methods, for
$$\underset{\{\varphi _i\}_{i=1}^N[1,+1]^N}{\mathrm{min}}H_\eta =\underset{\beta \mathrm{}}{lim}\frac{1}{\beta }\mathrm{log}Z(\beta )\underset{\beta \mathrm{}}{lim}F_\eta (\beta ),$$
(10)
where $`Z(\beta )`$ is the canonical partition function associated to $`H_\eta `$. Further, since $`H_\eta `$ contains quenched disorder we need to apply the replica formalism to analyze its ground states. If we let $`J=\{a_{i,s_i}^\mu \}`$ denote collectively the disorder variables and $`E_J(\mathrm{})`$ denote statistical average over $`J`$, the “typical free energy” $`F_\eta (\beta )`$ can be obtained from the identity
$$E_J[\mathrm{log}Z(\beta )]=\underset{n0}{lim}\frac{E_J[Z^n(\beta )]1}{n}.$$
(11)
A long but standard computation (see appendices in Refs ) leads to the following expression:
$`F_\eta (\beta )`$ $`=`$ $`{\displaystyle \frac{\alpha }{2\beta n}}\mathrm{Tr}\left\{\mathrm{log}\left[\left(1+{\displaystyle \frac{\beta }{\alpha }}\right)𝕀_n+{\displaystyle \frac{\beta }{\alpha }}\widehat{q}\right]\right\}+{\displaystyle \frac{\alpha \beta }{2n}}{\displaystyle \underset{ab}{\overset{1,n}{}}}r_{ab}q_{ab}+`$ (12)
$`{\displaystyle \frac{1}{\beta n}}\mathrm{log}\left\{\mathrm{Tr}_s\left[\mathrm{exp}\left({\displaystyle \frac{\alpha \beta ^2}{2}}{\displaystyle \underset{ab}{\overset{1,n}{}}}r_{ab}q_{ab}\right)\right]\right\}+{\displaystyle \frac{\eta }{2}}(1Q_a).`$
where $`𝕀_n`$ is the $`n`$ dimensional unit matrix, $`\widehat{q}`$ is the *overlap matrix* with elements $`q_{ab}=\varphi _i^a\varphi _i^b`$ ($`a,b=1,\mathrm{},n`$; $`ab`$), and $`Q_a=(1/N)_{i=1,N}(\varphi _i^a)^2`$ ($`a=1,\mathrm{},n`$). The quantities $`r_{ab}`$ and $`R_a`$ appear as Lagrange multipliers associated to $`q_{ab}`$ and $`Q_a`$, respectively.
Imposing the replica-symmetric (RS) Ansatz ($`q_{ab}=q`$ for all $`ab`$, $`Q_a=Q`$ for all $`a`$, and similarly for $`r_{ab}`$ and $`R_a`$) one obtains
$`F_\eta ^{(\mathrm{RS})}(\beta )`$ $`=`$ $`{\displaystyle \frac{\alpha }{2}}{\displaystyle \frac{1+q}{\alpha +\beta (Qq)}}+{\displaystyle \frac{\alpha }{2\beta }}\mathrm{log}\left[1+{\displaystyle \frac{\beta (Qq)}{\alpha }}\right]+{\displaystyle \frac{\beta }{2}}(RQrq)+`$ (13)
$`{\displaystyle \frac{1}{\beta }}\mathrm{log}{\displaystyle _1^1}\mathrm{exp}[\beta V(s;z)]𝑑s_z+{\displaystyle \frac{\eta }{2}}(1Q).`$
with
$$V(s;z)=\sqrt{\frac{\alpha r}{2}}zs\frac{\alpha \beta }{4}(Rr_1)s^2.$$
(14)
The ground state properties of $`H_\eta `$ in the $`N\mathrm{}`$ limit can now be studied by solving the the saddle point equations (obtained by setting equal to zero the derivatives of $`F_\eta ^{(\mathrm{RS})}(\beta )`$ with respect to $`Q`$, $`q`$, $`R`$ and $`r`$) in the $`\beta \mathrm{}`$ limit, since
$$\underset{N\mathrm{}}{lim}\underset{\{\varphi _i\}_{i=1}^N[1,+1]^N}{\mathrm{min}}\frac{H_\eta }{N}=\underset{\beta \mathrm{}}{lim}\frac{F_\eta ^{(\mathrm{RS})}(\beta )|_{\text{s.p.}}}{N},$$
(15)
the subscript s.p. indicating the function $`F_\eta ^{(\mathrm{RS})}(\beta )`$ computed at the saddle point values of $`Q`$, $`q`$, $`R`$ and $`r`$. This procedure yields the so-called RS solution.
For $`\eta =0`$ this solution is characterised by a phase transition at $`\alpha _c0.3374\mathrm{}`$ separating a symmetric phase ($`\alpha <\alpha _c`$) with $`H_0=0`$ from an asymmetric one ($`\alpha >\alpha _c`$) with $`H_0>0`$. The “spin susceptibility” $`\chi =\beta (Qq)/\alpha `$ diverges as $`\alpha \alpha _c^+`$. Also, it is known that this solution is stable against replica symmetry breaking for all values of the control parameter $`\alpha `$. A detailed account of this case can be found in Ref. .
We expect replica symmetry to breakdown for all $`\eta >0`$ at certain critical values of $`\alpha `$, denoted by $`\alpha _{AT}(\eta )`$. In order to test this prediction, we study the stability of the RS solution for generic $`\eta `$.
We shall see that, in the $`\beta \mathrm{}`$ limit a RSB phase arises. It is natural then to ask whether there is a critical temperature $`\beta _c`$ separating a high temperature behavior from a low temperature one. This question, even if not directly related to the Minority Game, may be of interest in its own and can be answered at the RS level. Without going into details, let it suffice to mention that setting $`q=0`$ in the RS saddle point equations – which is correct for the high temperature phase – one finds $`\beta _c=0`$ for all values of $`\alpha `$ and $`\eta `$.
### 3.2 Entropy of the RS solution for $`\eta =1`$
In random Ising spin systems useful indications about the stability of the RS solution are provided by the zero temperature entropy, namely
$$S_\eta ^{(\mathrm{RS})}(\beta \mathrm{})=\underset{\beta \mathrm{}}{lim}\beta ^2\frac{F_\eta ^{(\mathrm{RS})}(\beta )}{\beta ^2}.$$
(16)
This quantity is non-negative due to the discreteness of the configuration space of the model (i.e., $`\{\pm 1\}^N`$). If its zero temperature limit turns out to be negative, the corresponding solution is unstable and further steps in the approximation are needed.
In our case, this point is more subtle than it seems. In fact, $`H_\eta `$ is defined for continuous variables $`\varphi _i[1,+1]`$, and not for Ising (i.e., discrete) variables. The corresponding configuration space is not a discrete set of points, but rather a continuum. Therefore, in principle, the zero temperature entropy need not be non-negative.
For $`\eta =1`$, however, $`H_\eta `$ attains its minima on the corners of the phase space<sup>7</sup><sup>7</sup>7This is a consequence of the fact that $`H_1`$ is an harmonic function of $`\varphi _i`$, i.e. $`_\varphi ^2H_1=0`$. This implies that extrema occurs on the corners. $`[1,1]^N`$, i.e., $`\varphi _i=\pm 1`$ for all $`i`$. Hence the zero temperature entropy calculation is revealing. Leaving details of the computation aside, the final result is that for $`\alpha >1/\pi `$
$$S_1^{(\mathrm{RS})}(\beta \mathrm{})=\frac{\alpha }{2}\left[\frac{C}{\alpha +C}\mathrm{log}\left(1+\frac{C}{\alpha }\right)\right],C=\frac{\alpha }{\sqrt{\pi \alpha }1}.$$
(17)
This is negative for all $`\alpha >1/\pi `$ and it diverges at $`\alpha =1/\pi `$. For $`\alpha <1/\pi `$ one finds $`S_1=\mathrm{}`$. This means that RSB occurs for all values of $`\alpha `$ at $`\eta =1`$, or, in other words, that the study of NE requires RSB.
### 3.3 The number of Nash equilibria
A crucial feature of the occurrence of replica symmetry breaking is gained from the study of the number of minima of $`H_1`$. We show here that an exponentially (in $`N`$) large number of such minima occurs, i.e., that the game possesses an exponentially large number of NE. Our analytic results will be supported by numerical investigations.
In order to compute the number of Nash equilibria (NE) we use the fact that NE are in pure strategies<sup>5</sup>, or, that at any NE $`\varphi _i=\pm 1`$ for all $`i`$. Keeping this in mind, we start by considering that NE satisfy the condition
$$\varphi _i\left[\overline{u_i^\mu (+1,s_i)}\overline{u_i^\mu (1,s_i)}\right]=2\left[\overline{(\xi _i^\mu )^2}\overline{A^\mu \xi _i^\mu }\varphi _i\right]0,i.$$
(18)
Hence an indicator function for NE (using $`2\overline{(\xi _i^\mu )^2}1`$) is
$$I_{NE}(\{\varphi _i\})=\underset{i=1}{\overset{N}{}}\theta \left(12\overline{A^\mu \xi _i^\mu }\varphi _i\right)$$
(19)
($`I_{NE}=1`$ if $`\{\varphi _i\}`$ is a NE and $`=0`$ otherwise) and the number of NE is just obtained summing over all configuration $`\{\varphi _i\}\{\pm 1\}^N`$ (an operation which we denote by $`\text{Tr}_\varphi `$). Hence $`𝒩_{NE}=\text{Tr}_\varphi I_{NE}\{\varphi _i\}`$. Following we take the average over the disorder and introduce the integral representation of the $`\theta `$ function. We arrive at<sup>8</sup><sup>8</sup>8 The reader is warned that the $`\mathrm{\Omega }`$ appearing here is in no relation with the $`\mathrm{\Omega }^\mu `$ introduced in Section 2.
$$E_J\left(𝒩_{NE}\right)=\frac{N^2\alpha ^3}{(2\pi )^2}_{\mathrm{}}^{\mathrm{}}𝑑\gamma 𝑑\mathrm{\Gamma }𝑑\omega 𝑑\mathrm{\Omega }\mathrm{exp}\left[N\mathrm{\Sigma }(\gamma ,\mathrm{\Gamma },\omega ,\mathrm{\Omega })\right]$$
(20)
with
$$\mathrm{\Sigma }(\alpha ,\gamma ,\mathrm{\Gamma },\omega ,\mathrm{\Omega })=\alpha \omega \gamma +\alpha ^2\mathrm{\Omega }\mathrm{\Gamma }\frac{\alpha }{2}\mathrm{log}\left[(1+\gamma )^2+2\mathrm{\Gamma }\right]+\mathrm{log}\left[1+\mathrm{erf}\left(\frac{1\omega }{2\sqrt{\mathrm{\Omega }}}\right)\right]$$
(21)
Eq. (20) is dominated by the saddle point of $`\varphi `$, which is attained at $`\omega =1\gamma `$, $`\mathrm{\Omega }=\frac{1\gamma }{\alpha (1+\gamma )}`$ and $`\mathrm{\Gamma }=\frac{\gamma ^2(1+\gamma )}{2(1\gamma )}`$, where $`\gamma `$ is the root of the equation
$$\frac{\gamma ^2}{4}\frac{\alpha (1+\gamma )}{1\gamma }=\mathrm{log}\left(\gamma \sqrt{\frac{\alpha (1+\gamma )}{1\gamma }}\right)\mathrm{log}\left\{\alpha \gamma ^2\left[1+\mathrm{erf}\left(\frac{\gamma }{2}\sqrt{\frac{\alpha (1+\gamma )}{1\gamma }}\right)\right]\right\}.$$
(22)
In terms of the solution $`\gamma ^{}`$ of this equation we have
$$\mathrm{\Sigma }(\alpha )=\frac{\alpha \gamma ^{}}{2}(2\gamma ^{})\frac{\alpha }{2}\mathrm{log}\left(\frac{1+\gamma ^{}}{1\gamma ^{}}\right)+\mathrm{log}\left[1+\mathrm{erf}\left(\frac{\gamma ^{}}{2}\sqrt{\frac{\alpha (1+\gamma ^{})}{1\gamma ^{}}}\right)\right].$$
(23)
The behavior of $`\mathrm{\Sigma }`$ as a function of $`\alpha `$ is shown in Fig. 1. As expected (see ), as $`\alpha 0`$ the number of NE grows as $`2^N`$. Numerical results from exact enumeration for $`N20`$ are in very good agreement, which shows that the so-called annealed approximation used here (i.e., taking the average of $`𝒩_{NE}`$) is sufficient and one does not need to introduce replicas (to compute the average of $`\mathrm{log}𝒩_{NE}`$) at this level.
### 3.4 de Almeida-Thouless line
A more thorough analysis can be obtained using the de Almeida-Thouless (AT) protocol . In order to investigate the stability of the RS ground states of $`H_\eta `$ against RSB we compute the matrix of the second derivatives of the general expression for the free energy, Eq. (12), with respect to $`q_{ab}`$ and $`r_{ab}`$. The conditions for RSB are then obtained by studying the effect of fluctuations in the direction of RSB. This analysis results in an instability line – i.e. a family of points where the RS solution becomes unstable – in the parameter space $`(\alpha ,\eta )`$, called the AT line and denoted by $`\alpha _{AT}(\eta )`$. The resulting equation has to be solved (numerically) together with the RS saddle point equations.
An outline of the calculation is reported in the appendix. We have studied particularly the so called replicon mode, namely those eigenvectors of the stability matrix which are symmetric under interchange of all but two of the indices. The replicon mode is typically responsible for the onset of the RSB instability. Points on the AT line are found to satisfy the following stability condition:
$$\alpha \left[1\eta \left(1+\frac{\beta (Qq)}{\alpha }\right)\right]^2=1.$$
(24)
The “susceptibility” $`\chi \beta (Qq)/\alpha `$ remains finite as $`\beta \mathrm{}`$, so that the zero temperature behaviour can be safely detected. Results are reported in Fig. 2. For $`\eta =0`$ replica symmetry is preserved for all $`\alpha `$. The point $`\alpha _c=0.3374\mathrm{}`$ where the second order phase transition occurs in the standard MG, separates a line of first order phase transitions, for $`\alpha <\alpha _c`$, from a second order line. For $`\eta =0^+`$ one finds RSB for $`\alpha _{AT}=1^+`$. Finally, for $`\eta =1`$, RS is broken for all $`\alpha `$.
### 3.5 Replica symmetry breaking
The one step breaking of replica permutation symmetry is expressed by the Parisi Ansatz for the $`q_{ab}`$’s and the $`Q_a`$’s, where an additional parameter denoted by $`m`$ is introduced: $`Q_a=Q`$ (all $`a`$), $`q_{ab}=q_1`$ (all $`ab`$ such that $`|ab|m`$) and $`q_{ab}=q_0`$ (otherwise). The “free energy” is the same as in Eq. (12), but this time the overlap matrix has to be parameterised as
$$\widehat{q}=q_0ϵ_nϵ_n^T+(q_1q_0)𝕀_{\frac{n}{m}}ϵ_mϵ_m^T+(Qq_1)𝕀_n,$$
(25)
where $`ϵ_n`$ is the $`n`$-dimensional column vector with all components equal to one (so that $`ϵ_nϵ_n^T`$ is the $`n`$ dimensional matrix with all elements equal to one) and we have used the standard tensor product.
We need to consider the matrix
$$\widehat{T}=\left(1+\frac{\beta }{\alpha }\right)𝕀_n+\frac{\beta }{\alpha }\widehat{q},$$
(26)
that is
$$\widehat{T}=\left[1+\frac{\beta }{\alpha }(Qq_1)\right]𝕀_n+\frac{\beta }{\alpha }(1+q_0)ϵ_nϵ_n^T+\frac{\beta }{\alpha }(q_1q_0)𝕀_{\frac{n}{m}}ϵ_mϵ_m^T.$$
(27)
Using the identities
$`𝕀_n`$ $``$ $`𝕀_{\frac{n}{m}}𝕀_m`$
$`ϵ_nϵ_n^T`$ $``$ $`ϵ_{\frac{n}{m}}ϵ_{\frac{n}{m}}^Tϵ_mϵ_m^T`$ (28)
we can decompose $`\widehat{T}`$ into tensor products and write its determinant $`|\widehat{T}|`$ in a straightforward way (we need $`|\widehat{T}|`$ because $`\mathrm{Tr}(\mathrm{log}\widehat{T})=\mathrm{log}|\widehat{T}|`$). We get
$`|\widehat{T}|`$ $`=`$ $`[1+{\displaystyle \frac{\beta }{\alpha }}(Qq_1)+n{\displaystyle \frac{\beta }{\alpha }}(1+q_0)+m{\displaystyle \frac{\beta }{\alpha }}(q_1q_0)]\times `$
$`\times \left[1+{\displaystyle \frac{\beta }{\alpha }}(Qq_1)+m{\displaystyle \frac{\beta }{\alpha }}(q_1q_0)\right]^{\frac{n}{m}1}\left[1+{\displaystyle \frac{\beta }{\alpha }}(Qq_1)\right]^{n\frac{n}{m}}.`$
This means that the first factor on the r.h.s. is an eigenvalue of $`\widehat{T}`$ with multiplicity one, the second one is an eigenvalue with multiplicity $`\frac{n}{m}1`$, and the third one is an eigenvalue with multiplicity $`n\frac{n}{m}`$.
Putting the latter formula into Eq. (12) and taking the $`n0`$ limit as requested by the replica trick, we obtain the one step replica symmetry broken free energy $`F_\eta ^{(1\mathrm{R}\mathrm{S}\mathrm{B})}(\beta )`$, whose final expression is
$`F_\eta ^{(1\mathrm{R}\mathrm{S}\mathrm{B})}(\beta )`$ $`=`$ $`{\displaystyle \frac{\alpha }{2}}{\displaystyle \frac{1+q_0}{\alpha +\beta (Qq_1)+m\beta (q_1q_0)}}+{\displaystyle \frac{\eta }{2}}(1Q)+`$
$`+{\displaystyle \frac{\alpha }{2\beta }}\mathrm{log}\left[1+{\displaystyle \frac{m\beta (q_1q_0)}{\alpha +\beta (Qq_1)}}\right]+`$
$`+{\displaystyle \frac{\alpha }{2\beta m}}\mathrm{log}\left[1+{\displaystyle \frac{m\beta (q_1q_0)}{\alpha +\beta (Qq_1)}}\right]+`$
$`+{\displaystyle \frac{\alpha \beta }{4}}[RQ+(m1)r_1q_1mr_0q_0]+`$
$`{\displaystyle \frac{1}{m\beta }}\mathrm{log}\left({\displaystyle _1^1}ds\mathrm{exp}[\beta V(s;y,z)]\right)^m_y_z,`$
where $`\mathrm{}_x`$ denotes again the average over the unit gaussian variables $`x`$ and
$$V(s;y,z)=\sqrt{\frac{\alpha r_0}{2}}zs\sqrt{\frac{\alpha (r_1r_0)}{2}}ys\frac{\alpha \beta }{4}(Rr_1)s^2.$$
(30)
$`F_\eta ^{(1\mathrm{R}\mathrm{S}\mathrm{B})}(\beta )`$ depends on seven parameters: the three overlap matrix elements $`Q`$, $`q_0`$, $`q_1`$, their related Lagrange multipliers $`R`$, $`r_0`$ and $`r_1`$, and $`m`$. Their values have to be determined self-consistently from the seven saddle point equations obtained by setting to zero the derivatives of the free energy with respect to the above parameters. These equations can be solved numerically.
One finds three different regimes in the $`(\alpha ,\eta )`$ plane when $`\beta \mathrm{}`$ (Figure 2):
1. For $`\alpha <\alpha _00.09012\mathrm{}`$ (all $`\eta >0`$) one has $`H_\eta =0`$. The solution does not depend on $`\eta `$ as long as $`\eta >0`$. The self-overlap is $`Q=1`$ signalling that agents play pure strategies ($`\varphi _i=\pm 1`$) but off diagonal overlaps $`q_1>q_0`$ are both less than $`1`$. This suggests that NE are organized in a complex geometric structure. The parameter $`m`$ attains a finite value.
2. For $`\alpha _0<\alpha <\alpha _1(\eta )`$ (all $`\eta >0`$) the solution has $`H_\eta >0`$, it is independent of $`\eta `$ (for $`\eta >0`$) and $`1=Q=q_1>q_0`$. The spin susceptibility $`\chi =\beta (Qq_1)/\alpha `$ attains a finite value in the limit $`\beta \mathrm{}`$, which diverges as $`\alpha \alpha _0`$. Again agents play pure strategies and $`q_0<1`$ is the typical overlap between two NE. The parameter $`m`$ vanishes as $`1/\beta `$ (indeed the $`\beta m`$ is finite as $`\beta \mathrm{}`$). The line $`\alpha _1(\eta )`$ is determined by the solution of
$$\frac{\eta }{2}=\frac{1}{\alpha +\beta (Qq_1)}.$$
(31)
3. In between the line $`\alpha _1(\eta )`$ and the stability line $`\alpha _{\text{AT}}(\eta )`$ the solution has $`H_\eta >0`$ and $`1>Q=q_1>q_0`$. Hence agents do not play pure strategies. The solution in this region depends on $`\eta `$.
We stress again that NE are in pure strategies, since at $`\eta =1`$ for all values of $`\alpha `$ one finds $`Q=1`$.
Figure 3 shows that the one-step calculation for $`H_1/N`$ agrees very well with numerical simulations and it represents a considerable improvement over the replica symmetric result<sup>9</sup><sup>9</sup>9Note that numerical results refer to a typical NE which need not be the ground state of $`H_1`$. Further steps of RSB, most probably infinitely many, are likely to be needed to recover exact results. However, already the one step calculation provides a rather good approximation.
## 4 Conclusion
Summarizing, we have analyzed the solution of the minority game by means of statistical mechanics methods. Our starting point has been the study of Refs , where the NE of the game have been mapped onto the ground states of a disordered hamiltonian with RSB, suggesting the existence of a very large number of them. First, we have computed (both analitically and numerically) the number of NE, showing that actually they are exponentially many in $`N`$ (number of players). Then, we probed the stability of the replica symmetric theory developed in Ref. . After showing the necessity of RSB by simple entropy considerations, we have calculated the instability line (AT line) using the de Almeida-Thouless method. Finally, we have derived the broken-replica-symmetry solution, drawing the complete phase diagram of the model. All our results are in excellent agreement with computer experiments.
To our knowledge, the minority game is the first example of a market game that requires the full use of spin glass theory in order to uncover its behaviour. Remarkably, many features actually observed in real markets can be recovered within the simple setup of the minority game . Understanding real markets is among the most challenging theoretical problems ahead of us. This work suggests that statistical mechanics of disordered systems may be a valuable tool in this endeavour.
Acknowledgments We acknowledge frequent discussions with D. Challet, S. Franz, L. Giada, F. Ricci-Tersenghi, R. Zecchina and Y.C. Zhang which helped us considerably in the developement of this work.
## Appendix A Calculation of the AT line
The stability matrix has dimension $`n(n1)\times n(n1)`$ and is given by
$$C=\left(\begin{array}{cc}A^{(ab,cd)}& D^{(ab,cd)}\\ D^{(ab,cd)}& B^{(ab,cd)}\end{array}\right)$$
(32)
where
$$A^{(ab,cd)}=\frac{^2(nF)}{q_{ab}q_{cd}},B^{(ab,cd)}=\frac{^2(nF)}{r_{ab}r_{cd}}\mathrm{and}D^{(ab,cd)}=\frac{^2(nF)}{q_{ab}r_{cd}}$$
(33)
($`F`$ denotes shortly the replica-symmetric free energy.)
Introducing the “perturbation” of the RS solution in the form
$$\delta q_{ab}=\zeta _{ab}\mathrm{and}\delta r_{ab}=x\zeta _{ab}$$
(34)
with the condition $`_b\zeta _{ab}=0`$ for all $`a`$, it is possible to show that this condition is satisfied for all $`n`$ by
$`\zeta _{ab}=\zeta `$ $`(a,b)(1,2)`$
$`\zeta _{1b}=\zeta _{2b}={\displaystyle \frac{1}{2}}(3n)\zeta `$ $`b1,2`$
$`\zeta _{12}={\displaystyle \frac{1}{2}}(2n)(3n)\zeta `$ $`(a,b)=(1,2)`$ (35)
$`\zeta _{aa}=0`$ $`a.`$
The relevant eigenvalue equations (the so-called *replicon mode*) are given by
$`{\displaystyle \underset{cd}{}}\left(A^{(ab,cd)}+xD^{(ab,cd)}\right)\zeta _{cd}`$ $`=`$ $`\lambda \zeta _{ab}`$
$`{\displaystyle \underset{cd}{}}\left(xA^{(ab,cd)}+D^{(ab,cd)}\right)\zeta _{cd}`$ $`=`$ $`\lambda x\zeta _{ab}`$ (36)
One needs to find an expression for the matrix elemnts $`A`$, $`B`$ and $`D`$.
There are three different types of matrix elements in the replica symmetric state, corresponding to the cases $`(a,b)=(c,d)`$, $`a=c`$ and $`(a,b)(c,d)`$ respectively. For the $`A^{(ab,cd)}`$ they are $`A^{(ab,ab)}`$, $`A^{(ab,ad)}`$ and $`A^{(ab,cd)}`$. It is simple to show that in the replica symmetric state
$`A^{(ab,ab)}`$ $`=`$ $`\alpha \beta (E_{ab^2}+E_{aa}^2)`$
$`A^{(ab,ad)}`$ $`=`$ $`\alpha \beta [E_{ab}(E_{ab}+E_{aa})]`$
$`A^{(ab,cd)}`$ $`=`$ $`2\alpha \beta E_{ab}^2`$ (37)
where
$$E_{ab}=\beta q[\alpha +\beta (Qq)]^2\mathrm{and}E_{aa}=E_{ab}+[\alpha +\beta (Qq)]^1.$$
(38)
For the $`B^{(ab,cd)}`$ we find
$`B^{(ab,ab)}`$ $`=`$ $`\alpha ^2\beta ^3\left(s^2^2_zs^2_z^2\right)`$
$`B^{(ab,ad)}`$ $`=`$ $`\alpha ^2\beta ^3\left(s^2s^2_zs^2_z^2\right)`$ (39)
$`B^{(ab,cd)}`$ $`=`$ $`\alpha ^2\beta ^3\left(s^4_zs^2_z^2\right)`$
As for $`D^{(ab,cd)}`$, one finds the general result
$$D^{(ab,cd)}=\alpha \beta \left(\delta _{ac}\delta _{bd}+\delta _{ad}\delta _{bc}\right).$$
(40)
It is important to notice that the relevant combinations of matrix elements appearing in the eigenvalue equations are of the form $`A^{(ab,cd)}2A^{(ab,ad)}+A^{(ab,cd)}`$, and that the eigenvalues can be shown to depend only on
$$a=A^{(ab,cd)}2A^{(ab,ad)}+A^{(ab,cd)}\mathrm{and}b=B^{(ab,cd)}2B^{(ab,ad)}+B^{(ab,cd)}$$
(41)
via the simple formula
$$\lambda _\pm =\frac{1}{2}\left(a+b\pm \sqrt{(ab)^2+4}\right).$$
(42)
Putting things together and solving the eigenvalue equations, one finds that one of the eigenvalue (namely $`\lambda _{}`$) is constant in sign (at least at low temperatures). The second one, instead, changes sign and signals the onset of RSB instability. The equation corresponding to $`\lambda _+=0`$ in the end reads
$$\frac{\alpha \beta ^2}{\alpha ^2[1+\beta (Qq)/\alpha ]^2}(s^2s^2)^2_z=1.$$
(43)
Calculating explicitly the averages appearing in the above formula one arrives at the AT line reported in the text:
$$\alpha [1\eta (1+\beta (Qq)/\alpha )]^2=1.$$
(44)
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# Visualization of correlation cascade in spatio-temporal chaos using wavelets
## I Introduction
It has long been a subject of discussion how to capture and describe complex dynamical processes of strongly correlated many-body systems that are ubiquitous in nature. Probably one of the most important problems of such kind at present is the one faced in neuroscience. In what way biological information processing is encoded into spikes of neurons in the brain is the greatest mystery of present neuroscience. In order to capture subtle correlation between activities of neurons, information-theoretical techniques have been developed and applied to experimental signals \[Rieke et al., 1998\].
On the other hand, we may mention Richardson’s cascade picture of fully-developed fluid turbulence as a successful example of capturing the essence of a strongly correlated system, which stated that large eddies gradually fragment into smaller eddies, and eventually dissipate due to viscosity. Apart from the difficult problem of intermittency, this picture was mathematically formulated into the Kolmogorov 1941 theory, which gave the famous $`5/3`$ energy spectrum \[Frisch, 1995\]. This success was due to the fact that in fluid turbulence, fortunately, a simple hierarchy in spatial scale (or equivalently, in temporal / energy scales) is formed as a result of complex nonlinear dynamics of a number of modes. Since the success in fluid turbulence, attempts have been made to capture the dynamics of other kinds of many-body systems such as spatio-temporal chaos \[Ikeda & Matsumoto, 1989\], chaotic systems with large degrees of freedom \[Kaneko, 1994\], and stock markets \[Arnéodo et al., 1998a\], based on the picture that energy, or more generally information or causality, flows from scales to scales.
In analyzing space-scale properties of various kinds of signals, the wavelet transform is expected to be a useful tool due to its localized nature both in space and scale, and many applications have been made \[Mallat, 1999\]. Of course, there are many studies that tried to capture the cascade process of fully-developed turbulence directly using wavelets \[Argoul et al., 1989; Yamada & Ohkitani, 1991; Arnéodo et al., 1998b\]. But most of the studies treated one-dimensional time sequences of the velocity field measured at a single point in fluid, and rather a limited number of studies actually investigated the temporal evolution of correlation between scales of spatially extended fields \[Toh, 1995\].
In this paper, we propose a simple visualization method of such spatio-temporal correlation between scales of spatially extended fields using wavelets, and apply it to two typical spatio-temporally chaotic systems. As such chaotic systems, we investigate ordinary diffusively coupled complex Ginzburg-Landau oscillators \[Kuramoto, 1984; Bohr et al., 1998\] and somewhat unusual non-locally coupled complex Ginzburg-Landau oscillators \[Kuramoto, 1995\]. While the former system exhibits spatially smooth amplitude patterns, the latter system exhibits fractal spatial patterns. Therefore, the dynamical processes behind those systems are expected to differ from each other significantly, and it is interesting to see whether their difference can be detected by our method or not.
## II Coupled complex Ginzburg-Landau oscillators
The ordinary diffusively (i.e., locally) coupled complex Ginzburg-Landau equation \[Kuramoto, 1984; Bohr et al., 1998\] is given by
$$\dot{W}(x,t)=W(1+ic_2)|W|^2W+D(1+ic_1)^2W,$$
(1)
which can be derived, for example, from equations of oscillatory media in the vicinity of their Hopf bifurcation points by the center-manifold reduction technique. Here, $`W(x,t)`$ is a complex amplitude of an oscillator at position $`x`$ and at time $`t`$, $`c_1`$ and $`c_2`$ are real parameters, and $`D`$ is a diffusion constant. It is well known that this equation exhibits spatio-temporal chaos in some appropriate parameter region. We call this equation “LCGL equation” hereafter.
In this paper, we treat only spatially one-dimensional cases. We fix the length of the system at $`L=1`$, and assume a periodic boundary condition. The diffusion constant is set at $`D=0.0035`$, and the parameters are fixed at $`c_1=2`$ and $`c_2=2`$. With these values, the spatially uniform solution of the LCGL equation is unstable, and the system exhibits spatio-temporal chaos. Typical time scale of the system is roughly estimated as the time needed for an uncoupled free oscillator to go around its limit cycle, and is given by $`2\pi /c_23`$. The numerical simulation was done in wavenumber space using $`N=2^{10}2^{14}`$ modes by the pseudo-spectral method with a time step of $`0.01`$ (Euler integration).
Due to the existence of a diffusion term, the solution of the LCGL equation necessarily possesses a characteristic minimal length scale, below which fluctuations are strongly depressed. Thus, the amplitude pattern $`|W(x,t)|`$ of the solution has a smoothly modulated shape as displayed in Fig. 1. The shortest wavelength determined by the diffusion constant is about $`\sqrt{D}0.061/16`$. For the sake of comparison, this value is equated with the coupling length of the non-locally coupled system that we explain later. Figure 2 displays temporal evolution of the amplitude pattern, and Fig. 3 shows its power spectrum. Since short-wavelength components are strongly depressed, the power spectrum decays exponentially.
As another spatio-temporally chaotic system, we investigate a different type of coupled complex Ginzburg-Landau oscillators, which was first introduced by Kuramoto . Instead of a diffusive interaction, it has a non-local interaction; each oscillator feels a non-local mean field of other oscillators through a kernel $`g(|x|)`$, which is a decreasing function of $`|x|`$. Hereafter we use $`g(|x|)=g_0\mathrm{exp}(|x|/\gamma )`$ as the kernel, where $`\gamma `$ gives the coupling range. $`g_0`$ is some appropriate normalization constant, and assumed to be $`2\gamma `$. The equation for this system of non-locally coupled oscillators is given by
$`\dot{W}`$ $`(x,t)=W(1+ic_2)|W|^2W`$ (4)
$`+K(1+ic_1){\displaystyle 𝑑x^{}g(|x^{}x|)\left[W(x^{},t)W(x,t)\right]},`$
which we call “NCGL equation” hereafter. This equation can be derived, for example, from a model of oscillatory biological cells that are interacting through some diffusive chemical substance with a finite decay rate.
Although the difference between this NCGL equation and the LCGL equation is only the last interaction term, it was shown that this NCGL equation exhibits spatio-temporal chaos with remarkably different features from the LCGL equation, such as fractal amplitude patterns, power-law spatial correlation, and power-law spectrum. We again fix the length of the system at $`L=1`$, and assume a periodic boundary condition. The coupling range is fixed at $`\gamma =1/16`$, which is equal to the shortest characteristic length $`\sqrt{D}1/16`$ of the LCGL equation. (More precisely, by assuming the smoothness of the amplitude field $`W(x^{},t)`$ in the interaction term and expanding it up to the second order in $`x^{}x`$ around $`x`$, we obtain an effective diffusion constant $`D_{eff}=K\gamma ^2`$ for the NCGL equation. However, the smoothness of the amplitude field is not always guaranteed, hence this correspondence is only formal.) The parameter values are again set at $`c_1=2`$ and $`c_2=2`$. With these values, the uniform solution of the system is always unstable, and the system behaves in a chaotic manner. Further, the behavior of the system strongly depends on the remaining coupling strength $`K`$; the amplitude pattern of the system is smooth for large $`K`$ ($`1.5`$), while almost discontinuous for small $`K`$ ($`0.5`$). Between these values, there exists a parameter region where the amplitude pattern is fractal, and its spatial correlation exhibits power-law behavior. We fix $`K=0.9`$ hereafter, where the system is exactly in this “anomalous” spatio-temporally chaotic regime. The numerical simulation was done in real space using $`N=2^{10}2^{14}`$ oscillators by the 4th-order Runge-Kutta method with a time step of $`0.05`$.
Figure 4 displays a typical snapshot of the solution of the NCGL equation, and Fig. 5 shows its temporal evolution. In contrast to the case of the LCGL equation, the amplitude pattern is not completely smooth, but is composed of patches of smooth coherent regions and strongly disordered regions. The power spectrum displayed in Fig. 6 also indicates the difference clearly. It exhibits power-law decay rather than the exponential decay in the case of the LCGL equation, which implies the strong anomaly of the amplitude field. Actually, it was shown that the amplitude pattern of the NCGL equation is fractal, and its fractal dimension varies with the coupling strength $`K`$. These properties are essentially the results of the absence of length scale shorter than the coupling range $`\gamma `$, and can be explained using a simple multiplicative stochastic model to a certain extent \[Kuramoto & Nakao, 1996\].
Thus, the difference in the interaction term greatly changes the amplitude pattern of the system, even though the oscillators and characteristic length scales are the same. Especially, the power-law behavior of the power spectrum reminds us of the $`5/3`$ energy spectrum of fluid turbulence. (However, the exponent changes with the coupling strength $`K`$ for our NCGL equation.) The power-law behavior of the energy spectrum in fluid turbulence is related to the cascade process of breakdown of vortices. Thus, we may naively expect that our NCGL equation also possesses some kind of cascade process. Of course, the NCGL equation is strongly dissipative and there is no such conservative quantity like energy or enstrophy. However, we still expect a cascade process of some quantity from long-wavelength components to shorter-wavelength components, which may be called causality or information, or merely correlation. In the following part of this paper, we develop a simple method to visualize the spatio-temporal correlation between fluctuations at different scales, in order to prove the existence of such cascade process.
## III Wavelet-based correlation analysis
Let us consider decomposing the spatial pattern into various scales using wavelets, and analyzing the temporal correlation between them. We decompose a spatial pattern $`V(x)`$ using an orthonormal wavelet basis as
$$V(x)=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}a_{j,k}\psi _{j,k}(x),$$
(5)
where $`a_{j,k}`$ is an expansion coefficient, and $`\psi _{j,k}(x)`$ is a child wavelet generated from the mother wavelet $`\psi (x)`$ by translation and dilation as
$$\psi _{j,k}(x)=\sqrt{2^j}\psi (2^jxk),$$
(6)
where $`j`$ is a scale parameter, and $`k`$ is a translation parameter. The child wavelets are mutually orthogonal as
$$\psi _{j,k},\psi _{j^{},k^{}}=_{\mathrm{}}^{\mathrm{}}\psi _{j,k}(x^{})\psi _{j^{},k^{}}^{}(x^{})𝑑x^{}=\delta _{j,j^{}}\delta _{k,k^{}},$$
(7)
and form a complete orthonormal set. Since the child wavelet $`\psi _{j,k}(x)`$ is localized both in space and scale roughly at position $`k/2^j`$ and scale $`1/2^j`$, the expansion coefficient given by
$$a_{j,k}=V,\psi _{j,k}=_{\mathrm{}}^{\mathrm{}}V(x^{})\psi _{j,k}^{}(x^{})𝑑x^{}$$
(8)
quantifies the magnitude of fluctuation of $`V(x)`$ around this space-scale point. Hereafter we use Meyer’s wavelet. It is a real analytic function that decays faster than any power function as $`|x|\mathrm{}`$, and its moment of any order vanishes, i.e.,
$$_{\mathrm{}}^{\mathrm{}}x^n\psi (x)𝑑x=0,(n0).$$
(9)
Further, it has a smooth Fourier transform with a compact support. These features are very preferable from the viewpoint of physicists. For the details of Meyer’s wavelet and efficient numerical algorism, see Refs. \[Yamada & Ohkitani, 1991; Mallat, 1999\]. Meyer’s mother wavelet and one of its child wavelet are shown in Fig. 7.
In the practical numerical calculation, the amplitude pattern is a periodic function on $`[0,1)`$ represented by $`N=2^n`$ discrete points. Therefore, we extend the pattern over the whole real axis by a periodic extrapolation in order to apply the wavelet transform. Consequently, in the actual numerical calculation, the first subscript $`j`$ in the summation of Eq. (5) indicating the scale runs from $`0`$ to $`n1`$, and the second subscript $`k`$ indicating the translation runs from $`0`$ to $`2^j1`$.
Rather than using the obtained expansion coefficient directly, we define a new coarse-grained field $`b_j(x)(0x<1)`$ by a piecewise constant interpolation from the expansion coefficient $`a_{j,k}`$ as
$$b_j(x)=\mathrm{log}|a_{j,k}|^2(\frac{k}{2^j}x<\frac{k+1}{2^j},k=0,1,\mathrm{},2^j1).$$
(10)
Here we take a square of the coefficient, because we are interested in the absolute intensity of the fluctuation at a certain position and scale. Further taking the logarithm is merely for convenience sake here; we obtain similar results without taking the logarithm. However, there also exists some discussion claiming that taking the logarithm is more appropriate when the system under consideration possesses a random multiplicative process \[Arnéodo et al., 1998a; Arnéodo et al., 1998b\], and the NCGL equation is actually considered this case. Figures 8 and 9 display the coarse-grained fields $`b_j(x)`$ for the LCGL and NCGL equations. The obtained data for the LCGL equation are very noisy for large $`j`$ values, since components whose wavelengths are shorter than the dissipation scale $`\sqrt{D}1/16`$ decay very quickly. On the other hand, we can observe clear positional correlation between scales down to a very small wavelength in the case of the NCGL equation.
Let $`V_1(x)`$ be the amplitude pattern at time $`t`$, i.e., $`V_1(x)=|W(x,t)|`$, and $`V_2(x)`$ be the amplitude pattern evolved from $`V_1(x)`$ for a period of $`\mathrm{\Delta }t`$, i.e., $`V_2(x)=|W(x,t+\mathrm{\Delta }t)|`$. From these amplitude patterns, we obtain coarse-grained fields $`b_{j_1}^1(x)`$ and $`b_{j_2}^2(x)`$ at scales $`j_1`$ and $`j_2`$. We then define deviations of these fields from their mean values as
$$\mathrm{\Delta }b_{j_1}^1(x)=b_{j_1}^1(x)\overline{(b_{j_1}^1,I)},$$
(11)
and
$$\mathrm{\Delta }b_{j_2}^2(x)=b_{j_2}^2(x)\overline{(b_{j_2}^2,I)},$$
(12)
where $`I`$ is a constant-valued function $`I(x)1`$, and the inner product of $`f(x)`$ and $`g(x)`$ is defined as
$$(f,g)=_0^1f(x)g(x)𝑑x.$$
(13)
Now we define a kind of cross-correlation matrix $`C_{j_1,j_2}(\mathrm{\Delta }t)`$ from these $`\mathrm{\Delta }b_{j_1}^1(x)`$ and $`\mathrm{\Delta }b_{j_2}^2(x)`$ as
$$C_{j_1,j_2}(\mathrm{\Delta }t)=\frac{\overline{(\mathrm{\Delta }b_{j_1}^1,\mathrm{\Delta }b_{j_2}^2)}}{\left[\overline{(\mathrm{\Delta }b_{j_1}^1,\mathrm{\Delta }b_{j_1}^1)}\overline{(\mathrm{\Delta }b_{j_2}^2,\mathrm{\Delta }b_{j_2}^2)}\right]^{1/2}}.$$
(14)
From this definition, $`C_{j_1,j_2}(\mathrm{\Delta }t)=C_{j_2,j_1}(\mathrm{\Delta }t)`$ obviously follows. The overlines in the above equations indicate temporal average, assuming the stationarity of the spatio-temporal chaos under consideration.
When $`\mathrm{\Delta }t=0`$, $`V_2(x)`$ represents the same amplitude pattern as $`V_1(x)`$. Thus $`\mathrm{\Delta }b_{j_1}^1(x)`$ and $`\mathrm{\Delta }b_{j_2}^2(x)`$ are identical, and $`C_{j_1,j_2}`$ is merely a symmetric matrix. When $`\mathrm{\Delta }t0`$, $`V_2(x)`$ differs from $`V_1(x)`$ to some degree depending on $`\mathrm{\Delta }t`$. Therefore $`C_{j_1,j_2}`$ is no longer symmetric, and its asymmetry is expected to characterize the variation of spatio-temporal correlation between scales. In order to see this asymmetry clearly, let us decompose the matrix $`C_{j_1,j_2}`$ into symmetric and antisymmetric parts as
$$D_{j_1,j_2}(\mathrm{\Delta }t)=\frac{1}{2}\left\{C_{j_1,j_2}(\mathrm{\Delta }t)+C_{j_2,j_1}(\mathrm{\Delta }t)\right\},$$
(15)
and
$$E_{j_1,j_2}(\mathrm{\Delta }t)=C_{j_1,j_2}(\mathrm{\Delta }t)C_{j_2,j_1}(\mathrm{\Delta }t).$$
(16)
The symmetric part $`D_{j_1,j_2}`$ is invariant under time-reversal transform $`\mathrm{\Delta }t\mathrm{\Delta }t`$, while the antisymmetric part $`E_{j_1,j_2}`$ is not invariant (antisymmetric). Thus the matrix $`E_{j_1,j_2}`$ is expected to quantify the variation of spatio-temporal correlation between scales that is not symmetric to time-reversal, namely, to give a certain measure of the flow of information or correlation between scales.
In Fig. 10, the symmetric part $`D_{j_1,j_2}`$ and the antisymmetric part $`E_{j_1,j_2}`$ of the correlation matrix obtained for the LCGL equation are displayed using color codes for several values of the time difference $`\mathrm{\Delta }t`$. The top row displays results obtained for the same time ($`\mathrm{\Delta }t=0`$), and the lower rows display results for larger values of the time difference $`\mathrm{\Delta }t`$. The vertical axis of each graph represents the scale $`j_1`$ of the pattern at the reference time, and the horizontal axis represents the scale $`j_2`$ of the pattern $`\mathrm{\Delta }t`$ after the reference time. First, for $`\mathrm{\Delta }t=0`$, the antisymmetric part is merely zero and only the symmetric part takes positive values. The diagonal components naturally take largest values, since the self-correlation of the fluctuation at the same scale is strongest. Besides, it can be seen that there exists rather strong correlation between adjacent scales in the dissipation region determined by the diffusion constant $`D`$ (the region with the scale slightly smaller than $`j4`$). When the time difference $`\mathrm{\Delta }t`$ becomes a little larger, the symmetric part diminishes. But still the diagonal components and the dissipation region maintain higher correlations than the others. Now, if we look at the antisymmetric part, there appear small localized regions in the neighborhood of the dissipation region across the diagonal, which have negative and positive correlations, respectively. This implies that fluctuation at a certain scale at the reference time (vertical axis) possesses relatively strong correlation to fluctuation at slightly smaller scale $`\mathrm{\Delta }t`$ after the reference time (horizontal axis). In this case, it can be considered as visualizing the dissipation due to the diffusion term. As the time difference $`\mathrm{\Delta }t`$ becomes further large, both the symmetric and antisymmetric parts tend to take smaller values, and correlation between the patterns vanishes. The intensity of the antisymmetric part is typically $`20\%`$ of the symmetric part at its maximum.
Figure 11 displays the symmetric part $`D_{j_1,j_2}`$ and the antisymmetric part $`E_{j_1,j_2}`$ of the correlation matrix obtained for the NCGL equation for several values of the time difference $`\mathrm{\Delta }t`$ using color codes, as in the case of Fig. 10. When there is no time difference ($`\mathrm{\Delta }t=0`$), only the symmetric part takes finite values, and its diagonal components take largest values as in the case of the LCGL equation. However, the region where the correlation between adjacent scales takes its peak value is located at much smaller scale. When the time difference $`\mathrm{\Delta }t`$ becomes a little larger, the difference becomes more remarkable; there appear wide coherent regions with positive and negative correlation, which spread over scales shorter than the characteristic length of the system ($`j4`$) determined by the coupling range $`\gamma `$. This asymmetry is maintained up to a considerably large value of $`\mathrm{\Delta }t`$, which indicates that long-wavelength components at the reference time maintain strong correlations widely to shorter-wavelength components at later times. As we increase the time difference $`\mathrm{\Delta }t`$ further, the correlation between scales gradually vanishes from the long-wavelength region. But it takes much longer than the case of the LCGL equation. These facts suggest the existence of cascade-like propagation of some quantity from long wavelength to shorter wavelength. The maximum intensity of the antisymmetric part is about $`10\%`$ of the symmetric part.
In order to see the temporal evolution of the antisymmetric part of the correlation matrix in more detail, antisymmetric components $`E_{j,j+1}`$ between adjacent scales $`j`$ and $`j+1`$ normalized by the maximum value at each scale are displayed in Figs. 12 and 13 for scales smaller than the characteristic length. In the case of the LCGL equation, there is no clear ordering of the temporal variation of correlation. Each curve has its maximum at some small value of $`\mathrm{\Delta }t`$, and vanishes quickly. In the case of the NCGL equation, on the other hand, there exists a clear time ordering of the curves in scale; the correlation between long-wavelength fluctuation (small $`j`$ values) becomes large earlier and decreases quickly, while the correlation between short-wavelength fluctuation (large $`j`$ values) becomes large later and decreases very slowly. (Note the difference of time scale between the graphs for the LCGL and NCGL equations.)
## IV Discussion
As we demonstrated, our simple method based on the correlation matrix seems to visualize the spatio-temporal correlation between scales of spatio-temporal chaos successfully. Our method clearly visualized the difference of dynamical process between two systems with different natures. Especially, for the non-locally coupled system, it verified the existence of cascade-like propagation of correlation from long-wavelength components to shorter-wavelength components.
First, the use of the wavelet transform is essential for our results. It seems difficult to obtain similar results using a method that flattens out the spatial information completely, like the Fourier transform. As can clearly be seen from the definition of the inner product given in Eq. (13), our method captures the simultaneous fluctuation of two fields at the same position in space. For example, if we use a different definition of the inner product that loses the spatial information by taking spatial average of each field first, we cannot observe the correlation between scales so clearly as in Figs. 10 and 11. Similarly, if we change the definition of the inner product so as to multiply the fields with their origins shifted by half the system size, the correlation matrix almost vanishes. (This tendency is more remarkable for the NCGL equation. In the case of the LCGL equation, though not so clear, we obtain similar asymmetric correlation matrices with these modified definitions of the inner product. This is because with the values of $`D`$ and $`L`$ we used in our calculation, the ratio of the diffusion constant to the system size is considerably large, and some coherence over the whole system still remains. If we further decrease the value of $`D`$, the correlation matrix tends to vanish as in the case of the NCGL equation.)
Conversely, the fact that we succeeded in visualization using the definition of the inner product given in Eq. (13) implies that the cascade process occurs in spatially localized regions in our system. However, it is frequently seen in other systems that the spatial structure of the pattern collapses while moving constantly, as in the case of convective instability. In order to detect such phenomena, it would be possible to generalize the definition of the inner product in such a way that the distance between the origins of the fields is increased with $`\mathrm{\Delta }t`$.
Precisely speaking, the result visualized by our method is not a flow of causality, but merely a flow of correlation. Namely, it indicates that fluctuation at the reference time at some scale varies in unison with fluctuation at the later time at some other scale, but it does not necessarily mean that the former one actually affects the latter one. Of course, it is natural to interpret it as causality in our cases. But in order to be exact, it will be possible to perturb a certain mode at some scale, e.g. by using a sinusoidal wave, and visualize the propagation of its aftereffect to prove that it is indeed a causal relationship.
Finally, though we did not give a detailed discussion in this paper, we will be able to investigate the dynamical process of the spatio-temporal chaos of non-locally coupled oscillators in more detail by quantitatively analyzing the results obtained by our method. For example, from the peak positions of the $`E_{j,j+1}`$ curves shown in Fig.13, we will be able to study the dependence of characteristic time scale of the underlying cascade process on the coupling strength. Such a study will also be an interesting future subject.
## ACKNOWLEDGMENTS
H. N. gratefully acknowledges M. Hayashi for useful discussion, and S. Amari for providing an excellent environment for scientific study. He also thanks K. Ishioka, Y. Taniguchi, C. Liu, S. Kato, Y. Kitano, and N. Kobayashi for their warm hospitality during his stay in University of Tokyo, and T. Takami and T. Mizuguchi for providing a nice visualization tool. This work is partly supported by the JSPS research fellowship for young scientists, and partly by RIKEN.
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# Local Spin Correlations in Heisenberg Antiferromagnets
## I Introduction
This paper deals with the problem of calculating correlation functions, at $`T=0`$, for the $`S=1/2`$ Heisenberg antiferromagnet
$$H=2J\underset{ij}{}𝐒_i𝐒_j$$
(1)
where the sum is over all nearest neighbor pairs. We consider explicitly the linear chain ($`d=1`$), the square lattice ($`d=2`$) and the simple cubic lattice ($`d=3`$). The correlation functions (correlators) are defined as
$$C(r)4𝐒_i𝐒_{i+r}_0=C_l(r)+2C_t(r)$$
(2)
where the average is a ground state expectation value, $`r`$ is the distance between sites in units of the lattice parameter, and the factor $`4`$ is included for convenience. It is also convenient to separate the correlator into a longitudinal part $`C_l(r)`$ and a transverse part $`C_t(r)`$
$`C_l(r)`$ $``$ $`4S_i^zS_{i+r}^z_0`$ (4)
$`C_t(r)`$ $``$ $`2S_i^xS_{i+r}^x+S_i^yS_{i+r}^y_0`$ (5)
$`=`$ $`S_i^+S_{i+r}^{}+S_i^{}S_{i+r}^+_0`$ (6)
$`C_l`$ and $`C_t`$ will differ if the Hamiltonian is generalized to include Ising anisotropy
$$H=2J\underset{ij}{}[S_i^zS_j^z+\lambda (S_i^xS_j^x+S_i^yS_j^y)]$$
(7)
as we shall do, or if the ground state of the isotropic Hamiltonian exhibits spontaneous symmetry breaking. We shall see that this occurs in dimension $`d>1`$.
The correlators characterize the nature of the ground state of the system, and hence an accurate knowledge of their values can be important for testing approximate analytic theories. Surprisingly, apart from the 1-d case, knowledge of their values is limited.
We use the method of linked cluster expansions in which the Hamiltonian is written as
$$H=H_0+\lambda V$$
(8)
with the Ising part taken as the unperturbed Hamiltonian and the remainder as a perturbation. To improve the convergence of the series, we also include a local staggered field term $`t_i(1)^iS_i^z`$ in $`H_0`$, and subtract it from $`V`$, and adjust the strength $`t`$ to get best convergence in the series. The basic idea of the method has been discussed before so we only give brief details here. To compute series for $`C(r)`$ in powers of $`\lambda `$ we add a field term to $`H`$
$$H=H_0+\lambda V+h\underset{i}{}𝐒_i𝐒_{i+r}$$
(9)
compute the ground state energy in the form
$$_0(\lambda ,h)=E_0(\lambda )+hNC(r)/4+O(h^2)$$
(10)
and hence extract series in $`\lambda `$ for $`C(r)`$. For the longitudinal correlator the field term is $`h_iS_i^zS_{i+r}^z`$. Examples are give in the following sections.
Two important questions are not addressed in this work:
(i) the behaviour of correlators at large distances, and the asymptotic behaviour and associated critical exponents.
(ii) correlators at finite temperature.
## II The 1-Dimensional Case
The $`S=\frac{1}{2}`$ antiferromagnetic Heisenberg chain is exactly solvable by the Bethe ansatz, and the ground state energy and elementary excitations are given by simple analytic expressions. However the wavefunction is sufficiently complex that little exact information is available on correlators. In fact only the first two are known exactly, and are
$`𝐒_i𝐒_{i+1}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(14\mathrm{ln}2)=0.443147\mathrm{}`$ (11)
$`𝐒_i𝐒_{i+2}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(116\mathrm{ln}2+9\zeta (3))=0.182039\mathrm{}`$ (13)
The first result comes from the ground state energy, while the second is obtained via the strong coupling limit of the Hubbard model. There is no spontaneous symmetry breaking in the ground state of the isotropic spin chain, and hence longitudinal and transverse correlators are equal, and obtainable directly from (2).
The first serious attempts to obtain further results for the antiferromagnetic chain were by Bonner and Fisher, who used exact diagonalizations for systems up to $`N=10`$ spins, and by Kaplan and co-workers who extended this to $`N=18`$. These results suffer from large finite-size effects, and need to be extrapolated to the therodynamic limit via a finite-size scaling ansatz. In this way Kaplan et al. estimated the value of correlators up to 8th neighbors, with confidence limits of about 1% in $`C(8)`$. We show these values in Table I. Subsequently Lin and Campell extended the exact diagonalizations to $`N=30`$. By use of the empirical scaling relation
$$C_N(r)=C_{\mathrm{}}(r)f(r/N)$$
(14)
with
$$f(y)=[1+0.28822\mathrm{sinh}^2(1.673y)]^{1.75}$$
(15)
they estimated correlators up to $`r=15`$, i.e. 15th neighbors. However the accuracy of this scaling is perhaps doubtful since it is known that there are logarithmic terms which slow convergence.
The development of the Density Matrix Renormalization Group (DMRG) method allows much longer chains to be treated with high numerical accuracy. Hallberg et al. have used DMRG to compute correlators for Heisenberg chains up to $`N=70`$ spins, with a scaling function similar to (14) used to extrapolate to the thermodynamic limit. The data were shown to be consistent with the asymptotic behaviour
$$C(r)(1)^r(\mathrm{ln}r)^{1/2}/r$$
(16)
predicted by field theory.
We have described the series method briefly in the Introduction. Using this approach we have computed expansions in $`\lambda `$, for both the total and longitudinal correlators for distances $`r=1,2,\mathrm{},10`$. The maximum order is $`\lambda ^{24}`$ for $`r=1`$ and $`\lambda ^{16}`$ for $`r=10`$. We note that the longitudinal correlators, and the total correlators for $`r`$ even, contain only even powers of $`\lambda `$. We also note that the series are rather erratic, both in sign and magnitude of the coefficients. This had already been noted by Walker who expanded the ground state energy, and hence $`C(1)`$, to order $`\lambda ^{14}`$. Rather than quote all series here we will make them available to any reader on request. Table II shows the coefficients for the series for $`C(4)`$, $`C_l(4)`$ and $`C_t(4)`$. We note that, as expected, the series for the transverse correlator $`C_t(r)`$ starts with a term $`\lambda ^r`$.
The series have been evaluated for fixed $`\lambda `$ by means of integrated differential approximants, and the values of correlators for $`r=1,2,\mathrm{},6`$ are shown in Figure 1. The analysis becomes less precise as the weakly singular point $`\lambda =1`$ is approached. We also show in the Figure the extrapolated exact diagonalization results. As can be seen from the Figure, and from Table I, the agreement is very good. It is clear that in 1-dimension the series method is not able to match the precision of either scaled finite lattice or DMRG results, but in higher dimension these latter methods are not competitive. Furthermore, as we shall show, the series analysis can be made more precise in $`d2`$ because the stronger singularity at $`\lambda =1`$ can be removed by a transformation, and the Ising expansions used here are more suitable for $`d2`$, where the ground state has long range Neél order.
We should mention here also the work of Singh et al. who used exactly the same method as ours to compute the structure factors
$$S_{zz}=\frac{1}{4}\underset{r=\mathrm{}}{\overset{\mathrm{}}{}}[(1)^rC_l(r)4<S_0^z>^2]$$
(17)
and
$$S_+=\frac{1}{2}\underset{r=\mathrm{}}{\overset{\mathrm{}}{}}(1)^rC_t(r)$$
(18)
for the $`S=\frac{1}{/}2`$ antiferromagnetic chain. Our correlator series, when summed, agree with their results.
## III The 2-Dimensional Case
There has been much interest, in recent years, in the nature of the ground state of the Heisenberg antiferromagnet on the square lattice. There is considerable evidence, from exact diagonalizations and quantum Monte Carlo calculations that the ground state breaks rotational symmetry, giving rise to a staggered magnetization in some direction. This is generally referred to as a quantum Néel state, with Néel type order reduced to approx. 60% of its classical value by quantum fluctuations. The situation is summarized in recent reviews. The focus has generally been on the ground state energy and staggered magnetization, although some short range correlators have also been computed.
In any finite system there can be no spontaneous symmetry breaking and hence the exact diagonalization and Monte Carlo studies cannot distinguish between the longitudinal and transverse correlators for the isotropic case. Furthermore if $`C_l`$ or $`C_t`$ are computed by these methods the values will not yield correct results for the thermodynamic limit, where $`C_lC_t`$.
Other approaches, such as spin wave (SW) theory, variational methods, or perturbation series about the Ising limit start from a broken symmetry state, which is preserved during the calculation. It seems highly likely, although we know of no proof, that these approaches will yield the correct symmetry-broken state of the infinite isotropic system.
We have computed series expansions for a number of local correlators for the square lattice $`S=\frac{1}{2}`$ antiferromagnet. The expansions start from the Ising limit and are carried through order 14,9,9,9,7 for $`C(𝐫)`$, $`C_l(𝐫)`$ with $`𝐫=(0,1)`$, (1,1), (2,0), (3,0), (4,0) The series coefficients (for $`t=0`$) are given Table III.
In analysing the series it is advantageous to transform to a new variable
$$\delta =1(1\lambda )^{1/2},$$
(19)
to remove the singularity at $`\lambda =1`$. Spin wave theory predicts a square root singularity of this type. This transformation was first proposed by Huse and was also used in earlier work on the square lattice case. We then use both integrated first-order inhomogeneous differential approximants and Padé approximants to extrapolate the series to the isotropic point $`\delta =1`$ ($`\lambda =1`$). The results are shown as functions of $`\lambda `$ in Figure 2 for $`𝐫`$=(1,0), (1,1). We also show the transverse correlator, obtained from Eq. 6. In the Ising limit the total and longitudinal correlators are equal and the transverse correlator is zero. As we increase the transverse coupling, the longitudinal correlators decrease in magnitude while, as expected, the transverse correlators increase, while the total correlator increases in magnitude for nearest neighbors, but is reduced for second neighbors. The behaviour of further correlators is similar, and is not shown. It is also noteworthy that at $`\lambda =1`$ the longitudinal and transverse correlators remain unequal, reflecting the symmetry broken ground state. In Fig. 3 we show a comparison between our series results and other methods for the nearest neighbor correlators. For small $`\lambda `$ all methods are in close agreement, but near the isotropic point linear spin wave theory become poor for longitudinal (and transverse) correlators, whereas exact finite lattice diagonalizations have longitudinal and transverse correlators equal at $`\lambda =1`$. Third order spin-wave theory is much better, being almost indistinguishable from the series results over the whole range of $`\lambda `$.
In Table IV we give numerical estimates of all the correlators at the isotropic point, obtained by our series method and by exact diagonalization/Monte Carlo on finite lattices and linear spin wave theory. We believe that the expressions in Ref. contain minor errors, and should read, for 0 and $`𝐫`$ on the same sublattice
$`𝐒_0𝐒_r`$ $`=`$ $`S^2+S\left(1{\displaystyle \frac{2}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{1\mathrm{cos}𝐤𝐫}{\sqrt{1\lambda ^2\gamma _k^2}}}\right)+\mathrm{}`$ (20)
$`S_0^zS_r^z`$ $`=`$ $`S^2+S\left(1{\displaystyle \frac{2}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{1}{\sqrt{1\lambda ^2\gamma _k^2}}}\right)+\mathrm{}`$ (22)
while for 0 and $`𝐫`$ on different sublattices:
$`𝐒_0𝐒_r`$ $`=`$ $`S^2S\left(1{\displaystyle \frac{2}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{1\lambda \gamma _k\mathrm{cos}𝐤𝐫}{\sqrt{1\lambda ^2\gamma _k^2}}}\right)+\mathrm{}`$ (23)
$`S_0^zS_r^z`$ $`=`$ $`S^2+S\left(1{\displaystyle \frac{2}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{1}{\sqrt{1\lambda ^2\gamma _k^2}}}\right)+\mathrm{}`$ (25)
where the notation is as in Ref. , and $`\lambda `$ is the anisotropy parameter.
We note from Eqn. LABEL:Csame and LABEL:Cdiff and Table IV that first order spin-wave theory gives a longitudinal correlator which is independent of distance, clearly an artifact of the approximation. The total correlator is however very consistent with the series results. The picture is considerably improved in higher order spin-wave theory where, for example, 3rd order spin wave theory gives 3-figure agreement with series for all of $`C`$, $`C_l`$ and $`C_t`$ for nearest neighbors. We have not attempted to carry this out for further neighbors, and are unaware of any work along these lines.
## IV The 3-Dimensional Case
We have used the same series approach to calculate correlators for the $`S=\frac{1}{/}2`$ antiferromagnet on the simple cubic lattice. The magnetically ordered ground state will again be reflected in a difference between longitudinal and transverse correlators at the isotropic limit.
Expansions, starting from the Ising limit, have been obtained for $`C(𝐫)`$, $`C_l(𝐫)`$ for the five cases $`𝐫=`$(1,0,0), (1,1,0), (2,0,0), (3,0,0), (4,0,0) to order 12,7,9,7,7 respectively. We have again used a staggered field term $`t_i(1)^iS_i^z`$ to improve convergence. The series coefficients (for $`t=0`$) are given in Table V. The series is extrapolated in a similar way as that for the square lattice. Figure 4 shows the nearest and next-nearest neighbor correlators as functions of the anisotropy parameters. This is qualitatively similar to Figure 2, but clearly shows that in 3-dimensions transverse correlators are reduced and the difference between transverse and longitudinal correlators is increased for all values of the anisotropy parameter.
In Table VI we give numerical estimates of all correlators at the isotropic point, and a comparison with 1st order spin-wave theory. It is apparent that the correlators fall off more slowly with distance than in the 2-dimensional case, reflecting the greater stability of antiferromagnetic long-range order in the ground state in 3-dimensions. It is also apparent that the transverse correlators are, relatively, much weaker in 3-dimensions, consistent with weaker quantum fluctuations. Linear spin-wave theory gives reasonable results for the total correlators, but again suffers from the defect of having longitudinal correlators independent of distance. Third-order spin-wave theory gives results for nearest neighbor correlators in excellent agreement with the series results.
## V Conclusions
We have used series methods to obtain numerical estimates for short-distance ground state correlation functions for the $`S=\frac{1}{/}2`$ Heisenberg antiferromagnet on square and simple cubic lattices. Despite their importance in characterising the nature of the antiferromagnetic ground state, there appears to have been little previous work on the subject.
The series approach is able to provide rather precise estimates for correlators up to at least 4 lattice spacings. The results reflect the known breaking of rotational symmetry in the ground state, in that longitudinal and transverse correlators remain unequal even in the isotropic Hamiltonian limit. Exact diagonalizations and Monte Carlo calculations on finite lattices are unable to account for this and hence will not yield correct estimates for longitudinal and transverse correlators separately. In 3-dimensions no results are available from diagonalizations or quantum Monte Carlo, beyond nearest neighbors.
We have shown that first-order spin wave theory gives rather poor estimates but 2nd and 3 rd order spin wave theory gives excellent agreement with series results for nearest-neighbor correlators. Higher order spin wave results have not been obtained for further correlators, to our knowledge.
As a test of the method we also computed correlation series for the 1-dimensional case. The results were quantitatively accurate, but less precise than the DMRG method.
This approach can also be used to calculate correlators for more complex models involving competing interactions. For example, we have studied the $`J_1J_2`$ model, which has a quantum critical point at $`J_2/J_10.38`$, where the Néel order is destroyed and the system enters a magnetically disordered spin-liquid phase. We find that the difference between longitudinal and transverse correlators remains nonzero in the Neél phase, but vanishes at the quantum critical point, indicating a restoration of full rotational symmetry in the ground state at that point. We expect this method to prove useful in other problems of this type.
###### Acknowledgements.
This work forms part of a research project supported by a grant from the Australian Research Council. We thank Dr. Oleg Sushkov for stimulating our interest in this problem. The computation has been performed on Silicon Graphics Power Challenge and Convex machines. We thank the New South Wales Centre for Parallel Computing for facilities and assistance with the calculations.
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# Theory of finite pseudoalgebras
## 1. Introduction
Since the seminal papers of Belavin, Polyakov and Zamolodchikov \[BPZ\] and of Borcherds \[Bo1\] there has been a great deal of work towards understanding the algebraic structures underlying the notion of operator product expansion (OPE) of chiral fields in conformal field theory.
In physics literature the OPE of local chiral fields $`\phi `$ and $`\psi `$ is written in the form \[BPZ\]:
(1.1)
$$\phi (z)\psi (w)=\underset{j\mathrm{}}{}\frac{\phi (w)_{(j)}\psi (w)}{(zw)^{j+1}},$$
where $`\phi (w)_{(j)}\psi (w)`$ are some new fields, which may be viewed as bilinear products of fields $`\phi `$ and $`\psi `$ for all $`j`$ (see e.g. \[K2\] for a rigorous interpretation of (1.1)). If now $`V`$ is a space of pairwise local chiral fields which contains $`1`$, is invariant with respect to the derivative $`=_w`$, and is closed under all $`j`$th products, $`j`$, we obtain an algebraic structure which physicists (respectively mathematicians) call a chiral (respectively vertex) algebra. In more abstract terms, $`V`$ is a module over $`[]`$ with a marked element $`1`$ and infinitely many bilinear over $``$ products $`\phi _{(j)}\psi `$, $`j`$, satisfying a certain system of identities, first written down by Borcherds \[Bo1\]. (An equivalent system of axioms, which is much easier to verify, may be found in \[K2\].)
One of the important features of the OPE (1.1) is that its singular part encodes the commutation relations of fields, namely one has (see e.g. \[K2\]):
(1.2)
$$[\phi (z),\psi (w)]=\underset{j0}{}\left(\phi (w)_{(j)}\psi (w)\right)_w^j\delta (zw)/j!,$$
where $`\delta (zw)=_nz^nw^{n1}`$ is the delta-function. This leads to the notion of a Lie conformal algebra, which is a $`[]`$-module with $``$-bilinear products $`\phi _{(j)}\psi `$ for all non-negative integers $`j`$, subject to certain identities \[K2\]. In order to write down these identities in a compact form, it is convenient to consider the formal Fourier transform of (1.2), called the $`\lambda `$-bracket (where $`\lambda `$ is an indeterminate):
$$[\phi _\lambda \psi ]=\underset{j0}{}\frac{\lambda ^j}{j!}(\phi _{(j)}\psi ).$$
Then a *Lie conformal algebra* $`L`$ is defined as a $`[]`$-module endowed with a $``$-linear map
$$LL[\lambda ]L,ab[a_\lambda b]$$
satisfying the following axioms \[DK\] $`(a,b,cL)`$:
$`\begin{array}{cc}\text{(sesquilinearity)}\hfill & [a_\lambda b]=\lambda [a_\lambda b],[a_\lambda b]=(+\lambda )[a_\lambda b],\hfill \\ \text{(skew-commutativity)}\hfill & [b_\lambda a]=[a_\lambda b],\hfill \\ \text{(Jacobi identity)}\hfill & [a_\lambda [b_\mu c]]=[[a_\lambda b]_{\lambda +\mu }c]+[b_\mu [a_\lambda c]].\hfill \end{array}`$
In the past few years a structure theory \[DK\], representation theory \[CK, CKW\] and cohomology theory \[BKV\] of finite (i.e., finitely generated as $`[]`$-modules) Lie conformal algebras have been worked out. For example, one of the main results of \[DK\] states that any finite simple Lie conformal algebra is isomorphic either to the *Virasoro* conformal algebra:
(1.3)
$$\mathrm{Vir}=[]\mathrm{},[\mathrm{}_\lambda \mathrm{}]=(+2\lambda )\mathrm{}$$
or to the *current* conformal algebra associated to a simple finite-dimensional Lie algebra $`𝔤`$:
(1.4)
$$\mathrm{Cur}𝔤=[]𝔤,[a_\lambda b]=[a,b],a,b𝔤.$$
The objective of the present paper is to develop a theory of “multi-dimensional” Lie conformal algebras, i.e. a theory where the algebra of polynomials $`[]`$ is replaced by a “multi-dimensional” associative algebra $`H`$. In order to explain the definition, let us return to the singular part (1.2) of the OPE. Choosing a set of generators $`a^i`$ of the $`[]`$-module $`L`$, we can write:
$$[a_\lambda ^ia^j]=_kQ_k^{ij}(\lambda ,)a^k,$$
where $`Q_k^{ij}`$ are some polynomials in $`\lambda `$ and $``$. The corresponding singular part of the OPE is:
$$[a^i(z),a^j(w)]=_kQ_k^{ij}(_w,_t)(a^k(t)\delta (zw))|_{t=w}.$$
Letting $`P_k^{ij}(x,y)=Q_k^{ij}(x,x+y)`$, we can rewrite this in a more symmetric form:
(1.5)
$$[a^i(z),a^j(w)]=_kP_k^{ij}(_z,_w)(a^k(w)\delta (zw)).$$
We thus obtain an $`H=[]`$-bilinear map (i.e., a map of $`HH`$-modules):
$$LL(HH)_HL,ab[ab]$$
(where $`H`$ acts on $`HH`$ via the comultiplication map $`\mathrm{\Delta }()=1+1`$), defined by
$$[a^ia^j]=_kP_k^{ij}(1,1)_Ha^k.$$
Hence the notion of $`\lambda `$-bracket $`[a_\lambda b]`$ is equivalent to the notion of the $``$-bracket $`[ab]`$ introduced by Beilinson and Drinfeld \[BD\], the relation between the two brackets being given by letting $`\lambda =1`$. For example, the Virasoro conformal algebra (1.3) corresponds to the Virasoro pseudoalgebra
(1.6)
$$\mathrm{Vir}=[]\mathrm{},[\mathrm{}\mathrm{}]=(11)_{[]}\mathrm{}.$$
It is natural to introduce the general notion of a conformal algebra as a $`[]`$-module $`L`$ endowed with a $``$-linear map $`LL[\lambda ]L`$, $`aba_\lambda b`$ satisfying the sesquilinearity property:
$$(a)_\lambda b=\lambda (a_\lambda b),a_\lambda (b)=(+\lambda )(a_\lambda b).$$
Such a conformal algebra is called associative (respectively commutative) if
$$a_\lambda (b_\mu c)=(a_\lambda b)_{\lambda +\mu }c(\text{respectively}b_\lambda a=a_\lambda b),$$
and the $`\lambda `$-product of an associative conformal algebra defines a $`\lambda `$-bracket
$$[a_\lambda b]=a_\lambda bb_\lambda a,$$
making it a Lie conformal algebra \[K4, DK\].
As above, we have the equivalent notion of a $``$-product on an $`H=[]`$-module $`L`$, which is an $`H`$-bilinear map
(1.7)
$$LL(HH)_HL,abab.$$
Now it is clear that the notion of a $``$-product can be defined by (1.7) for any Hopf algebra $`H`$ by making use of the comultiplication $`\mathrm{\Delta }:HHH`$ to define $`(HH)_HL`$. A *pseudoalgebra* is a (left) $`H`$-module $`L`$ endowed with an $`H`$-bilinear map (1.7). The name is motivated by the fact that this is an algebra in a pseudotensor category (introduced in \[L\], \[BD\]). Accordingly, the $``$-product will be called a pseudoproduct.
One is able to define a pseudoproduct as soon as a structure of a bialgebra is given on $`H`$. However, in order to generalize the equivalence of a pseudoalgebra and an $`H`$-conformal algebra structure on an $`H`$-module $`L`$, we need $`H`$ to be a Hopf algebra. In this case any element of $`HH`$ can be uniquely written as a (finite) sum:
$$_i(h_i1)\mathrm{\Delta }(f_i),\text{where }h_i\text{ are linearly independent.}$$
Hence the pseudoproduct on $`L`$ can be written in the form:
(1.8)
$$ab=_i(h_i1)_Hc_i.$$
The corresponding $`H`$-conformal algebra structure is then a $``$-linear map $`LLHL`$ given by
(1.9)
$$ab=_ih_ic_i.$$
Every element $`x`$ of $`H^{}`$ then defines an $`x`$-product $`LLL`$:
(1.10)
$$a_xb=_ix,S(h_i)c_i,$$
where $`S`$ is the antipode of $`H`$.
The $`H`$-bilinearity property of the pseudoproduct (1.8) is, of course, easily translated to certain sesquilinearity properties of the products (1.9) and (1.10). In particular, in the case $`H=[]`$, the product (1.9) is the $`\lambda `$-product if we let $`\lambda =`$, and the product (1.10) for $`x=t^j`$ is the $`j`$th product described above, where $`H^{}[[t]]`$, $`t,=1`$. The equivalence of these three structures (discussed in Section 9) is very useful in the study of pseudoalgebras.
In order to define associativity of a pseudoproduct, we extend it from $`LLH^2_HL`$ to $`(H^2_HL)LH^3_HL`$ and to $`L(H^2_HL)H^3_HL`$ by letting:
$`(f_Ha)b`$ $`=_i(f1)(\mathrm{\Delta }\mathrm{id})(g_i)_Hc_i,`$
$`a(f_Hb)`$ $`=_i(1f)(\mathrm{id}\mathrm{\Delta })(g_i)_Hc_i,\text{where}ab=_ig_i_Hc_i.`$
Then the associativity property is given by the usual equality (in $`H^3_HL`$):
$$(ab)c=a(bc).$$
The easiest example of a pseudoalgebra is a current pseudoalgebra, defined as follows. Let $`H^{}`$ be a Hopf subalgebra of $`H`$ and let $`A`$ be an $`H^{}`$-pseudoalgebra (for example, if $`H^{}=`$, then $`A`$ is an ordinary algebra over $``$). Then the associated *current* $`H`$-pseudoalgebra is $`\mathrm{Cur}A=H_H^{}A`$ with the pseudoproduct
$$(f_H^{}a)(g_H^{}b)=((fg)_H1)(ab).$$
The $`H`$-pseudoalgebra $`\mathrm{Cur}A`$ is associative iff the $`H^{}`$-pseudoalgebra $`A`$ is.
The most important example of an associative $`H`$-pseudoalgebra is the pseudoalgebra of all pseudolinear endomorphisms of a finitely generated $`H`$-module $`V`$, which is denoted by $`\mathrm{Cend}V`$ (see Section 10). A pseudolinear endomorphism of $`V`$ is a $``$-linear map $`\varphi :V(HH)_HV`$ such that
$`\varphi (hv)`$ $`=((1h)_H1)\varphi (v),hH,vV.`$
The space $`\mathrm{Cend}V`$ of all such $`\varphi `$ becomes a (left) $`H`$-module if we define
$`(h\varphi )(v)`$ $`=((h1)_H1)\varphi (v).`$
The definition of a pseudoproduct on $`\mathrm{Cend}V`$ is especially simple when $`V`$ is a free $`H`$-module, $`V=HV_0`$, where $`V_0`$ is a finite-dimensional vector space over $``$ with a trivial action of $`H`$. Then $`\mathrm{Cend}V`$ is isomorphic to $`HH\mathrm{End}V_0`$, with $`H`$ acting by left multiplication on the first factor, with the following pseudoproduct:
$$(faA)(gbB)=_i(fga_i^{})_H(1ba_i^{\prime \prime }AB),$$
where $`\mathrm{\Delta }(a)=_ia_i^{}a_i^{\prime \prime }`$.
The main objects of our study are Lie pseudoalgebras. The corresponding pseudoproduct is conventionally called pseudobracket and denoted by $`[ab]`$. Given an associative pseudoalgebra with pseudoproduct $`ab`$ we may give it a structure of a Lie pseudoalgebra by defining the pseudobracket
$$[ab]=ab(\sigma _H\mathrm{id})ba,$$
where $`\sigma :HHHH`$ is the permutation of factors. It is immediate to see that this pseudobracket satisfies the following skew-commutativity and Jacobi identity axioms:
(1.11) $`[ba]`$ $`=(\sigma _H\mathrm{id})[ab],`$
(1.12) $`[a[bc]]`$ $`=[[ab]c]+((\sigma \mathrm{id})_H\mathrm{id})[b[ac]].`$
It is important to point out here that the above pseudobracket and both identities are well defined, provided that the Hopf algebra $`H`$ is cocommutative. A pseudoalgebra with pseudoproduct $`[ab]`$ satisfying identities (1.11) and (1.12) is called a *Lie pseudoalgebra*. We will always assume that $`H`$ is cocommutative when talking about Lie pseudoalgebras. Of course, the simplest examples of Lie pseudoalgebras are $`\mathrm{Cur}A`$, where $`A`$ is a $`H^{}(H)`$ Lie pseudoalgebra ($`=`$ Lie algebra if $`H^{}=`$). It is needless to say that in the case $`H=[]`$, $`\mathrm{\Delta }()=1+1`$, the $`H`$-conformal algebras associated to Lie pseudoalgebras are nothing else but the Lie conformal algebras discussed above.
We will explain now the connection of the notion of a Lie pseudoalgebra to the more classical notion of a differential Lie algebra studied in \[R1\]\[R4\], \[C\], \[NW\] and many other papers (see Section 7). Let $`Y`$ be a commutative associative algebra over $``$ with compatible left and right actions of the Hopf algebra $`H`$. Then, given a Lie pseudoalgebra $`L`$, we let $`𝒜_YL=Y_HL`$ with the obvious left $`H`$-module structure and the following Lie algebra (over $``$) structure:
$$[(x_Ha),(y_Hb)]=_i(xf_i)(yg_i)_Hc_i\text{if}[ab]=_i(f_ig_i)_Hc_i.$$
Provided that $`L`$ is a free $`H`$-module, the Lie algebra $`𝒜_YL`$ is a free $`Y`$-module, hence $`𝒜_YL`$ is a differential Lie algebra in the sense of \[NW\]. The most classical case is again $`H=[]`$, when $`Y`$ is simply a commutative associative algebra with a (left and right) derivation $``$, and we get the differential Lie algebras of Ritt \[R1\]\[R4\]. Thus, the notion of a Lie pseudoalgebra is reminiscent of the notion of a group scheme: each Lie pseudoalgebra $`L`$, which is free as an $`H`$-module, gives rise to a functor $`𝒜`$ from the category of commutative associative algebras with compatible left and right actions of $`H`$ to the category of differential Lie algebras ($`=`$ category of formal differential groups).
For example, the functor $`𝒜`$ corresponding to the Virasoro pseudoalgebra (1.6) associates to any commutative associative algebra $`Y`$ with a derivation the differential Lie algebra $`Y`$ with bracket $`[u,v]=uv^{}u^{}v`$, called the substitutional Lie algebra by Ritt. The current pseudoalgebra $`\mathrm{Cur}𝔤`$, where $`𝔤`$ is a Lie algebra over $``$, associates to $`Y`$ the obvious differential Lie algebra $`Y𝔤`$. Thus, a result of \[DK\] asserts that any simple finite differential Lie algebra with “constant coefficients” is isomorphic either to the substitutional Lie algebra or to $`Y𝔤`$ where $`𝔤`$ is a simple finite-dimensional Lie algebra. In the rank $`1`$ case, but without the constant coefficients assumption, this is the main result of \[R1\].
The main tool in the study of pseudoalgebras is the *annihilation algebra* $`𝒜_XL`$, where $`X=H^{}`$ is the associative algebra dual to the coalgebra $`H`$. We find it remarkable that the annihilation algebra of the associative pseudoalgebra $`\mathrm{Cend}H=HH`$ is nothing else but the Drinfeld double (with the obvious comultiplication) of the Hopf algebra $`H`$. Note that in the associative case $`Y`$ need not be commutative in order to define the functor $`𝒜_Y`$, but in the Lie algebra case it must be. So, in order to construct the annihilation Lie algebra we again use cocommutativity of $`H`$.
Recall that, by Kostant’s Theorem 2.1, any cocommutative Hopf algebra $`H`$ is a smash product of a group algebra $`[\mathrm{\Gamma }]`$ and the universal enveloping algebra $`U(𝔡)`$ of a Lie algebra $`𝔡`$. In Sections 5 and 13.7 we show that the theory of pseudoalgebras over a smash product of $`[\mathrm{\Gamma }]`$ and any Hopf algebra $`H`$ reduces to that over $`H`$. This allows us in many cases to assume, without loss of generality, that $`H`$ is the universal enveloping algebra of a Lie algebra $`𝔡`$.
However, for most of our results we have to assume that $`𝔡`$ is finite dimensional. In this case the algebra $`H=U(𝔡)`$ is Noetherian, and the annihilation algebra $`𝒜_XL`$ is linearly compact, provided that $`L`$ is *finite* (i.e., finitely generated as an $`H`$-module). Recall that a topological Lie algebra is called linearly compact if its underlying topological vector space is a topological product of finite-dimensional vector spaces with the discrete topology (see Section 6).
In Section 11 we prove “reconstruction” theorems, which claim that, under some mild assumptions, a Lie pseudoalgebra is completely determined by its annihilation Lie algebra along with the action of $`𝔡`$. This reduces the classification of finite simple Lie pseudoalgebras to the well developed structure theory of linearly compact Lie algebras, which goes back to E. Cartan (see \[G1, G2\] and Section 6).
We turn now to examples of finite Lie pseudoalgebras beyond the rather obvious examples of current Lie pseudoalgebras. The first example is the generalization of the Virasoro pseudoalgebra (1.6) defined for $`H=[]`$ (which is the universal enveloping algebra of the $`1`$-dimensional Lie algebra) to the case $`H=U(𝔡)`$, where $`𝔡`$ is any finite-dimensional Lie algebra. This is the Lie pseudoalgebra $`W(𝔡)=H𝔡`$ with pseudobracket
$$[(1a)(1b)]=(11)_H(1[a,b])+(b1)_H(1a)(1a)_H(1b).$$
Since the associated annihilation algebra $`𝒜_XW(𝔡)X𝔡`$ is isomorphic to the Lie algebra of formal vector fields on the Lie group $`D`$ with Lie algebra $`𝔡`$, it is natural to call $`W(𝔡)`$ the pseudoalgebra of all vector fields. In fact we develop (in Section 8) a formalism of pseudoforms similar to the usual formalism of differential forms, which may be viewed as the beginning of a “pseudo differential geometry”.
This allows us to define the remaining three series of finite simple Lie pseudoalgebras: $`S(𝔡,\chi )`$, $`H(𝔡,\chi ,\omega )`$ and $`K(𝔡,\theta )`$. The annihilation algebras of the simple Lie pseudoalgebras $`W(𝔡)`$, $`S(𝔡,\chi )`$, $`H(𝔡,\chi ,\omega )`$ and $`K(𝔡,\theta )`$ are isomorphic to the four series of Lie–Cartan linearly compact Lie algebras $`W_N`$, $`S_N`$, $`P_N`$ (which is an extension of $`H_N`$ by a $`1`$-dimensional center) and $`K_N`$, where $`N=dim𝔡`$. However the Lie pseudoalgebras $`S(𝔡,\chi )`$, $`H(𝔡,\chi ,\omega )`$ and $`K(𝔡,\theta )`$ depend on certain parameters $`\chi ,\omega `$ and $`\theta `$, due to inequivalent actions of $`𝔡`$ on the annihilation algebra. The parameter $`\chi `$ is a $`1`$-dimensional representation of $`𝔡`$, i.e., $`\chi 𝔡^{}`$ such that $`\chi ([𝔡,𝔡])=0`$. The parameter $`\omega `$ is an element of $`𝔡^{}𝔡^{}`$ such that $`\omega ^{N/2}0`$ and $`\mathrm{d}\omega +\chi \omega =0`$ in the case $`H(𝔡,\chi ,\omega )`$, when $`N`$ is even. The parameter $`\theta 𝔡^{}`$ is such that $`\theta (\mathrm{d}\theta )^{(N1)/2}0`$ in the case $`K(𝔡,\theta )`$, when $`N`$ is odd. In the cases $`H(𝔡,\chi ,\omega )`$, $`K(𝔡,\theta )`$, these parameters are in one-to-one correspondence with “nondegenerate” skew-symmetric solutions $`\alpha =r+s11s`$ ($`r𝔡𝔡`$, $`s𝔡`$) of a modification of the classical Yang–Baxter equation, which is a special case of the dynamical classical Yang–Baxter equation (see \[Fe, ES\]).
The central result of the paper is the classification of finite simple Lie pseudoalgebras over the Hopf algebra $`H=U(𝔡)`$. As usual, a Lie pseudoalgebra $`L`$ is called *simple* if it is nonabelian (i.e., $`[LL]0`$) and its only ideals are $`0`$ and $`L`$. Our Theorem 13.10 states that any such Lie pseudoalgebra is isomorphic either to a current pseudoalgebra $`\mathrm{Cur}𝔤=\mathrm{Cur}_{}^H𝔤`$ over a simple finite-dimensional Lie algebra $`𝔤`$, or to a current pseudoalgebra $`\mathrm{Cur}_H^{}^HL^{}`$ over one of the Lie pseudoalgebras $`L^{}=W(𝔡^{})`$, $`S(𝔡^{},\chi ^{})`$, $`H(𝔡^{},\chi ^{},\omega ^{})`$ or $`K(𝔡^{},\theta ^{})`$, where $`H^{}=U(𝔡^{})`$ and $`𝔡^{}`$ is a subalgebra of $`𝔡`$.
A Lie pseudoalgebra $`L`$ is called *semisimple* if it contains no nonzero abelian ideals. One also defines in the usual way the derived pseudoalgebra, solvable and nilpotent pseudoalgebras, and for a finite Lie pseudoalgebra $`L`$ one has the solvable radical $`\mathrm{Rad}L`$ (so that $`L/\mathrm{Rad}L`$ is semisimple).
Our Theorem 13.15 states that any finite semisimple Lie $`U(𝔡)`$-pseudoalgebra is a direct sum of finite simple Lie pseudoalgebras and of Lie pseudoalgebras of the form $`A\mathrm{Cur}𝔤`$, where $`A`$ is a subalgebra of $`W(𝔡)`$ and $`𝔤`$ is a simple finite-dimensional Lie algebra. In addition, in Theorem 13.18 we show that any subalgebra of $`W(𝔡)`$ is simple, and in Corollary 13.27 we give a complete list of all these subalgebras. (A more concise formulation of Theorem 13.10 is that any finite simple Lie pseudoalgebra over $`U(𝔡)`$ is either a current pseudoalgebra $`\mathrm{Cur}𝔤`$ over a simple finite-dimensional Lie algebra $`𝔤`$, or a nonzero subalgebra of $`W(𝔡)`$.)
Note, however, that Levi’s theorem on $`L`$ being a semidirect sum of $`L/\mathrm{Rad}L`$ and $`\mathrm{Rad}L`$ is not true even in the case $`dim𝔡=1`$. This stems from the fact that the cohomology of simple Lie pseudoalgebras with nontrivial coefficients is (highly) nontrivial (see Section 15 and \[BKV\]), in a sharp contrast with the Lie algebra case. For example, it follows from \[BKV\] that there are precisely five cases (three isolated examples and two families) of non-split extensions of $`\mathrm{Vir}`$ by $`\mathrm{Cur}`$. Translated into the language of differential Lie algebras, this result goes back to Ritt \[R3\].
Closely related to the present paper are the papers \[Ki\] and \[NW\], where (in our terminology) the annihilation algebras of rank $`1`$ over $`H`$ Lie pseudoalgebras, and of simple Lie pseudoalgebras of arbitrary finite rank, respectively, are studied. In fact, our Theorems 13.10 and 13.15 provide a completed form of the classification results of \[NW\] (in the “constant coefficients” case).
The structural results of the present paper in the simplest case $`dim𝔡=1`$ reproduce the results of \[DK\]. However, this case is much easier than the case $`dim𝔡>1`$, mainly due to the fact that only in this case is any finite torsionless $`H`$-module free.
Note also the close connection of our work to Hamiltonian formalism in the theory of nonlinear evolution equations (see the review \[DN2\], the book \[Do\] and references there, and also \[GD\], \[DN1\], \[Z\], \[M\], \[X\], and many other papers). In Section 16 we derive, as a corollary of Theorems 13.10 and 13.15, a classification of simple and semisimple linear Poisson brackets in any finite number of indeterminates.
In Section 14 we develop a representation theory of finite Lie pseudoalgebras. First, we prove an analogue of Lie’s Lemma that any weight space for an ideal of a Lie pseudoalgebra $`L`$ acting on a finite module is an $`L`$-submodule (Proposition 14.2). This implies an analogue of Lie’s Theorem that a solvable Lie pseudoalgebra has an eigenvector in any finite module (Theorem 14.3), and an analogue of Cartan–Jacobson Theorem that describes all finite Lie pseudoalgebras which have a finite faithful irreducible module (Theorem 14.5). Finally, we reduce the classification and construction of finite irreducible modules over semisimple Lie pseudoalgebras to that of irreducible modules over linearly compact Lie algebras of the type studied by Rudakov \[Ru1, Ru2\] (the complete classification will appear in a future publication). Note that complete reducibility fails already in the simplest case of Lie pseudoalgebras with $`dim𝔡=1`$ \[CKW\].
In Section 15 we define cohomology of Lie pseudoalgebras and show that it describes module extensions, abelian pseudoalgebra extensions, and pseudoalgebra deformations. We also relate this cohomology to the Gelfand–Fuchs cohomology \[Fu\]. These results generalize those of \[BKV\] in the $`dim𝔡=1`$ case.
Note that in the case $`dim𝔡=1`$ Lie pseudoalgebras are closely related to vertex algebras in a way similar to the relation of Lie algebras to universal enveloping algebras \[K2\]. We expect that, under certain conditions, there is a similar relation of “multi-dimensional” Lie pseudoalgebras to “multi-dimensional” vertex algebras defined in \[Bo2\]. In the case of a commutative Lie algebra $`𝔡`$ the Lie pseudoalgebras encode the OPE between ultralocal fields (as well as the linear Poisson brackets). However, it is not clear how Lie pseudoalgebras are related to the OPE of realistic quantum field theories.
In order to end the introduction on a more optimistic note, we would like to point out that in the definition of a Lie pseudoalgebra one may replace the permutation $`\sigma `$ by the map $`fg(gf)R`$ where $`R`$ is an R-matrix for $`H`$, hence one can take for $`H`$ any quasi-triangular Hopf algebra (defined in \[D\]). This observation, the appearance of the classical Yang–Baxter equation, and the fact that the annihilation algebra of the associative pseudoalgebra $`\mathrm{Cend}H`$ is the Drinfeld double of $`H`$, lead us to believe that there should be a deep connection between the theories of pseudoalgebras and quantum groups.
Unless otherwise specified, all vector spaces, linear maps and tensor products are considered over an algebraically closed field $`𝐤`$ of characteristic $`0`$.
### Acknowledgements
A major part of the present work was done in the fall of 1998 while two of the authors were visitors at ENS and in the spring of 1999 while they were visitors at IHES. In the spring of 2000 one of the authors was visiting MIT. The work has been completed in July 2000 while two of the authors were visiting ESI. We are grateful to these institutions for their hospitality. One of the authors wishes to thank S. P. Novikov for valuable discussions and references, and E. B. Vinberg for valuable correspondence. Finally, we thank the referees for many remarks which led to improvement of the exposition.
## 2. Preliminaries on Hopf Algebras
The goal of this section is to gather some facts and notation which will be used throughout the paper. The material in Sections 2.1 and 2.2 is standard and can be found, for example, in Sweedler’s book \[Sw\]. The material in Section 2.3 seems new.
### 2.1. Notation and basic identities
Let $`H`$ be a Hopf algebra with a coproduct $`\mathrm{\Delta }`$, a counit $`\epsilon `$, and an antipode $`S`$. We will use the following notation (cf. \[Sw\]):
(2.1) $`\mathrm{\Delta }(h)`$ $`=h_{(1)}h_{(2)},`$
(2.2) $`(\mathrm{\Delta }\mathrm{id})\mathrm{\Delta }(h)`$ $`=(\mathrm{id}\mathrm{\Delta })\mathrm{\Delta }(h)=h_{(1)}h_{(2)}h_{(3)},`$
(2.3) $`(S\mathrm{id})\mathrm{\Delta }(h)`$ $`=h_{(1)}h_{(2)},\text{etc.}`$
Note that notation (2.2) uses the coassociativity of $`\mathrm{\Delta }`$. The axioms of the antipode and the counit can be written as follows:
(2.4) $`h_{(1)}h_{(2)}`$ $`=h_{(1)}h_{(2)}=\epsilon (h),`$
(2.5) $`\epsilon (h_{(1)})h_{(2)}`$ $`=h_{(1)}\epsilon (h_{(2)})=h,`$
while the fact that $`\mathrm{\Delta }`$ is a homomorphism of algebras translates as:
(2.6)
$$(fg)_{(1)}(fg)_{(2)}=f_{(1)}g_{(1)}f_{(2)}g_{(2)}.$$
Equations (2.4) and (2.5) imply the following useful relations:
(2.7)
$$h_{(1)}h_{(2)}h_{(3)}=1h=h_{(1)}h_{(2)}h_{(3)}.$$
Let $`\mathrm{G}(H)`$ be the subset of group-like elements of $`H`$, i.e., $`gH`$ such that $`\mathrm{\Delta }(g)=gg`$. Then $`\mathrm{G}(H)`$ is a group, because $`S(g)g=gS(g)=1`$ for $`g\mathrm{G}(H)`$. Let $`\mathrm{P}(H)`$ be the subspace of primitive elements of $`H`$, i.e., $`pH`$ such that $`\mathrm{\Delta }(p)=p1+1p`$. This is a Lie algebra with respect to the commutator $`[p,q]=pqqp`$. Note that $`\mathrm{G}(H)`$ acts on $`\mathrm{P}(H)`$ by inner automorphisms: $`gpg^1\mathrm{P}(H)`$ for $`p\mathrm{P}(H)`$, $`g\mathrm{G}(H)`$.
The proof of the following theorem may be found in \[Sw\].
###### Theorem 2.1 (Kostant).
Let $`H`$ be a cocommutative Hopf algebra over $`𝐤`$ $`(`$an algebraically closed field of characteristic $`0)`$. Then $`H`$ is isomorphic $`(`$as a Hopf algebra$`)`$ to the smash product of the universal enveloping algebra $`U(\mathrm{P}(H))`$ and the group algebra $`𝐤[\mathrm{G}(H)]`$.
An associative algebra $`A`$ is called an $`H`$-differential algebra if it is also a left $`H`$-module such that the multiplication $`AAA`$ is a homomorphism of $`H`$-modules. In other words,
(2.8)
$$h(xy)=(h_{(1)}x)(h_{(2)}y)$$
for $`hH`$, $`x,yA`$. The smash product $`A\mathrm{}H`$ of an $`H`$-differential algebra $`A`$ with $`H`$ is the tensor product $`AH`$ of vector spaces but with a new multiplication:
(2.9)
$$(a\mathrm{}g)(b\mathrm{}h)=a(g_{(1)}b)\mathrm{}g_{(2)}h.$$
If both $`A`$ and $`H`$ are Hopf algebras, then $`A\mathrm{}H`$ is a Hopf algebra if we consider it as a tensor product of coalgebras. In the theorem above, $`U(\mathrm{P}(H))`$ is a $`𝐤[\mathrm{G}(H)]`$-differential algebra with respect to the adjoint action of $`\mathrm{G}(H)`$ on $`\mathrm{P}(H)`$.
It is worth mentioning that as a byproduct of Kostant’s Theorem, we obtain that the antipode of a cocommutative Hopf algebra is an involution, i.e., $`S^2=\mathrm{id}`$.
We will often be working with the Hopf algebra $`H=U(𝔡)`$, where $`𝔡`$ is a finite-dimensional Lie algebra. It is well known that this is a Noetherian domain, and any two nonzero elements $`f,gH`$ have a nonzero left (respectively right) common multiple. In particular, $`H=U(𝔡)`$ has a skew-field of fractions $`K`$.
###### Lemma 2.2.
Let $`H`$ be a Noetherian domain which has a skew-field of fractions $`K`$, and let $`L`$ be a finite $`H`$-module. Then there is a homomorphism $`i:LF`$ from $`L`$ to a free $`H`$-module $`F`$, whose kernel is the torsion submodule of $`L`$. If $`L`$ is torsion-free, then the module $`F`$ can be chosen in such a way that $`hFi(L)`$ for some nonzero $`hH`$ and $`i(L)/hF`$ is torsion.
###### Proof.
The kernel of the natural map $`\iota :LL_K:=K_HL`$ is the torsion of $`L`$. The image of $`L`$ under this map is contained inside a free $`H`$-submodule of $`L_K`$. In order to see this, let us consider a set of $`H`$-generators $`\{l_1,\mathrm{},l_n\}`$ of $`L`$, and a $`K`$-basis $`\{v_1,\mathrm{},v_k\}`$ of $`L_K`$. We can express the elements $`\iota (l_j)`$ as $`K`$-linear combinations of the $`v_i`$’s, and by rescaling elements of this basis by a common multiple of the denominators, we can assume the $`\iota (l_j)`$’s to be $`H`$-linear combinations of the $`v_i`$’s. Hence the image $`\iota (L)`$ is contained in the $`H`$-module $`F`$ spanned by the $`v_i`$’s, which is free by construction.
The fact that $`F/L`$ is torsion is clear because there exist nonzero elements $`h_iH`$ such that $`h_iv_iL`$. If $`h`$ is a common multiple of the $`h_i`$’s, then $`hF`$ is contained in $`L`$. On the other hand, the inclusion $`LF`$ implies $`hLhF`$, hence $`h(L/hF)=0`$ and $`L/hF`$ is torsion. ∎
### 2.2. Filtration and topology
We define an increasing sequence of subspaces of a Hopf algebra $`H`$ inductively by:
(2.10) $`\mathrm{F}^nH=0\text{for }n<0\text{,}\mathrm{F}^0H=𝐤[\mathrm{G}(H)],`$
(2.11) $`\begin{array}{cc}\hfill \mathrm{F}^nH=\mathrm{span}_𝐤\{hH|\mathrm{\Delta }(h)\mathrm{F}^0Hh& +h\mathrm{F}^0H\hfill \\ & +_{i=1}^{n1}\mathrm{F}^iH\mathrm{F}^{ni}H\}.\hfill \end{array}`$
It has the following properties (which are immediate from definitions):
(2.12) $`(\mathrm{F}^mH)(\mathrm{F}^nH)`$ $`\mathrm{F}^{m+n}H,`$
(2.13) $`\mathrm{\Delta }(\mathrm{F}^nH)`$ $`_{i=0}^n\mathrm{F}^iH\mathrm{F}^{ni}H,`$
(2.14) $`S(\mathrm{F}^nH)`$ $`\mathrm{F}^nH.`$
When $`H`$ is cocommutative, using Theorem 2.1, one can show that:
(2.15)
$$\underset{n}{}\mathrm{F}^nH=H.$$
(This condition is also satisfied when $`H`$ is a quantum universal enveloping algebra.) Provided that (2.15) holds, we say that a nonzero element $`aH`$ has degree $`n`$ if $`a\mathrm{F}^nH\mathrm{F}^{n1}H`$.
When $`H`$ is a universal enveloping algebra, we get its canonical filtration. Later in some instances we will also impose the following finiteness condition on $`H`$:
(2.16)
$$dim\mathrm{F}^nH<\mathrm{}n.$$
It is satisfied when $`H`$ is a universal enveloping algebra of a finite-dimensional Lie algebra, or its smash product with the group algebra of a finite group.
Now let $`X=H^{}:=\mathrm{Hom}_𝐤(H,𝐤)`$ be the dual of $`H`$. Recall that $`H`$ acts on $`X`$ by the formula ($`h,fH`$, $`xX`$):
(2.17)
$$hx,f=x,S(h)f,$$
so that $`X`$ is an associative $`H`$-differential algebra (see (2.8)). Moreover, $`X`$ is commutative when $`H`$ is cocommutative. Similarly, one can define a right action of $`H`$ on $`X`$ by
(2.18)
$$xh,f=x,fS(h),$$
and then we have
(2.19)
$$(xy)h=(xh_{(1)})(yh_{(2)}).$$
Associativity of $`H`$ implies that $`X`$ is an $`H`$-bimodule, i.e.
(2.20)
$$f(xg)=(fx)g,f,gH,xX.$$
Let $`X=\mathrm{F}_1X\mathrm{F}_0X\mathrm{}`$ be the decreasing sequence of subspaces of $`X`$ dual to $`\mathrm{F}^nH`$: $`\mathrm{F}_nX=(\mathrm{F}^nH)^{}`$. It has the following properties:
(2.21) $`(\mathrm{F}_mX)(\mathrm{F}_nX)`$ $`\mathrm{F}_{m+n}X,`$
(2.22) $`(\mathrm{F}^mH)(\mathrm{F}_nX)`$ $`\mathrm{F}_{nm}X,`$
and
(2.23)
$$\underset{n}{}\mathrm{F}_nX=0,\text{provided that (}\text{2.15}\text{) holds}.$$
We define a topology of $`X`$ by considering $`\{\mathrm{F}_nX\}`$ as a fundamental system of neighborhoods of $`0`$. We will always consider $`X`$ with this topology, while $`H`$ with the discrete topology. It follows from (2.23) that $`X`$ is Hausdorff, provided that (2.15) holds. By (2.21) and (2.22), the multiplication of $`X`$ and the action of $`H`$ on it are continuous; in other words, $`X`$ is a topological $`H`$-differential algebra.
We define an antipode $`S:XX`$ as the dual of that of $`H`$:
(2.24)
$$S(x),h=x,S(h).$$
Then we have:
(2.25)
$$S(ab)=S(b)S(a)\text{for}a,bX\text{or}H.$$
We will also define a comultiplication $`\mathrm{\Delta }:XX\widehat{}X`$ as the dual of the multiplication $`HHH`$, where $`X\widehat{}X:=(HH)^{}`$ is the completed tensor product. Formally, we will use the same notation for $`X`$ as for $`H`$ (see (2.1)–(2.3)), writing for example $`\mathrm{\Delta }(x)=x_{(1)}x_{(2)}`$ for $`xX`$. By definition, for $`x,yX`$, $`f,gH`$, we have:
(2.26) $`xy,f`$ $`=xy,\mathrm{\Delta }(f)=x,f_{(1)}y,f_{(2)},`$
(2.27) $`x,fg`$ $`=\mathrm{\Delta }(x),fg=x_{(1)},fx_{(2)},g.`$
We have:
(2.28) $`S(\mathrm{F}_nX)`$ $`\mathrm{F}_nX,`$
(2.29) $`\mathrm{\Delta }(\mathrm{F}_nX)`$ $`_{i=0}^n\mathrm{F}_iX\widehat{}\mathrm{F}_{ni}X.`$
If $`H`$ satisfies the finiteness condition (2.16), then the filtration of $`X`$ satisfies
(2.30)
$$dimX/\mathrm{F}_nX<\mathrm{}n,$$
which implies that $`X`$ is linearly compact (see Section 6 below).
By a basis of $`X`$ we will always mean a topological basis $`\{x_i\}`$ which tends to $`0`$, i.e., such that for any $`n`$ all but a finite number of $`x_i`$ belong to $`\mathrm{F}_nX`$. Let $`\{h_i\}`$ be a basis of $`H`$ (as a vector space) compatible with the filtration. Then the set of elements $`\{x_i\}`$ of $`X`$ defined by $`x_i,h_j=\delta _{ij}`$ is called the dual basis of $`X`$. If $`H`$ satisfies (2.16), then $`\{x_i\}`$ is a basis of $`X`$ in the above sense, i.e., it tends to $`0`$. We have for $`gH`$, $`yX`$:
$$g=_ig,x_ih_i,y=_iy,h_ix_i,$$
where the first sum is finite, and the second one is convergent in $`X`$.
###### Example 2.3.
Let $`H=U(𝔡)`$ be the universal enveloping algebra of an $`N`$-dimensional Lie algebra $`𝔡`$. Fix a basis $`\{_i\}`$ of $`𝔡`$, and for $`I=(i_1,\mathrm{},i_N)_+^N`$ let $`^{(I)}=_1^{i_1}\mathrm{}_N^{i_N}/i_1!\mathrm{}i_N!`$. Then $`\{^{(I)}\}`$ is a basis of $`H`$ (the Poincaré–Birkhoff–Witt basis). Moreover, it is easy to see that
(2.31)
$$\mathrm{\Delta }(^{(I)})=\underset{J+K=I}{}^{(J)}^{(K)}.$$
If $`\{t_I\}`$ is the dual basis of $`X`$, defined by $`t_I,^{(J)}=\delta _{I,J}`$, then (2.31) implies $`t_Jt_K=t_{J+K}`$. Therefore, $`X`$ can be identified with the ring $`𝒪_N=𝐤[[t_1,\mathrm{},t_N]]`$ of formal power series in $`N`$ indeterminates. Then the action of $`H`$ on $`𝒪_N`$ is given by differential operators.
###### Lemma 2.4.
If $`\{h_i\}`$, $`\{x_i\}`$ are dual bases in $`H`$ and $`X`$, then
(2.32)
$$\mathrm{\Delta }(x)=_ixS(h_i)x_i=_ix_iS(h_i)x$$
for any $`xX`$.
###### Proof.
For $`f,gH`$, we have:
$`_ixS(h_i)x_i,fg`$ $`=_ixS(h_i),fx_i,g`$
$`=xS(g),f`$ $`=x,fg=\mathrm{\Delta }(x),fg,`$
which proves the first identity. The second one is proved in the same way. ∎
### 2.3. Fourier transform
For an arbitrary Hopf algebra $`H`$, we introduce a map $`:HHHH`$, called the Fourier transform, by the formula
(2.33)
$$(fg)=(f1)(S\mathrm{id})\mathrm{\Delta }(g)=fg_{(1)}g_{(2)}.$$
It follows from (2.7) that $``$ is a vector space isomorphism with an inverse given by
(2.34)
$$^1(fg)=(f1)\mathrm{\Delta }(g)=fg_{(1)}g_{(2)}.$$
Indeed, using the coassociativity of $`\mathrm{\Delta }`$ and (2.7), we compute
$$^1(fg_{(1)}g_{(2)})=fg_{(1)}(g_{(2)})_{(1)}(g_{(2)})_{(2)}=fg_{(1)}g_{(2)}g_{(3)}=fg.$$
The significance of $``$ is in the identity
(2.35)
$$fg=^1(fg)=(fg_{(1)}1)\mathrm{\Delta }(g_{(2)}),$$
which, together with properties (2.12)–(2.14) of the filtration of $`H`$, implies the next result.
###### Lemma 2.5.
(i) Every element of $`HH`$ can be uniquely represented in the form $`_i(h_i1)\mathrm{\Delta }(l_i)`$, where $`\{h_i\}`$ is a fixed $`𝐤`$-basis of $`H`$ and $`l_iH`$. In other words, $`HH=(H𝐤)\mathrm{\Delta }(H)`$.
(ii) We have:
(2.36)
$$(\mathrm{F}^nH𝐤)\mathrm{\Delta }(H)=\mathrm{F}^n(HH)\mathrm{\Delta }(H)=(𝐤\mathrm{F}^nH)\mathrm{\Delta }(H),$$
where $`\mathrm{F}^n(HH)=_{i+j=n}\mathrm{F}^iH\mathrm{F}^jH`$.
In particular, for any $`H`$-module $`W`$, we have:
(2.37)
$$(\mathrm{F}^nH𝐤)_HW=\mathrm{F}^n(HH)_HW=(𝐤\mathrm{F}^nH)_HW.$$
###### Proof.
For $`hHH`$ we have:
$$h=_i(h_i1)\mathrm{\Delta }(l_i)=^1(_ih_il_i)\text{iff}_ih_il_i=(h).$$
This proves (i).
To prove (2.36), it is enough to show that $`\mathrm{F}^n(HH)(\mathrm{F}^nH𝐤)\mathrm{\Delta }(H)`$. This follows from the above equation and the fact that $`(\mathrm{F}^n(HH))\mathrm{F}^n(HH)\mathrm{F}^nHH`$. ∎
The Fourier transform $``$ has the following properties (which are easy to check using (2.4)–(2.6)):
(2.38) $`\left((fg)\mathrm{\Delta }(h)\right)`$ $`=(fg)(1h),`$
(2.39) $`(hfg)`$ $`=(h1)(fg),`$
(2.40) $`(fhg)`$ $`=(1h_{(2)})(fg)(h_{(1)}1),`$
(2.41) $`_{12}_{13}_{23}`$ $`=_{23}_{12}.`$
Here in (2.41), we use the standard notation $`_{12}=\mathrm{id}`$ acting on $`HHH`$.
## 3. Pseudotensor Categories and Pseudoalgebras
In this section, we review some definitions of Beilinson and Drinfeld \[BD\]; we also use the exposition in \[BKV, Section 12\].
The theory of conformal algebras \[K2\] is in many ways analogous to the theory of Lie algebras. The reason is that in fact conformal algebras can be considered as Lie algebras in a certain “pseudotensor” category, instead of the category of vector spaces. A pseudotensor category \[BD\] is a category equipped with “polylinear maps” and a way to compose them (such categories were first introduced by Lambek \[L\] under the name multicategories). This is enough to define the notions of Lie algebra, representations, cohomology, etc.
As an example, consider first the category $`𝒱ec`$ of vector spaces (over $`𝐤`$). For a finite nonempty set $`I`$ and a collection of vector spaces $`\{L_i\}_{iI}`$, $`M`$, we can define the space of polylinear maps from $`\{L_i\}_{iI}`$ to $`M`$ as
$$\mathrm{Lin}(\{L_i\}_{iI},M)=\mathrm{Hom}(_{iI}L_i,M).$$
The symmetric group $`S_I`$ acts among these spaces by permuting the factors in $`_{iI}L_i`$.
For any surjection of finite sets $`\pi :JI`$ and a collection $`\{N_j\}_{jJ}`$, we have the obvious compositions of polylinear maps
(3.1) $`\mathrm{Lin}(\{L_i\}_{iI},M){\displaystyle \underset{iI}{}}\mathrm{Lin}(\{N_j\}_{jJ_i},L_i)\mathrm{Lin}(\{N_j\}_{jJ},M),`$
(3.2) $`\varphi \times \{\psi _i\}_{iI}\varphi (_{iI}\psi _i)\varphi (\{\psi _i\}_{iI}),`$
where $`J_i=\pi ^1(i)`$ for $`iI`$.
The compositions have the following properties:
If $`KJ`$, $`\{P_k\}_{kK}`$ is a family of objects and $`\chi _j\mathrm{Lin}(\{P_k\}_{kK_j},N_j)`$, then $`\varphi \left(\left\{\psi _i(\{\chi _j\}_{jJ_i})\right\}_{iI}\right)=\left(\varphi (\{\psi _i\}_{iI})\right)(\{\chi _j\}_{jJ})\mathrm{Lin}(\{P_k\}_{kK},M)`$.
For any object $`M`$ there is an element $`\mathrm{id}_M\mathrm{Lin}(\{M\},M)`$ such that for any $`\varphi \mathrm{Lin}(\{L_i\}_{iI},M)`$ one has $`\mathrm{id}_M(\varphi )=\varphi (\{\mathrm{id}_{L_i}\}_{iI})=\varphi `$.
The compositions (3.1) are equivariant with respect to the natural action of the symmetric group.
###### Definition 3.1 (\[BD\]).
A pseudotensor category is a class of objects $``$ together with vector spaces $`\mathrm{Lin}(\{L_i\}_{iI},M)`$, equipped with actions of the symmetric groups $`S_I`$ among them and composition maps (3.1), satisfying the above three properties.
###### Remark 3.2.
For a pseudotensor category $``$ and objects $`L,M`$, let $`\mathrm{Hom}(L,M)=\mathrm{Lin}(\{L\},M)`$. This gives a structure of an ordinary (additive) category on $``$ and all $`\mathrm{Lin}`$ are functors $`(^{})^I\times 𝒱ec`$, where $`^{}`$ is the dual category of $``$.
###### Remark 3.3.
The notion of pseudotensor category is a straightforward generalization of the notion of operad. By definition, an operad is a pseudotensor category with only one object.
###### Definition 3.4.
A Lie algebra in a pseudotensor category $``$ is an object $`L`$ equipped with $`\beta \mathrm{Lin}(\{L,L\},L)`$ satisfying the following properties.
$`\beta =\sigma _{12}\beta ,`$ where $`\sigma _{12}=(12)S_2`$.
$`\beta (\beta (,),)=\beta (,\beta (,))\sigma _{12}\beta (,\beta (,)),`$ where now $`\sigma _{12}=(12)`$ is viewed as an element of $`S_3`$.
It is instructive to think of a polylinear map $`\varphi \mathrm{Lin}(\{L_i\}_{i=1}^n,M)`$ as an operation with $`n`$ inputs and $`1`$ output, as depicted in Figure 1. The skew-commutativity and Jacobi identity for a Lie algebra $`(L,\beta )`$ are represented pictorially in Figures 2 and 3.
###### Definition 3.5.
A representation of a Lie algebra $`(L,\beta )`$ is an object $`M`$ together with $`\rho \mathrm{Lin}(\{L,M\},M)`$ satisfying
$$\rho (\beta (,),)=\rho (,\rho (,))\sigma _{12}\rho (,\rho (,)).$$
Similarly, one can define cohomology of a Lie algebra $`(L,\beta )`$ with coefficients in a module $`(M,\rho )`$ (cf. \[BKV\]).
###### Definition 3.6.
An $`n`$-cochain of a Lie algebra $`(L,\beta )`$, with coefficients in a module $`(M,\rho )`$ over it, is a polylinear operation $`\gamma \mathrm{Lin}(\{\underset{n}{\underset{}{L,\mathrm{},L}}\},M)`$ which is skew-symmetric, i.e., satisfying for all $`i=1,\mathrm{},n1`$ the identity shown in Figure 4. The differential $`d\gamma `$ of a cochain $`\gamma `$ is defined by Figure 5. The same computation as in the ordinary Lie algebra case shows that $`d^2=0`$. The cohomology of the resulting complex is called the cohomology of $`L`$ with coefficients in $`M`$ and is denoted by $`\mathrm{H}^{}(L,M)`$.
###### Example 3.7.
A Lie algebra in the category of vector spaces $`𝒱ec`$ is just an ordinary Lie algebra. The same is true for representations and cohomology.
###### Example 3.8.
Let $`H`$ be a cocommutative bialgebra with a comultiplication $`\mathrm{\Delta }`$ and a counit $`\epsilon `$. Then the category $`^l(H)`$ of left $`H`$-modules is a symmetric tensor category. Hence, $`^l(H)`$ is a pseudotensor category with polylinear maps
(3.3)
$$\mathrm{Lin}(\{L_i\}_{iI},M)=\mathrm{Hom}_H(_{iI}L_i,M).$$
The composition of polylinear maps is given by (3.2). An algebra (e.g., Lie or associative) in the category $`^l(H)`$ will be called an $`H`$-differential algebra: this is an ordinary algebra which is also a left $`H`$-module and such that the product (or the bracket) is a homomorphism of $`H`$-modules, see (2.8).
One can also define the notions of associative algebra or commutative algebra in a pseudotensor category, their representations and analogues of the Hochschild, cyclic, or Harrison cohomology.
###### Definition 3.9.
An associative algebra in a pseudotensor category $``$ is an object $`A`$ and a product $`\mu \mathrm{Lin}(\{A,A\},A)`$ satisfying
$`\mu (\mu (,),)=\mu (,\mu (,))`$,
see Figure 6. The algebra $`(A,\mu )`$ is called commutative if, in addition, $`\mu `$ satisfies
$`\mu =\sigma _{12}\mu ,`$ where $`\sigma _{12}=(12)S_2`$.
###### Remark 3.10.
In order to define the notion of an associative algebra in a pseudotensor category, one does not use the actions of the symmetric groups among the spaces of polylinear maps. One can relax the definition of a pseudotensor category by forgetting these actions. Then what we call a “pseudotensor category” should be termed a “symmetric pseudotensor category”, while there is a more general notion of a “braided” one (cf. \[So\]).
###### Proposition 3.11.
Let $`(A,\mu )`$ be an associative algebra in a pseudotensor category $``$. Define $`\beta \mathrm{Lin}(\{A,A\},A)`$ as the commutator $`\beta :=\mu \sigma _{12}\mu `$, see Figure 7. Then $`(A,\beta )`$ is a Lie algebra in $``$.
Let $`H`$ be a cocommutative bialgebra with a comultiplication $`\mathrm{\Delta }`$. We introduce a pseudotensor category $`^{}(H)`$ with the same objects as $`^l(H)`$ (i.e., left $`H`$-modules) but with another pseudotensor structure \[BD\]:
(3.4)
$$\mathrm{Lin}(\{L_i\}_{iI},M)=\mathrm{Hom}_{H^I}(_{iI}L_i,H^I_HM).$$
Here $`_{iI}`$ is the tensor product functor $`^l(H)^I^l(H^I)`$. For a surjection $`\pi :JI`$, the composition of polylinear maps is defined as follows:
(3.5)
$$\varphi \left(\{\psi _i\}_{iI}\right)=\mathrm{\Delta }^{(\pi )}\left(\varphi \right)\left(_{iI}\psi _i\right).$$
Here $`\mathrm{\Delta }^{(\pi )}`$ is the functor $`^l(H^I)^l(H^J)`$, $`MH^J_{H^I}M`$, where $`H^I`$ acts on $`H^J`$ via the iterated comultiplication determined by $`\pi `$.
Explicitly, let $`n_jN_j`$ $`(jJ)`$, and write
(3.6)
$$\psi _i\left(_{jJ_i}n_j\right)=\underset{r}{}g_i^r_Hl_i^r,g_i^rH^{J_i},l_i^rL_i,$$
where, as before, $`J_i=\pi ^1(i)`$ for $`iI`$. Let
(3.7)
$$\varphi \left(_{iI}l_i^r\right)=\underset{s}{}f^{rs}_Hm^{rs},f^{rs}H^I,m^{rs}M.$$
Then, by definition,
(3.8)
$$\left(\varphi \left(\{\psi _i\}_{iI}\right)\right)\left(_{jJ}n_j\right)=\underset{r,s}{}(_{iI}g_i^r)\mathrm{\Delta }^{(\pi )}(f^{rs})_Hm^{rs},$$
where $`\mathrm{\Delta }^{(\pi )}:H^IH^J`$ is the iterated comultiplication determined by $`\pi `$. For example, if $`\pi :\{1,2,3\}\{1,2\}`$ is given by $`\pi (1)=\pi (2)=1`$, $`\pi (3)=2`$, then $`\mathrm{\Delta }^{(\pi )}=\mathrm{\Delta }\mathrm{id}`$; if $`\pi (1)=1`$, $`\pi (2)=\pi (3)=2`$, then $`\mathrm{\Delta }^{(\pi )}=\mathrm{id}\mathrm{\Delta }`$.
The symmetric group $`S_I`$ acts among the spaces $`\mathrm{Lin}(\{L_i\}_{iI},M)`$ by simultaneously permuting the factors in $`_{iI}L_i`$ and $`H^I`$. This is the only place where we need the cocommutativity of $`H`$; for example, the permutation $`\sigma _{12}=(12)S_2`$ acts on $`(HH)_HM`$ by
$$\sigma _{12}\left((fg)_Hm\right)=(gf)_Hm,$$
and this is well defined only when $`H`$ is cocommutative.
One can generalize the above construction for (quasi)triangular bialgebras as follows.
###### Remark 3.12.
Let $`H`$ be a triangular bialgebra with a universal R-matrix $`R`$. Recall that $`R`$ is an invertible element of $`HH`$ satisfying the following equations:
(3.9) $`\sigma (R)`$ $`=R^1,`$
(3.10) $`\sigma (\mathrm{\Delta }(h))R`$ $`=R\mathrm{\Delta }(h)hH,`$
(3.11) $`(\mathrm{id}\mathrm{\Delta })R`$ $`=R_{13}R_{12},`$
(3.12) $`(\mathrm{\Delta }\mathrm{id})R`$ $`=R_{13}R_{23},`$
where $`\sigma `$ is the permutation $`\sigma (fg)=gf`$, and we use the standard notation $`R_{12}=R\mathrm{id}HHH`$, etc. Then we define a pseudotensor category $`^{}(H)`$ as above but with a modified action of the symmetric groups. It is enough to describe the action of the transposition $`\sigma _{12}=(12)S_2`$ on $`(HH)_HM`$; it is given by
$$\sigma _{12}\left((fg)_Hm\right)=(gf)R_Hm.$$
This is well defined because of (3.10), and $`\sigma _{12}^2=\mathrm{id}`$ because of (3.9). Since any permutation is a product of transpositions, this can be extended to an action of the symmetric group among the spaces of polylinear maps; due to (3.11), (3.12), this action is compatible with compositions.
If $`H`$ is quasitriangular, i.e., if we drop relation (3.9), we will get an action of the braid group instead of the symmetric one and a “braided” pseudotensor category (cf. Remark 3.10).
The following notion will be the main object of our study.
###### Definition 3.13.
A Lie $`H`$-pseudoalgebra (or just a Lie pseudoalgebra) is a Lie algebra $`(L,\beta )`$ in the pseudotensor category $`^{}(H)`$ as defined above.
Examples of Lie pseudoalgebras will be given in Sections 4 and 8 below. One can also define associative $`H`$-pseudoalgebras as associative algebras $`(A,\mu )`$ in the pseudotensor category $`^{}(H)`$. It is convenient to define the general notion of an algebra in $`^{}(H)`$ as follows.
###### Definition 3.14.
An $`H`$-pseudoalgebra (or just a pseudoalgebra) is a left $`H`$-module $`A`$ together with an operation $`\mu \mathrm{Hom}_{HH}(AA,(HH)_HA)`$, called the pseudoproduct.
We will denote the pseudoproduct $`\mu (ab)(HH)_HA`$ of two elements $`a,bA`$ by $`ab`$. It has the following defining property:
For $`a,bA`$, $`f,gH`$, one has
(3.13)
$$fagb=((fg)_H1)(ab).$$
Explicitly, if
(3.14)
$$ab=_i(f_ig_i)_He_i,$$
then $`fagb=_i(ff_igg_i)_He_i`$.
To describe explicitly the associativity condition for a pseudoproduct $`\mu `$, we need to compute the compositions $`\mu (\mu (,),)`$ and $`\mu (,\mu (,))`$ in $`^{}(H)`$. Let $`ab`$ be given by (3.14), and let
(3.15)
$$e_ic=_{i,j}(f_{ij}g_{ij})_He_{ij}.$$
Then $`(ab)c\mu (\mu (ab)c)`$ is the following element of $`H^3_HA`$ (cf. (3.8)):
(3.16)
$$(ab)c=_{i,j}(f_if_{ij}^{}{}_{(1)}{}^{}g_if_{ij}^{}{}_{(2)}{}^{}g_{ij})_He_{ij}.$$
Similarly, if we write
(3.17) $`bc`$ $`=_i(h_il_i)_Hd_i,`$
(3.18) $`ad_i`$ $`=_{i,j}(h_{ij}l_{ij})_Hd_{ij},`$
then
(3.19) $`a(bc)`$ $`=_{i,j}(h_{ij}h_il_{ij}^{}{}_{(1)}{}^{}l_il_{ij}^{}{}_{(2)}{}^{})_Hd_{ij}.`$
Now a pseudoproduct $`ab`$ is associative iff it satisfies
(3.20)
$$a(bc)=(ab)c$$
in $`H^3_HA`$, where the compositions $`(ab)c`$ and $`a(bc)`$ are given by the above formulas.
The pseudoproduct $`ab`$ is commutative iff it satisfies
(3.21)
$$ba=(\sigma _H\mathrm{id})(ab),$$
where $`\sigma :HHHH`$ is the permutation $`\sigma (fg)=gf`$. Explicitly,
(3.22)
$$ba=_i(g_if_i)_He_i,$$
if $`ab`$ is given by (3.14). Note that the right-hand side of (3.21) is well defined due to the cocommutativity of $`H`$.
In the case of a Lie pseudoalgebra $`(L,\beta )`$, we will call the pseudoproduct $`\beta `$ a pseudobracket, and we will denote it by $`[ab]`$. Let us spell out its properties ($`a,b,cL`$, $`f,gH`$):
(3.23)
$$[fagb]=((fg)_H1)[ab].$$
(3.24)
$$[ba]=(\sigma _H\mathrm{id})[ab].$$
(3.25)
$$[a[bc]]((\sigma \mathrm{id})_H\mathrm{id})[b[ac]]=[[ab]c]$$
in $`H^3_HL`$, where the compositions $`[[ab]c]`$ and $`[a[bc]]`$ are defined as above.
###### Proposition 3.15.
Let $`(A,\mu )`$ be an associative $`H`$-pseudoalgebra. Define a pseudobracket $`\beta `$ as the commutator $`[ab]=ab(\sigma _H\mathrm{id})(ba)`$. Then $`(A,\beta )`$ is a Lie $`H`$-pseudoalgebra (cf. Proposition 3.11).
The definitions of representations of Lie pseudoalgebras or associative pseudoalgebras are obvious modifications of the above.
###### Definition 3.16.
A representation of an associative $`H`$-pseudoalgebra $`A`$ is a left $`H`$-module $`M`$ together with an operation $`\rho \mathrm{Lin}(\{A,M\},M)`$, written as $`ac\rho (ac)(HH)_HM`$, which satisfies (3.20) for $`a,bA`$, $`cM`$.
###### Definition 3.17.
A representation of a Lie $`H`$-pseudoalgebra $`L`$ is a left $`H`$-module $`M`$ together with an operation $`\rho \mathrm{Lin}(\{L,M\},M)`$, written as $`ac\rho (ac)`$, which satisfies
(3.26)
$$a(bc)((\sigma \mathrm{id})_H\mathrm{id})(b(ac))=[ab]c$$
for $`a,bL`$, $`cM`$.
## 4. Some Examples of Lie Pseudoalgebras
In this section we give some examples of Lie pseudoalgebras, and discuss their relationship with previously known objects. Other important examples — the pseudoalgebras of vector fields — are treated in detail in Section 8.
### 4.1. Conformal algebras
The (Lie) conformal algebras introduced by Kac \[K2\] are exactly the (Lie) $`𝐤[]`$-pseudoalgebras, where $`𝐤[]`$ is the Hopf algebra of polynomials in one variable $``$. The explicit relation between the $`\lambda `$-bracket of \[DK\] and the pseudobracket of Section 3 is:
$$[a_\lambda b]=_ip_i(\lambda )c_i[ab]=_i(p_i()1)_{𝐤[]}c_i.$$
This correspondence has been explained in detail in the introduction.
Similarly, for $`H=𝐤[_1,\mathrm{},_N]`$ we get conformal algebras in $`N`$ indeterminates, see \[BKV, Section 10\]. We may say that for $`N=0`$, $`H`$ is $`𝐤`$; then a $`𝐤`$-conformal algebra is the same as a Lie algebra, cf. Example 3.7.
On the other hand, when $`H=𝐤[\mathrm{\Gamma }]`$ is the group algebra of a group $`\mathrm{\Gamma }`$, one obtains the $`\mathrm{\Gamma }`$-conformal algebras studied in \[GK\]. This is a special case of a more general construction described in Section 5 below.
### 4.2. Current pseudoalgebras
Let $`H^{}`$ be a Hopf subalgebra of $`H`$, and let $`A`$ be an $`H^{}`$-pseudoalgebra. Then we define the current $`H`$-pseudoalgebra $`\mathrm{Cur}_H^{}^HA\mathrm{Cur}A`$ as $`H_H^{}A`$ by extending the pseudoproduct $`ab`$ of $`A`$ using the $`H`$-bilinearity. Explicitly, for $`a,bA`$, we define
$$\begin{array}{cc}\hfill (f_H^{}a)(g_H^{}b)& =((fg)_H1)(ab)\hfill \\ & =_i(ff_igg_i)_H(1_H^{}e_i),\hfill \end{array}$$
if $`ab=_i(f_ig_i)_H^{}e_i`$. Then $`\mathrm{Cur}_H^{}^HA`$ is an $`H`$-pseudoalgebra which is Lie or associative when $`A`$ is so.
An important special case is when $`H^{}=𝐤`$: given a Lie algebra $`𝔤`$, let $`\mathrm{Cur}𝔤=H𝔤`$ with the following pseudobracket
$$[(fa)(gb)]=(fg)_H(1[a,b]).$$
Then $`\mathrm{Cur}𝔤`$ is a Lie $`H`$-pseudoalgebra.
### 4.3. $`H`$-pseudoalgebras of rank $`1`$
Let $`L=He`$ be a Lie pseudoalgebra which is a free $`H`$-module of rank $`1`$. Then, by $`H`$-bilinearity, the pseudobracket on $`L`$ is determined by $`[ee]`$, or equivalently, by an $`\alpha HH`$ such that $`[ee]=\alpha _He`$.
###### Proposition 4.1.
$`L=He`$ with the pseudobracket $`[ee]=\alpha _He`$ is a Lie $`H`$-pseudoalgebra iff $`\alpha HH`$ satisfies the following equations:
(4.1) $`\alpha =\sigma (\alpha ),`$
(4.2) $`(\alpha 1)(\mathrm{\Delta }\mathrm{id})(\alpha )=(1\alpha )(\mathrm{id}\mathrm{\Delta })(\alpha )(\sigma \mathrm{id})\left((1\alpha )(\mathrm{id}\mathrm{\Delta })(\alpha )\right).`$
Similarly, $`A=Ha`$ with a pseudoproduct $`aa=\alpha _Ha`$ is an associative $`H`$-pseudoalgebra iff $`\alpha `$ satisfies
$$(\alpha 1)(\mathrm{\Delta }\mathrm{id})(\alpha )=(1\alpha )(\mathrm{id}\mathrm{\Delta })(\alpha ).$$
###### Proof.
Follows immediately from definitions. Indeed, if $`[ee]=\alpha _He`$, then:
$`[[ee]e]`$ $`=(\alpha 1)(\mathrm{\Delta }\mathrm{id})(\alpha )_He,`$
$`[e[ee]]`$ $`=(1\alpha )(\mathrm{id}\mathrm{\Delta })(\alpha )_He.`$
###### Lemma 4.2.
Let $`H=U(𝔡)`$ be the universal enveloping algebra of a Lie algebra $`𝔡`$. Then any solution $`\alpha HH`$ of equations (4.1), (4.2) is of the form $`\alpha =r+s11s`$, where $`r𝔡𝔡`$, $`s𝔡`$.
In this case (4.2) is equivalent to the following system of equations:
(4.3) $`[r,\mathrm{\Delta }(s)]=0,`$
(4.4) $`([r_{12},r_{13}]+r_{12}s_3)+\text{cyclic}=0.`$
$`(`$As usual, $`r_{12}=r1`$, $`s_3=11s`$, etc., and “cyclic” here and further means applying the two nontrivial cyclic permutations on $`HHH.)`$
###### Proof.
Using an argument similar to that of \[Ki\], we will show that if $`\alpha `$ satisfies (4.2) then $`\alpha H(𝔡+𝐤)`$. Then (4.1) will imply the first claim, that $`\alpha (𝔡+𝐤)(𝔡+𝐤)`$.
Let $`\{_1,\mathrm{},_N\}`$ be a basis of $`𝔡`$, and let us consider the corresponding Poincaré–Birkhoff–Witt basis of $`H=U(𝔡)`$ given by elements $`^{(I)}:=_1^{i_1}\mathrm{}_N^{i_N}/i_1!\mathrm{}i_N!`$, where $`I=(i_1,\mathrm{},i_N)_+^N`$. In this basis the comultiplication takes the simple form (2.31). We can write $`\alpha =_I\alpha _I^{(I)}`$, $`\alpha _IH`$. Equation (4.2) then becomes:
(4.5)
$$\underset{I}{}\alpha \mathrm{\Delta }(\alpha _I)^{(I)}=\underset{I,J,K}{}(\alpha _{J+K}\alpha _I^{(J)}\alpha _I^{(J)}\alpha _{J+K})^{(I)}^{(K)}.$$
Let $`p`$ be the maximal value of $`|I|=i_1+\mathrm{}+i_N`$ for $`I`$ such that $`\alpha _I0`$. We want to show that $`p1`$. Among all $`I`$ such that $`|I|=p`$ there will be some (nonzero) $`\alpha _I`$ of maximal degree $`d`$. Then without loss of generality we can change the basis $`_1,\mathrm{},_N`$ and assume that the coefficient $`\alpha _{(p,0,\mathrm{},0)}`$ is nonzero and of degree $`d`$. If $`p>1`$, then no nonzero term in the left-hand side of (4.5) has a third tensor factor of degree $`2p`$ or $`2p1`$ since $`2p1>p`$. Hence, terms from the right-hand side of degree $`2p`$ (respectively $`2p1`$) in the third tensor factor must cancel against each other.
Terms having degree $`2p`$ in the third tensor factor cancel, since they give the following sum:
(4.6)
$$\underset{|I|=|K|=p}{}\alpha _K\alpha _I[^{(I)},^{(K)}],$$
which in the third tensor factor has degree $`2p1`$ and lower. Note also that their coefficients have total degree at most $`2d`$.
Terms having third tensor factors of degree $`2p1`$, besides (4.6), arise when we choose $`|I+K|=2p1`$. Those with $`|I|=p1,|K|=p`$ can be expressed in terms of commutators as above, and hence only contribute to lower degree. So, we only need to account for terms with $`|I|=p,|K|=p1`$.
Let us focus on such terms having a third tensor factor proportional to $`_1^{2p1}`$, whose coefficient must be zero. They occur in (4.5) only when $`I=(p,0,\mathrm{},0)`$, $`K=(p1,0,\mathrm{},0)`$. When $`J=0`$, things cancel as above. The only other nonzero terms are the following:
$$\underset{j}{}(\alpha _{K+\epsilon _j}\alpha _I_j\alpha _I_j\alpha _{K+\epsilon _j})^{(I)}^{(K)},$$
where $`\{\epsilon _j\}`$ is the standard basis of $`^N`$.
We have seen that all other contributions have coefficients of degree at most $`2d`$, so the sum $`_j(\alpha _{K+\epsilon _j}\alpha _I_j\alpha _I_j\alpha _{K+\epsilon _j})`$ must lie inside $`\mathrm{F}^{2d}(HH)`$. All $`\alpha _{K+\epsilon _j}`$ are of degree at most $`d`$ and $`\alpha _I_j`$ are of degree exactly $`d+1`$, hence $`_j\alpha _{K+\epsilon _j}\alpha _I_j`$ must lie in $`\mathrm{F}^{2d}(HH)`$ too. But this implies that $`\alpha _{K+\epsilon _j}\mathrm{F}^{d1}H`$ for all $`j`$, so in particular $`\alpha _I\mathrm{F}^{d1}H`$, which is a contradiction.
This proves that $`\alpha (𝔡+𝐤)(𝔡+𝐤)`$. Now if $`\alpha =r+s_1s_2`$, where $`r𝔡𝔡`$, $`s𝔡`$, then we have:
$$(\mathrm{\Delta }\mathrm{id})(\alpha )=r_{13}+r_{23}+s_1+s_2s_3,$$
and (4.2) becomes
(4.7)
$$\left([r_{12},r_{13}+s_1+s_2]+r_{12}s_3\right)+\text{cyclic}=0.$$
Comparing the terms in $`𝔡𝔡𝐤`$, we see that (4.7) is equivalent to the system (4.3, 4.4). ∎
Note that when $`\alpha =r𝔡𝔡`$, $`s=0`$, (4.4) is exactly the classical Yang–Baxter equation
(4.8)
$$[r_{12},r_{13}]+[r_{12},r_{23}]+[r_{13},r_{23}]=0.$$
Eq. (4.4) is a special case of the dynamical classical Yang–Baxter equation (see \[Fe, ES\]).
## 5. $`(H\mathrm{}𝐤[\mathrm{\Gamma }])`$-Pseudoalgebras
Let again $`H`$ be a cocommutative Hopf algebra. Let $`\mathrm{\Gamma }`$ be a group acting on $`H`$ by automorphisms, and let $`\stackrel{~}{H}=H\mathrm{}𝐤[\mathrm{\Gamma }]`$ be the smash product of $`H`$ with the group algebra of $`\mathrm{\Gamma }`$. As an associative algebra this is the semidirect product of $`H`$ with $`𝐤[\mathrm{\Gamma }]`$, while as a coalgebra it is the tensor product of coalgebras.
We will denote the action of $`\mathrm{\Gamma }`$ on $`H`$ by $`gf`$ for $`g\mathrm{\Gamma }`$, $`fH`$; then $`gf=gfg^1`$. Then a left $`\stackrel{~}{H}`$-module $`L`$ is the same as an $`H`$-module together with an action of $`\mathrm{\Gamma }`$ on it which is compatible with that of $`H`$, i.e., such that $`(gf)l=g(f(g^1l))`$ for $`g\mathrm{\Gamma }`$, $`fH`$, $`lL`$.
In this section we will study the relationship between the pseudotensor categories $`^{}(\stackrel{~}{H})`$ and $`^{}(H)`$. In particular, we will show that an $`\stackrel{~}{H}`$-pseudoalgebra is the same as an $`H`$-pseudoalgebra on which the group $`\mathrm{\Gamma }`$ acts by preserving the pseudoproduct.
Let us start by defining maps $`\delta _I:\stackrel{~}{H}^IH^I_H\stackrel{~}{H}`$ for each finite nonempty set $`I`$. It is enough to define $`\delta _I`$ on elements of the form $`_{iI}f_ig_i`$ where $`f_iH`$, $`g_i\mathrm{\Gamma }`$, in which case we let
$$\delta _I(_{iI}f_ig_i)=\{\begin{array}{cc}(_{iI}f_i)_Hg,\hfill & \text{if all }g_i\text{ are equal to some }g,\hfill \\ 0,\hfill & \text{if some of }g_i\text{ are different}.\hfill \end{array}$$
It is easy to see that $`\delta _I`$ is a homomorphism of both left $`H^I`$-modules and of right $`\stackrel{~}{H}`$-modules.
This allows us to define a pseudotensor functor $`\delta :^{}(\stackrel{~}{H})^{}(H)`$ as follows. For an object $`L`$ (a left $`\stackrel{~}{H}`$-module), we let $`\delta (L)L`$ be the left $`H`$-module obtained by restricting the action of $`\stackrel{~}{H}`$ to $`H\stackrel{~}{H}`$. For a polylinear map $`\varphi \mathrm{Lin}(\{L_i\}_{iI},M)`$ in $`^{}(\stackrel{~}{H})`$, i.e., for a homomorphism of left $`\stackrel{~}{H}^I`$-modules
$$\varphi :_{iI}L_i\stackrel{~}{H}^I_{\stackrel{~}{H}}M,$$
we let $`\delta (\varphi )`$ be the composition
$$\delta (\varphi ):_{iI}L_i\stackrel{\mathit{\varphi }}{}\stackrel{~}{H}^I_{\stackrel{~}{H}}M\stackrel{\delta _I_{\stackrel{~}{H}}\mathrm{id}}{}(H^I_H\stackrel{~}{H})_{\stackrel{~}{H}}M\stackrel{}{}H^I_HM.$$
This is a homomorphism of left $`H^I`$-modules, i.e., a polylinear map in $`^{}(H)`$. Moreover, since the maps $`\delta _I`$ are compatible with the actions of the symmetric groups and with the comultiplication of $`\stackrel{~}{H}`$, it follows that $`\delta `$ is compatible with the actions of the symmetric groups and with compositions of polylinear maps, i.e., it is a pseudotensor functor.
As usual, the action of $`\mathrm{\Gamma }`$ on $`H`$ can be extended to an action of $`\mathrm{\Gamma }`$ on $`H^I`$ by using the comultiplication $`\mathrm{\Delta }^{(I)}(g)=_{iI}g`$. Hence, $`\mathrm{\Gamma }`$ also acts on $`H^I_HM`$ by the formula
$$g\left((_{iI}f_i)_Hm\right)=(_{iI}gf_i)_Hgm,g\mathrm{\Gamma },f_iH,mM.$$
Then it is easy to see that $`\psi =\delta (\varphi )`$ has the following property:
(5.1)
$$\psi (_{iI}gl_i)=g\psi (_{iI}l_i),g\mathrm{\Gamma },l_iL_i,$$
in other words, it commutes with the action of $`\mathrm{\Gamma }`$.
We let $`_\mathrm{\Gamma }^{}(H)`$ be the subcategory of $`^{}(H)`$ with objects left $`\stackrel{~}{H}`$-modules, and with polylinear maps those polylinear maps $`\psi `$ of $`^{}(H)`$ that commute with the action of $`\mathrm{\Gamma }`$ (see (5.1)). This is a pseudotensor category, and $`\delta `$ is a pseudotensor functor from $`^{}(\stackrel{~}{H})`$ to $`_\mathrm{\Gamma }^{}(H)`$.
###### Theorem 5.1.
If $`\mathrm{\Gamma }`$ is a finite group, the functor $`\delta :^{}(\stackrel{~}{H})_\mathrm{\Gamma }^{}(H)`$ constructed above is an equivalence of pseudotensor categories.
###### Proof.
We will construct a pseudotensor functor $`\mathrm{\Sigma }`$ from $`_\mathrm{\Gamma }^{}(H)`$ to $`^{}(\stackrel{~}{H})`$. On objects $`L`$ we let $`\mathrm{\Sigma }(L)=L`$. In order to define it on polylinear maps, we need to find out how $`\varphi `$ can be recovered from $`\delta (\varphi )`$ and the action of $`\mathrm{\Gamma }`$.
Denote by $`\iota `$ the embedding $`H\stackrel{~}{H}`$, and let $`\pi _I`$ be the composition
$$\pi _I:H^I_H\stackrel{~}{H}\stackrel{\iota ^I_H\mathrm{id}}{}\stackrel{~}{H}^I_H\stackrel{~}{H}\stackrel{~}{H}^I_{\stackrel{~}{H}}\stackrel{~}{H}\stackrel{}{}\stackrel{~}{H}^I.$$
Explicitly, $`\pi _I`$ is given by the formula
$$\pi _I((_{iI}f_i)_Hg)=_{iI}f_ig,f_iH,g\mathrm{\Gamma }.$$
This is a homomorphism of both left $`H^I`$-modules and of right $`\stackrel{~}{H}`$-modules. Moreover, for $`f_iH`$, $`g_i\mathrm{\Gamma }`$, we have:
$$\pi _I\delta _I(_{iI}f_ig_i)=\{\begin{array}{cc}_{iI}f_ig_i,\hfill & \text{if all }g_i\text{ are equal},\hfill \\ 0,\hfill & \text{otherwise}.\hfill \end{array}$$
The crucial observation, which will allow us to invert $`\delta _I`$, is that for any $`h_i\stackrel{~}{H}`$, $`g_i\mathrm{\Gamma }`$, we have:
(5.2)
$$\underset{(g_i)\mathrm{\Gamma }^I/\mathrm{\Gamma }}{}(_{iI}g_i)(\pi _I\delta _I)(_{iI}g_i^1h_i)=_{iI}h_i.$$
Here $`\mathrm{\Gamma }`$ acts diagonally on $`\mathrm{\Gamma }^I`$ from the right; the left-hand side of (5.2) is invariant under $`(g_i)(g_ig)`$.
Given a polylinear map $`\psi \mathrm{Lin}(\{L_i\}_{iI},M)`$ in $`_\mathrm{\Gamma }^{}(H)`$, we can extend it to a map
$$\stackrel{~}{\psi }:_{iI}L_i\stackrel{𝜓}{}H^I_HM\stackrel{}{}(H^I_H\stackrel{~}{H})_{\stackrel{~}{H}}M\stackrel{\pi _I_{\stackrel{~}{H}}\mathrm{id}}{}\stackrel{~}{H}^I_{\stackrel{~}{H}}M.$$
(Note, however, that $`\stackrel{~}{\psi }`$ is not $`\stackrel{~}{H}^I`$-linear.) Now we define $`\mathrm{\Sigma }\psi :_{iI}L_i\stackrel{~}{H}^I_{\stackrel{~}{H}}M`$ by the formula:
(5.3)
$$(\mathrm{\Sigma }\psi )(_{iI}l_i)=\underset{(g_i)\mathrm{\Gamma }^I/\mathrm{\Gamma }}{}((_{iI}g_i)_{\stackrel{~}{H}}1)\stackrel{~}{\psi }(_{iI}g_i^1l_i).$$
It is easy to check that $`\mathrm{\Sigma }\psi `$ is $`\stackrel{~}{H}^I`$-linear, so it is a polylinear map in $`^{}(\stackrel{~}{H})`$. Moreover, $`\delta \mathrm{\Sigma }\psi =\psi `$. For a polylinear map $`\varphi \mathrm{Lin}(\{L_i\}_{iI},M)`$ in $`^{}(\stackrel{~}{H})`$, it is immediate from (5.2) and the $`\stackrel{~}{H}^I`$-linearity of $`\varphi `$ that $`\mathrm{\Sigma }\delta \varphi =\varphi `$. Therefore, $`\mathrm{\Sigma }:_\mathrm{\Gamma }^{}(H)^{}(\stackrel{~}{H})`$ is a pseudotensor functor inverse to $`\delta `$. ∎
###### Remark 5.2.
The above theorem holds also for infinite groups $`\mathrm{\Gamma }`$ if we restrict ourselves to polylinear maps $`\psi `$ of $`_\mathrm{\Gamma }^{}(H)`$ satisfying the following finiteness condition:
(5.4)
$$\psi (_{iI}g_il_i)0\text{for only a finite number of }(g_i)\mathrm{\Gamma }\mathrm{\Gamma }^I$$
for any fixed $`l_iL`$. (Note that, by (5.1), this condition does not depend on the choice of representatives $`(g_i)`$.) Indeed, the only place in the proof where we used the finiteness of $`\mathrm{\Gamma }`$ was to insure that the right-hand side of (5.3) is a finite sum.
If $`\psi =\delta (\varphi )`$ comes from a polylinear map $`\varphi `$ of $`^{}(\stackrel{~}{H})`$, then it satisfies (5.4), because $`\varphi `$ is $`\stackrel{~}{H}^I`$-linear and for any element $`h\stackrel{~}{H}^I`$ one has $`\delta _I((_{iI}g_i)h)0`$ for only a finite number of $`(g_i)\mathrm{\Gamma }\mathrm{\Gamma }^I`$.
Therefore, $`\delta :^{}(\stackrel{~}{H})_{\mathrm{\Gamma },\text{fin}}^{}(H)`$ is an equivalence of pseudotensor categories, where $`_{\mathrm{\Gamma },\text{fin}}^{}(H)`$ is the subcategory of $`_\mathrm{\Gamma }^{}(H)`$ consisting of polylinear maps $`\psi `$ satisfying (5.4).
###### Corollary 5.3.
A Lie $`\stackrel{~}{H}=(H\mathrm{}𝐤[\mathrm{\Gamma }])`$-pseudoalgebra $`L`$ is the same as a Lie $`H`$-pseudoalgebra $`L`$ on which the group $`\mathrm{\Gamma }`$ acts in a way compatible with the action of $`H`$, by preserving the $`H`$-pseudobracket:
(5.5)
$$[gagb]=g[ab]\text{for }g\mathrm{\Gamma }\text{}a,bL,$$
and satisfying the following finiteness condition:
(5.6) given $`a,bL`$, $`[gab]`$ is nonzero for only a finite number of $`g\mathrm{\Gamma }`$.
The $`\stackrel{~}{H}`$-pseudobracket of $`L`$ is given by the formula:
(5.7)
$$[a\stackrel{~}{}b]=\underset{g\mathrm{\Gamma }}{}\left((g^11)_{\stackrel{~}{H}}1\right)[gab],a,bL.$$
A similar statement holds for representations, as well as for associative pseudoalgebras.
This result, combined with Kostant’s Theorem 2.1, will allow us in many cases to reduce the study of $`H`$-pseudoalgebras to the case when $`H`$ is a universal enveloping algebra (see Section 13.7).
###### Example 5.4.
Let $`\mathrm{\Gamma }`$ be a subgroup of $`𝐤^{}`$ and let
$$H=𝐤[]\mathrm{}𝐤[\mathrm{\Gamma }]=\underset{m_+,\alpha \mathrm{\Gamma }}{}𝐤^mT_\alpha $$
with multiplication $`T_\alpha T_\beta =T_{\alpha \beta }`$, $`T_1=1`$, $`T_\alpha T_\alpha ^1=\alpha `$ and comultiplication $`\mathrm{\Delta }()=1+1`$, $`\mathrm{\Delta }(T_\alpha )=T_\alpha T_\alpha `$. Then the notion of a Lie $`H`$-pseudoalgebra is equivalent to that of a $`\mathrm{\Gamma }`$-conformal algebra (cf. \[K4\]).
###### Example 5.5.
Let now $`H=𝐤[]\times F(\mathrm{\Gamma })`$, where $`F(\mathrm{\Gamma })`$ is the function algebra of a finite abelian group $`\mathrm{\Gamma }`$. In other words, $`H=_{m_+,\alpha \mathrm{\Gamma }}𝐤^m\pi _\alpha `$ with multiplication $`\pi _\alpha \pi _\beta =\delta _{\alpha ,\beta }\pi _\alpha `$, $`\pi _\alpha =\pi _\alpha `$ and comultiplication $`\mathrm{\Delta }()=1+1`$, $`\mathrm{\Delta }(\pi _\alpha )=_{\gamma \mathrm{\Gamma }}\pi _{\alpha \gamma ^1}\pi _\gamma `$. Then one gets the notion of a $`\mathrm{\Gamma }`$-twisted conformal algebra (cf. \[K4\]).
## 6. A Digression to Linearly Compact Lie Algebras
We will view the base field $`𝐤`$ as a topological field with discrete topology. A topological vector space $``$ over $`𝐤`$ is called *linearly compact* if it is the space of all linear functionals on a vector space $`𝒱`$ with discrete topology, with the topology on $``$ defined by taking all subspaces $`\{U^{}|U𝒱,dimU<\mathrm{}\}`$ as a fundamental system of neighborhoods of $`0`$ in $``$. Here, as usual, $`U^{}`$ denotes the subspace of $``$ consisting of all linear functionals vanishing on $`U`$.
In general, given a topological vector space $`𝒲`$, we define a topology on $`𝒲^{}`$ by taking for the fundamental system of neighborhoods of $`0`$ the subspaces $`U^{}`$ where $`U`$ is a linearly compact subspace of $`𝒲`$.
Several equivalent definitions of linear-compactness are provided by the next proposition.
###### Proposition 6.1.
For a topological vector space $``$ over the topological field $`𝐤`$ the following statements are equivalent:
1. $``$ is the dual of a discrete vector space.
2. The topological dual $`^{}`$ of $``$ is a discrete topological space.
3. $``$ is the topological product of finite-dimensional discrete vector spaces.
4. $``$ is the projective limit of finite-dimensional discrete vector spaces.
5. $``$ has a collection of finite-codimensional open subspaces whose intersection is $`\{0\}`$, with respect to which it is complete.
###### Proof.
Can be found in \[G1\]. ∎
###### Remark 6.2.
For both discrete and linearly compact vector spaces, the canonical map from $``$ to $`^{}`$ is an isomorphism.
A linearly compact (associative or Lie) algebra is a topological (associative or Lie) algebra for which the underlying topological space is linearly compact.
The basic example of a linearly compact associative algebra is the algebra $`𝒪_N=𝐤[[t_1,\mathrm{},t_N]]`$ of formal power series over $`𝐤`$ in $`N1`$ indeterminates $`t_1,\mathrm{},t_N`$, with the usual formal topology for which $`(t_1,\mathrm{},t_N)^j`$, the powers of the ideal $`(t_1,\mathrm{},t_N)`$, form a fundamental system of neighborhoods of $`𝒪_N`$.
###### Remark 6.3.
The topological vector spaces $`𝒪_N`$ are isomorphic and characterized among linearly compact vector spaces by each of the following properties:
1. $`𝒪_N^{}`$ is countable-dimensional.
2. $`𝒪_N`$ has a filtration by open subspaces.
###### Remark 6.4.
(i) One defines a completed tensor product of two linearly compact vector spaces $`𝒱,𝒲`$ by $`𝒱\widehat{}𝒲=(𝒱^{}𝒲^{})^{}`$ where we put the discrete topology on $`𝒱^{}𝒲^{}`$. Then $`𝒱\widehat{}𝒲`$ is linearly compact.
(ii) With this definition, $`𝒪_{M+N}𝒪_M\widehat{}𝒪_N`$ as topological algebras.
(iii) Given a commutative associative linearly compact algebra $`𝒪`$ and a linearly compact Lie algebra $``$, their completed tensor product $`𝒪\widehat{}`$ is again a linearly compact Lie algebra.
The basic example of a linearly compact Lie algebra is the Lie algebra $`W_N`$ of continuous derivations of the topological algebra $`𝒪_N`$. The filtration
$$\mathrm{F}_j𝒪_N=(t_1,\mathrm{},t_N)^{j+1},j=1,0,1,\mathrm{}$$
of $`𝒪_N`$ induces the *canonical filtration* $`\mathrm{F}_jW_N`$ of $`W_N`$, where
$$\mathrm{F}_jW_N=\{DW_N|D(\mathrm{F}_i𝒪_N)\mathrm{F}_{i+j}𝒪_Ni\},j=1,0,1,\mathrm{}.$$
It is clear that $`W_N`$ consists of all linear differential operators of the form
$$D=\underset{i=1}{\overset{N}{}}P_i(t)\frac{}{t_i},\text{ where }P_i(t)𝒪_N,$$
and that $`\mathrm{F}_jW_N`$ $`(j1)`$ consists of those $`D`$ for which all $`P_i(t)`$ lie in $`\mathrm{F}_j𝒪_N`$.
Let $`E=_{i=1}^Nt_i\frac{}{t_i}`$ be the Euler operator. The spectrum of $`\mathrm{ad}E`$ consists of all integers $`j1`$, and, denoting by $`W_{N;j}`$ the $`j`$-th eigenspace of $`\mathrm{ad}E`$ we obtain the *canonical* $``$*-gradation*:
$$W_N=\underset{j1}{}W_{N;j},[W_{N;i},W_{N;j}]W_{N;i+j}.$$
The following fact is well known.
###### Lemma 6.5.
$`W_{N;0}𝔤𝔩_N(𝐤)`$ and one has the following isomorphism of $`𝔤𝔩_N(𝐤)`$-modules:
$$W_{N;j}𝐤^N(\mathrm{S}^{j+1}𝐤^N)^{}.$$
Furthermore, one has a decomposition into a direct sum of irreducible submodules: $`W_{N;j}=W_{N;j}^{}+W_{N;j}^{\prime \prime }`$, where $`W_{N;j}^{}(\mathrm{S}^j𝐤^N)^{}`$ $`(=0`$ if $`j=1)`$ and $`W_{N;j}^{\prime \prime }`$ the highest component of $`𝐤^N(\mathrm{S}^{j+1}𝐤^N)^{}`$. The subspace $`𝔭=W_{N;1}+W_{N;0}+W_{N;1}^{}`$ is a subalgebra of $`W_N`$ isomorphic to $`𝔰𝔩_{N+1}(𝐤)`$.
Let $`\mathrm{\Omega }_N=_{j=0}^N\mathrm{\Omega }_{N;j}`$ denote the algebra of differential forms over $`𝒪_N`$. The defining representation of $`W_N`$ on $`𝒪_N`$ extends uniquely to a representation on $`\mathrm{\Omega }_N`$ commuting with the differential $`\mathrm{d}`$.
Recall that a *volume form* is a differential $`N`$-form $`v=f(t_1,\mathrm{},t_N)\mathrm{d}t_1\mathrm{}\mathrm{d}t_N`$ such that $`f(0)0`$, a *symplectic form* is a closed $`2`$-form $`s=_{i,j=1}^Ns_{ij}(t_1,\mathrm{},t_N)\mathrm{d}t_i\mathrm{d}t_j`$ such that $`det(s_{ij}(0))0`$, and a *contact form* is a $`1`$-form $`c`$ such that $`c(\mathrm{d}c)^{(N1)/2}`$ is a volume form. The following facts are well known.
###### Lemma 6.6.
(i) Any volume form can be transformed by an automorphism of $`𝒪_N`$ to the standard volume form $`v_0=\mathrm{d}t_1\mathrm{}\mathrm{d}t_N`$.
(ii) A symplectic form exists iff $`N`$ is even, $`N=2n`$, and by an automorphism of $`𝒪_N`$ it can be transformed to the standard symplectic form $`s_0=_{i=1}^N\mathrm{d}t_i\mathrm{d}t_{n+i}`$.
(iii) A contact form exists iff $`N`$ is odd, $`N=2n+1`$, and by an automorphism of $`𝒪_N`$ it can be brought to the standard contact form $`c_0=\mathrm{d}t_N+_{i=1}^nt_i\mathrm{d}t_{n+i}`$.
Consider the following (closed) subalgebras of the Lie algebra $`W_N`$:
$`S_N(v)`$ $`=\{DW_N|Dv=0\}(N2),`$
$`H_N(s)`$ $`=\{DW_N|Ds=0\}(N\text{ even }2),`$
$`K_N(c)`$ $`=\{DW_N|Dc=fc\text{ for some }f𝒪_N\}(N\text{ odd }3).`$
Let also $`S_N=S_N(v_0)`$, $`H_N=H_N(s_0)`$, $`K_N=K_N(c_0)`$. Lemma 6.6 implies isomorphisms: $`S_N(v)S_N`$, $`H_N(s)H_N`$, $`K_N(c)K_N`$, $`S_2H_2`$.
The canonical filtration of $`W_N`$ induces *canonical filtrations* $`\mathrm{F}_jS_N(v):=\mathrm{F}_jW_NS_N(v)`$, etc. Note that $`dimW_N/\mathrm{F}_1W_N=N`$. A Lie subalgebra $``$ of $`W_N`$ is called transitive if $`dim/(\mathrm{F}_1W_N)=N`$. It is known that the Lie algebras $`W_N`$, $`S_N`$, $`H_N`$ and $`K_N`$ are transitive. In addition, the canonical filtrations $`\mathrm{F}_j`$ of these Lie algebras have the following transitivity property:
(6.1)
$$\mathrm{F}_{j+1}=\{a\mathrm{F}_j|[a,]\mathrm{F}_j\}.$$
Noting that $`Ev_0=Nv_0`$ and $`Es_0=2s_0`$, we conclude that $`\mathrm{ad}E`$ is an (outer) derivation of $`S_N`$ and $`H_N`$, hence the canonical gradation of $`W_N`$ induces *canonical* $``$*-gradations* $`S_N=_{j1}S_{N;j}`$ and $`H_N=_{j1}H_{N;j}`$.
Let $`E^{}=2t_N\frac{}{t_N}+_{i=1}^{N1}t_i\frac{}{t_i}`$. Then $`E^{}c_0=2c_0`$, hence $`E^{}K_N`$ and the eigenspace decomposition of $`\mathrm{ad}E^{}`$ defines the *canonical* $``$*-gradation* $`K_N=_{j2}K_{N;j}`$. The following facts are well known.
###### Lemma 6.7.
(i) $`S_{N;0}𝔰𝔩_N(𝐤)`$, $`H_{N;0}𝔰𝔭_N(𝐤)`$, $`K_{N;0}𝔠𝔰𝔭_{N1}(𝐤)`$ $`(𝔰𝔭_{N1}(𝐤)𝐤)`$.
(ii) The $`S_{N;0}`$-module $`S_{N;j}`$ is isomorphic to the highest component of the $`𝔰𝔩_N(𝐤)`$-module $`𝐤^N(\mathrm{S}^{j+1}𝐤^N)^{}`$.
(iii) The $`H_{N;0}`$-module $`H_{N;j}`$ is isomorphic to the $`(`$irreducible$`)`$ $`𝔰𝔭_N(𝐤)`$-module $`\mathrm{S}^{j+2}𝐤^N`$.
(iv) $`K_{N;0}=𝔰𝔭_{N1}(𝐤)𝐤E^{}`$ and the $`𝔰𝔭_{N1}(𝐤)`$-module $`K_{N;j}`$ decomposes into the following direct sum of irreducible modules:
$$K_{N;j}=\underset{i=0}{\overset{\left[\frac{j}{2}\right]+1}{}}K_{N;j}^{(i)},\text{ where }K_{N;j}^{(i)}\mathrm{S}^{j+22i}𝐤^{N1}.$$
The subspace $`𝔭=K_{N;2}+K_{N;1}+K_{N;0}+K_{N;1}^{(1)}+K_{N;2}^{(2)}`$ is a subalgebra of $`K_N`$ isomorphic to $`𝔰𝔭_{N+1}(𝐤)`$.
The following celebrated theorem goes back to E. Cartan (see \[G2\] for a relatively simple proof).
###### Theorem 6.8.
Any infinite-dimensional simple linearly compact Lie algebra is isomorphic to one of the topological Lie algebras $`W_N`$, $`S_N`$, $`H_N`$ or $`K_N`$.
Let $`𝔤`$ be a Lie algebra, and let $`𝔥`$ be its subalgebra of codimension $`N`$. Then $`F=\mathrm{Hom}_{U(𝔥)}(U(𝔤),𝐤)`$, with the product $`(f_1f_2)(u)=f_1(u_{(1)})f_2(u_{(2)})`$, is (canonically) isomorphic to the algebra of formal power series on $`(𝔤/𝔥)^{}`$ \[B2\], which is (non-canonically) isomorphic to the linearly compact algebra $`𝒪_N`$. The Lie algebra $`D`$ of continuous derivations of $`F`$ is then isomorphic to $`W_N`$. $`F`$ has a canonical $`𝔤`$-action induced by the left-multiplication $`𝔤`$-action on $`U(𝔤)`$, which gives us a homomorphism $`\gamma `$ of $`𝔤`$ to $`W_N`$. (This is non-canonical since the identification of $`F`$ with $`𝒪_N`$ is not canonical.)
We will use in the sequel the following theorem of Guillemin and Sternberg \[GS\] (see \[B2\] for a simple proof).
###### Proposition 6.9.
Let $`𝔤`$ be a Lie algebra, and let $`𝔥`$ be its subalgebra of codimension $`N`$. Provided that $`𝔥`$ contains no nonzero ideals of $`𝔤`$, the above-defined $`\gamma `$ is a Lie algebra isomorphism of $`𝔤`$ with a subalgebra of $`W_N`$ which maps $`𝔥`$ into $`\mathrm{F}_0W_N`$.
$`(`$Conversely, if $`𝔤W_N`$ maps $`𝔥`$ into $`\mathrm{F}_0W_N`$, then $`𝔥`$ doesn’t contain nonzero ideals of $`𝔤.)`$ Every Lie algebra homomorphism of $`𝔤`$ to $`W_N`$, which coincides with $`\gamma `$ modulo $`\mathrm{F}_0W_N`$, is conjugated to $`\gamma `$ via a unique automorphism of $`𝒪_N`$.
We have the following important property of the filtrations on $`H`$ and $`X=H^{}`$, defined in Section 2.2.
###### Lemma 6.10.
Let $`H=U(𝔤)\mathrm{}𝐤[\mathrm{\Gamma }]`$ be a cocommutative Hopf algebra, and $`X=H^{}`$. If $`h\mathrm{F}^iU(𝔤)H`$ but $`h\mathrm{F}^{i1}U(𝔤)`$, then $`h\mathrm{F}_nX=\mathrm{F}_{ni}X`$. In particular, for any $`h𝔤\{0\}`$ and for every open subspace $`UX`$, there is some $`n`$ such that $`h^nU=X`$. Similar statements hold for the right action of $`h`$ as well.
###### Proof.
By the construction of the filtrations it is evident that we can assume $`H=U(𝔤)`$. Then $`X𝒪_N`$ $`(N=dim𝔤)`$, and $`𝔤W_N`$ acts on it by linear differential operators. The rest of the proof is clear. ∎
The following result from \[G1, G2\] will be essential for our purposes.
###### Proposition 6.11.
(i) A linearly compact Lie algebra $``$ satisfies the descending chain condition on closed ideals if and only if it has a fundamental subalgebra, i.e., an open subalgebra containing no ideals of $``$.
(ii) When either of the assumptions of (i) holds, the noncommutative minimal closed ideals of $``$ are of the form $`𝒪_r\widehat{}𝔰`$ where $`𝔰`$ is a simple linearly compact Lie algebra and $`r_+`$.
We will also need the following examples of non-simple linearly compact Lie algebras:
$`CS_N(v)`$ $`=\{DW_N|Dv=av,a𝐤\},`$
$`CH_N(s)`$ $`=\{DW_N|Ds=as,a𝐤\}.`$
As before, we have isomorphisms $`CS_N(v)CS_NCS_N(v_0)`$ and $`CH_N(s)CH_NCH_N(s_0)`$. Note also that $`CS_N=𝐤ES_N`$ and $`CH_N=𝐤EH_N`$. Another important example of a non-simple linearly compact Lie algebra is the Poisson algebra $`P_N`$, which is $`𝒪_N`$ $`(N=2n)`$ endowed with the Poisson bracket:
$$\{f,g\}=\underset{i=1}{\overset{n}{}}\frac{f}{t_i}\frac{g}{t_{n+i}}\frac{f}{t_{n+i}}\frac{g}{t_i}.$$
It is a nontrivial central extension of $`H_N`$:
$$0𝐤P_N\stackrel{𝜑}{}H_N0,$$
where $`\phi (f)=_{i=1}^n\frac{f}{t_i}\frac{}{t_{n+i}}\frac{f}{t_{n+i}}\frac{}{t_i}`$.
We can describe also $`K_N`$ in a more explicit way, similar to the above description of $`P_N`$. For $`f,g𝒪_N`$, define
$$\{f,g\}^{}=\{f,g\}_{2n}+\frac{f}{t_{2n+1}}(E_{2n}g2g)(E_{2n}f2f)\frac{g}{t_{2n+1}},$$
where $`\{f,g\}_{2n}`$ is the Poisson bracket taken with respect to the variables $`t_1,\mathrm{},t_{2n}`$ and $`E_{2n}`$ is the Euler operator $`_{i=1}^{2n}t_i\frac{}{t_i}`$. If we define
$$\psi (f)=\underset{i=1}{\overset{n}{}}\left(\frac{f}{t_i}\frac{}{t_{n+i}}\frac{f}{t_{n+i}}\frac{}{t_i}\right)+\frac{f}{t_{2n+1}}E_{2n}(E_{2n}f2f)\frac{}{t_{2n+1}},$$
then we have $`\psi (f)g=\{f,g\}^{}+2\frac{f}{t_{2n+1}}g`$. It is easy to see that $`[\psi (f),\psi (g)]=\psi (\{f,g\}^{})`$ and $`\psi (f)c_0=2\frac{f}{t_{2n+1}}c_0`$. Thus $`K_N`$ is isomorphic to $`𝒪_N`$ with the bracket $`\{,\}^{}`$.
For a linearly compact Lie algebra $``$ denote by $`\mathrm{Der}`$ the Lie algebra of its continuous derivations and by $`\widehat{}`$ the universal central extension of $``$. Then we have:
###### Proposition 6.12.
(i) $`\mathrm{Der}W_N=W_N`$, $`\mathrm{Der}S_N=CS_N`$, $`\mathrm{Der}H_N=CH_N`$, $`\mathrm{Der}K_N=K_N`$.
(ii) $`\mathrm{Der}(𝒪_r\widehat{})=W_r1+𝒪_r\widehat{}\mathrm{Der}`$ for any simple linearly compact Lie algebra $``$.
(iii) The Lie algebras $`𝒪_r\widehat{}W_N`$, $`𝒪_r\widehat{}S_N`$ $`(`$for $`N>2)`$ and $`𝒪_r\widehat{}K_N`$ have no nontrivial central extensions. The universal central extension of $`𝒪_r\widehat{}H_N`$ is $`𝒪_r\widehat{}P_N`$.
(iv) If $`𝔤`$ is a simple finite-dimensional Lie algebra, then $`(𝒪_r𝔤)^\widehat{}=(𝒪_r𝔤)+(\mathrm{\Omega }_{r;1}/\mathrm{d}𝒪_r)`$ with the bracket
$$[fa,gb]^\widehat{}=fg[a,b]+(a|b)f\mathrm{d}gmod\mathrm{d}𝒪_r,$$
where $`(a|b)`$ is the Killing form on $`𝔤`$.
###### Proof.
For a proof of (iv) see \[Ka\].
In order to prove (ii), notice that if $`d`$ is a derivation of $`𝒪_r\widehat{}`$, then its action on $`1`$ is given by $`d(1x)=_ia_id_i(x)`$ for all $`x`$, where the $`a_i`$ form a topological basis of $`𝒪_r`$ and the $`d_i`$ are continuous derivations of $``$. Subtracting $`_ia_id_i`$ from $`d`$, we get a derivation $`\stackrel{~}{d}`$ acting trivially on $`1`$. We are going to show that if $``$ is simple then $`\stackrel{~}{d}`$ is of the form $`D1`$ where $`D\mathrm{Der}𝒪_r=W_r`$.
Let us fix $`P𝒪_r`$. Then $`\stackrel{~}{d}(Px)`$ can be written as $`_ia_if_i(x)`$, where $`f_i`$ are continuous $`𝐤`$-endomorphisms of $``$. From $`\stackrel{~}{d}([Px,1y])=[\stackrel{~}{d}(Px),1y]`$ we see that $`[f_i(x),y]=f_i([x,y])`$ for every $`x,y`$. This means that $`f_i`$ commutes with $`\mathrm{ad}y`$ for all $`y`$. By a Schur’s lemma argument \[G1, Proposition 4.4\] and simplicity of $``$, we conclude that the $`f_i`$ are multiples of the identity map, hence $`\stackrel{~}{d}(Px)=a_Px`$ for some $`a_P𝒪_r`$ and all $`x`$. It is now immediate to check that the mapping $`D:Pa_P`$ is indeed a derivation of $`𝒪_r`$, proving (ii).
In order to prove the rest of the statements, denote by $`𝔞`$ the $`0`$th component of the canonical $``$-gradation of $`=W_N`$, $`S_N`$, $`H_N`$ or $`K_N`$. This is a reductive subalgebra of $``$, hence $`\mathrm{Der}=V`$, where $`[𝔞,V]V`$. But $`[V,]`$, hence $`[𝔞,V]=0`$, i.e., any element $`DV`$ defines an endomorphism of $``$ viewed as an $`𝔞`$-module. Since $`EW_N`$ and $`E^{}K_N`$, we conclude that $`D`$ also preserves the canonical gradation of these Lie algebras and we may assume that $`D`$ acts trivially on the $`(1)`$st component. Using the transitivity of $`W_N`$ and $`K_N`$, we conclude that $`D=0`$. By Lemma 6.7, all components of the canonical $``$-gradation of $`S_N`$ and $`H_N`$ are inequivalent $`𝔞`$-modules, hence $`D`$ preserves this gradation in this case as well. Subtracting from $`D`$ a multiple of $`E`$, we may assume that $`D`$ acts trivially on the $`(1)`$st component and, using transitivity, we conclude that $`D`$ is a multiple of $`E`$. Thus (i) is proved.
Since $`𝔞`$ acts completely reducibly on the space $`Z^2`$ of $`2`$-cocycles on $`𝒪_r\widehat{}`$ with values in $`𝐤`$, and since $`𝔞`$ acts trivially on cohomology, we may choose a subspace $`U`$ of $`Z^2`$, complementary to the space of trivial $`2`$-cocycles, on which $`𝔞`$ acts trivially. Hence for any $`2`$-cocycle $`\alpha U`$ we have: $`\alpha (a,b)=0`$ if $`aM_1`$, $`bM_2`$ and $`M_i`$ are irreducible non-contragredient $`𝔞`$-submodules of $`\overline{}:=𝒪_r\widehat{}`$. Let $`\overline{}_j=𝒪_r_j`$ for short, where $`_j`$ is the $`j`$th component of the canonical gradation.
It follows from Lemma 6.7(ii) that all pairs of $`𝔞`$-submodules in $`\overline{}=𝒪_r\widehat{}S_N`$ are non-contragredient, except for the adjoint $`𝔞`$-submodules in $`\overline{}_0=𝒪_rS_{N;0}`$. Thus, we have:
$$\alpha (a,b)=0\text{ if }a\overline{S}_{N;i},b\overline{S}_{N;j},i0\text{ or }j0.$$
Taking now $`a\overline{S}_{N;1}`$, $`b\overline{S}_{n;1}`$ and $`c\overline{S}_{N;0}`$, the cocycle condition
$$\alpha ([a,b],c)+\alpha ([b,c],a)+\alpha ([c,a],b)=0$$
gives $`\alpha ([a,b],c)=0`$. Since $`\overline{S}_{N;0}=[\overline{S}_{N;1},\overline{S}_{N;1}]`$, we conclude that $`\alpha =0`$. Hence all central extensions of $`\overline{S}_N`$ are trivial.
Likewise, $`\alpha `$ is zero on any pair of subspaces $`\overline{W}_{N;i},\overline{W}_{N;j}`$, unless $`i+j=0`$, and on the pair $`\overline{W}_{N;1},\overline{W}_{N;1}^{\prime \prime }`$ (see Lemma 6.5). Choosing $`a\overline{W}_{N;1}`$, $`b\overline{W}_{N;1}^{\prime \prime }`$, $`c\overline{W}_{N;0}`$, we obtain, as above, from the cocycle condition, that $`\alpha `$ is zero on the pair $`\overline{W}_{N;0},[\overline{W}_{N;0},\overline{W}_{N;0}]`$. It follows from (iv) applied to the subalgebra $`𝒪_r𝔰𝔩_{N+1}(𝐤)`$ of $`\overline{W}_N`$ (see Lemma 6.5) that $`\alpha `$ is zero on this subalgebra if $`N>1`$. Thus any cocycle on $`\overline{W}_N`$ $`(N>1)`$ is trivial. In the case of $`\overline{W}_1`$ the cocycle $`\alpha `$ is trivial. The case of $`\overline{K}_N`$ is similar.
In the remaining case of $`\overline{H}_N`$ we show, as above, that the cocycle $`\alpha `$ is trivial on any pair $`\overline{H}_{N;i},\overline{H}_{N;j}`$ if $`ij`$. Using the cocycle condition for $`a\overline{H}_{N;k}`$, $`b\overline{H}_{N;k+1}`$ and $`c\overline{H}_{N;1}`$, and the fact that $`\overline{H}_{N;k}=[\overline{H}_{N;k+1},\overline{H}_{N;1}]`$, we conclude that $`\alpha `$ is trivial on any pair $`\overline{H}_{N;i},\overline{H}_{N;i}`$ as well, unless $`i=1`$. It is easy to see that this implies that $`\widehat{\overline{H}}_N=\overline{P}_N`$. ∎
## 7. $`H`$-Pseudoalgebras and $`H`$-Differential Algebras
In this section, $`H`$ will be a cocommutative Hopf algebra, and as before, $`X=H^{}`$.
### 7.1. The annihilation algebra
Let $`Y`$ be an $`H`$-bimodule which is a commutative associative $`H`$-differential algebra both for the left and for the right action of $`H`$ (see (2.8), (2.19)); for example, $`Y=X:=H^{}`$.
For a left $`H`$-module $`L`$, let $`𝒜_YL=Y_HL`$. We define a left action of $`H`$ on $`𝒜_YL`$ in the obvious way:
(7.1)
$$h(x_Ha)=hx_Ha,hH,xY,aL.$$
If, in addition, $`L`$ is an $`H`$-pseudoalgebra with a pseudoproduct $`ab`$, we can define a product on $`𝒜_YL`$ by the formula:
(7.2)
$$\begin{array}{cc}\hfill (x_Ha)(y_Hb)& =_i(xf_i)(yg_i)_He_i,\hfill \\ \hfill \text{if}ab& =_i(f_ig_i)_He_i.\hfill \end{array}$$
By (2.19) and the $`H`$-bilinearity (3.23) of the pseudoproduct, it is clear that (7.2) is well defined.
###### Proposition 7.1.
If $`L`$ is a Lie $`H`$-pseudoalgebra, then $`𝒜_YL`$ is a Lie $`H`$-differential algebra, i.e., a Lie algebra which is also a left $`H`$-module so that
(7.3)
$$h[\alpha ,\beta ]=[h_{(1)}\alpha ,h_{(2)}\beta ],\text{for }hH,\alpha ,\beta 𝒜_YL.$$
Similarly, if $`L`$ is an associative $`H`$-pseudoalgebra, then $`𝒜_YL`$ is an associative $`H`$-differential algebra. A similar statement holds for modules as well: if $`M`$ is an $`L`$-module, then $`𝒜_YM`$ is an $`𝒜_YL`$-module with a compatible $`H`$-action so that
(7.4)
$$h(am)=(h_{(1)}a)(h_{(2)}m)\text{for}hH,a𝒜_YL,m𝒜_YM.$$
###### Proof.
Equation (7.3) follows from (2.8). The skew-commutativity of the bracket (7.2) follows immediately from that of $`[ab]`$. The proof of the Jacobi identity is straightforward by using (3.25). Let us check for example that the associativity of $`L`$ is equivalent to that of $`𝒜_YL`$; the case of the Jacobi identity is similar.
We will use the notation from (3.14)–(3.19), and we will write $`a_xx_Ha`$ for $`aL`$, $`xY`$. Then we want to compute the products $`a_x(b_yc_z)`$ and $`(a_xb_y)c_z`$. By definition, if we have (3.17) and (3.18), then
$`b_yc_z`$ $`=_i(yh_i)(zl_i)_Hd_i`$
and
$`a_x(b_yc_z)`$ $`=_{i,j}\left(xh_{ij}\right)\left(\left((yh_i)(zl_i)\right)l_{ij}\right)_Hd_i`$
$`=_{i,j}(xh_{ij})(yh_il_{ij}^{}{}_{(1)}{}^{})(zl_il_{ij}^{}{}_{(2)}{}^{})_Hd_i.`$
Similarly, if we have (3.14) and (3.15), then
$$(a_xb_y)c_z=_{i,j}(xf_if_{ij}^{}{}_{(1)}{}^{})(yg_if_{ij}^{}{}_{(2)}{}^{})(zg_{ij})_He_i.$$
Now recalling (3.16) and (3.19), we see that the associativity of $`L`$ is equivalent to that of $`𝒜_YL`$. ∎
###### Definition 7.2.
The $`H`$-differential algebra $`𝒜(L)𝒜_XL:=X_HL`$ is called the annihilation algebra of the pseudoalgebra $`L`$. We will write $`a_xx_Ha`$ for $`aL`$, $`xX`$.
###### Remark 7.3.
When $`L`$ is an associative $`H`$-pseudoalgebra, one does not need the cocommutativity of $`H`$ or the commutativity of $`Y`$ in order to define $`𝒜_YL`$ (cf. Remark 3.10).
###### Lemma 7.4.
Let $`H=U(𝔡)\mathrm{}𝐤[\mathrm{\Gamma }]`$, and let $`M`$ be a left $`H`$-module. If an element $`aM`$ is $`U(𝔡)`$-torsion, i.e., if $`ha=0`$ for some $`hU(𝔡)\{0\}`$, then $`X_Ha=0`$. In particular, for $`H=U(𝔡)`$, we have: $`𝒜(M)𝒜(M/\mathrm{Tor}M)`$, where $`\mathrm{Tor}M`$ is the torsion submodule of $`M`$.
###### Proof.
We have $`0=x_Hha=xh_Ha`$ for every $`xX`$. Since the right action of $`h`$ on $`X`$ is surjective (see Lemma 6.10), it follows that $`x_Ha=0`$ for any $`xX`$. ∎
### 7.2. The functor $`𝒜_Y`$
Analyzing the proof of Proposition 7.1, one can notice that the definition of $`𝒜_Y`$ is a special case of a more general construction which we describe below.
First, recall that a commutative associative $`H`$-differential algebra $`Y`$ is the same as a commutative associative algebra in the pseudotensor category $`^l(H)`$ from Example 3.8. We denote by $`^b(H)`$ the category of $`H`$-bimodules, provided with a pseudotensor structure given by (3.3), but with $`\mathrm{Hom}_H`$ there replaced by $`\mathrm{Hom}_{HH}`$ which means homomorphisms of $`H`$-bimodules. The composition of polylinear maps in $`^b(H)`$ is given again by (3.2). Then $`H`$-differential algebras $`Y`$ considered above are exactly the commutative associative algebras in $`^b(H)`$.
Instead of one $`H`$-bimodule $`Y`$ one can use several: for any collections of objects $`Y_i^b(H)`$ and $`L_i^{}(H)`$ $`(iI)`$ we can consider the left $`H`$-modules $`𝒜_{Y_i}L_i=Y_i_HL_i`$ as objects of $`^l(H)`$. Assume we are given polylinear maps $`f\mathrm{Lin}(\{Y_i\}_{iI},Z)`$ in $`^b(H)`$ and $`\varphi \mathrm{Lin}(\{L_i\}_{iI},M)`$ in $`^{}(H)`$. Then we define a polylinear map $`f_H\varphi \mathrm{Lin}(\{Y_i_HL_i\}_{iI},Z_HM)`$ in $`^l(H)`$ as the following composition:
$$\begin{array}{c}_{iI}(Y_i_HL_i)\stackrel{}{}(_{iI}Y_i)_{H^I}(_{iI}L_i)\hfill \\ \hfill \stackrel{\mathrm{id}\varphi }{}(_{iI}Y_i)_{H^I}(H^I_HM)\stackrel{}{}(_{iI}Y_i)_HM\stackrel{f\mathrm{id}}{}Z_HM.\end{array}$$
###### Proposition 7.5.
The above definition is compatible with compositions of polylinear maps in $`^b(H)`$, $`^{}(H)`$, and $`^l(H)`$:
$$f(\{g_i\}_{iI})_H\varphi (\{\psi _i\}_{iI})=\left(f_H\varphi \right)\left(\{g_i_H\psi _i\}_{iI}\right).$$
The proof of this proposition is straightforward and is left to the reader.
###### Corollary 7.6.
Let $`(Y,\nu )`$ be a commutative associative algebra in $`^b(H)`$. For a finite nonempty set $`I`$, let $`\nu ^{(I)}:Y^IY`$ be the iterated multiplication $`\nu (\nu \mathrm{id})\mathrm{}(\nu \mathrm{id}\mathrm{}\mathrm{id})`$. Recall that for an object $`L`$ in $`^{}(H)`$, we define $`𝒜_Y(L):=Y_HL`$. For a polylinear map $`\varphi \mathrm{Lin}(\{L_i\}_{iI},M)`$ in $`^{}(H)`$, let $`𝒜_Y(\varphi ):=\nu ^{(I)}_H\varphi `$. Then $`𝒜_Y`$ is a pseudotensor functor from $`^{}(H)`$ to $`^l(H)`$.
As a special case of this corollary, we obtain Proposition 7.1.
Let us give another application of Proposition 7.5. An instance of an $`H`$-bimodule is $`H`$ itself (however, $`H`$ is not an $`H`$-differential algebra!). The coproduct $`\mathrm{\Delta }:HHH`$, the evaluation map $`ev:XH𝐤`$, and the isomorphism $`𝐤H\stackrel{}{}H`$ are all homomorphisms of $`H`$-bimodules, so the composition
(7.5)
$$\eta :XH\stackrel{\mathrm{id}\mathrm{\Delta }}{}XHH\stackrel{ev\mathrm{id}}{}𝐤H\stackrel{}{}H$$
is a polylinear map in $`^b(H)`$. Let again $`L`$ be a (Lie) pseudoalgebra and $`(M,\rho )`$ be an $`L`$-module, where $`\rho \mathrm{Lin}(\{L,M\},M)`$ in $`^{}(H)`$. Then $`\eta _H\rho \mathrm{Lin}(\{X_HL,H_HM\},H_HM)`$ is a polylinear map in $`^l(H)`$. In other words, we get a homomorphism of $`H`$-modules $`\eta _H\rho :𝒜(L)MM`$. Proposition 7.5 now implies:
###### Corollary 7.7.
The above map $`\eta _H\rho `$ provides $`M`$ with the structure of an $`𝒜(L)`$-module, and this structure is compatible with that of an $`H`$-module (cf. (7.4)).
For $`aL`$, $`xX`$, the action of $`a_xx_Ha`$ on an element $`mM`$ will be denoted by $`a_xm`$. This defines $`x`$-products $`a_xm:=a_xmM`$. When $`M=L`$ is the Lie pseudoalgebra with the adjoint action, these will be called $`x`$-brackets and denoted as $`[a_xb]`$. Then all the axioms of (Lie or associative) pseudoalgebras, representations, etc., can be reformulated in terms of the properties of the $`x`$-brackets or products — this will be done in Section 9. Although this may seem a mere tautology, it is more explicit and convenient in some cases.
Finally, let us give two more constructions.
###### Example 7.8.
The base field $`𝐤`$, with the action $`h1=\epsilon (h)`$ ($`hH`$), is a commutative associative $`H`$-differential algebra. Then for any Lie $`H`$-pseudoalgebra $`L`$, $`𝒜_𝐤L=𝐤_HL`$ is a Lie $`H`$-differential algebra. Explicitly, $`𝒜_𝐤LL/H_+L`$, where $`H_+=\{hH|\epsilon (h)=0\}`$ is the augmentation ideal. The Lie bracket in $`L/H_+L`$ is given by (cf. (7.2)):
(7.6)
$$[a\mathrm{mod}H_+L,b\mathrm{mod}H_+L]=_i\epsilon (f_i)\epsilon (g_i)e_i\mathrm{mod}H_+L,$$
if
(7.7)
$$[ab]=_i(f_ig_i)_He_i.$$
In the case when $`𝔡=𝐤`$ is $`1`$-dimensional, we recover the usual construction $`LL/L`$ that assigns a Lie algebra to any Lie conformal algebra \[K2\].
###### Remark 7.9.
Let $`Y`$ be a commutative associative $`H`$-differential algebra with a right action of $`H`$, and let $`L`$ be a Lie $`H`$-pseudoalgebra. We provide $`YL`$ with the following structure of a left $`H`$-module:
(7.8)
$$h(xa)=xh_{(1)}h_{(2)}a,hH,xY,aL.$$
Then define a Lie pseudobracket on $`YL`$ by the formula:
(7.9)
$$[(xa)(yb)]=_i\left(f_{i}^{}{}_{(1)}{}^{}g_{i}^{}{}_{(1)}{}^{}\right)_H\left((xf_{i}^{}{}_{(2)}{}^{})(yg_{i}^{}{}_{(2)}{}^{})e_i\right),$$
if $`[ab]`$ is given by (7.7). It is easy to check that (7.9) is well defined and provides $`YL`$ with the structure of a Lie $`H`$-pseudoalgebra. Moreover, $`𝒜_YL(YL)/H_+(YL)`$ as a Lie algebra (cf. Example 7.8).
In the case $`𝔡=𝐤`$, $`Y=𝐤[t,t^1]`$, $`=_t`$, the Lie $`𝐤[]`$-pseudoalgebra (= conformal algebra) $`YL`$ is known as an affinization of the conformal algebra $`L`$ \[K2\].
### 7.3. Relation to differential Lie algebras
Fix two positive integers $`N,r`$ and let $`𝒪_N=𝐤[[t_1,\mathrm{},t_N]]`$, $`=𝒪_N𝐤^r`$. A structure of a Lie algebra on $``$ is called local (and $``$ is called a local Lie algebra \[Ki\]) if the Lie bracket is given by matrix bi-differential operators. More explicitly, let $`\{e^i\}`$ be a basis of $`𝐤^r`$. Then for any $`x,y𝒪_N`$, the bracket in $``$ is given by:
(7.10)
$$[xe^i,ye^j]=_{k,l}(P_{kl}^{ij}x)(Q_{kl}^{ij}y)e^k,$$
where $`P_{kl}^{ij},Q_{kl}^{ij}`$ are differential operators with coefficients in $`𝒪_N`$. The number $`r`$ is called the rank of $``$.
A related notion is that of a differential Lie algebra \[R1\]\[R4\] (see also \[C\]). This is a Lie algebra structure on $`=Y𝐤^r`$, where $`Y`$ is any commutative associative $`H=𝐤[_1,\mathrm{},_N]`$-differential algebra, given by (7.10) for $`x,yY`$, $`P_{kl}^{ij},Q_{kl}^{ij}YH`$. One can allow a universal enveloping algebra $`H=U(𝔡)`$ ($`dim𝔡=N`$) in place of $`𝐤[_1,\mathrm{},_N]`$, cf. \[NW\].
Recall that for $`H=U(𝔡)`$, $`X=H^{}`$ is a commutative associative $`H`$-differential algebra that can be identified with $`𝒪_N`$ for $`N=dim𝔡`$. Moreover, the action of $`H`$ (and of $`XH`$) on $`X`$ is given by differential operators in this identification. Therefore a differential Lie algebra for $`Y=X`$ is the same as a local Lie algebra.
Then the results of Section 7.1 immediately imply:
###### Proposition 7.10.
Let $`L=H𝐤^r`$ be a Lie pseudoalgebra which is a free $`H`$-module of rank $`r`$. Let $`Y`$ be an $`H`$-bimodule which is a commutative associative $`H`$-differential algebra both for the left and for the right action of $`H`$ (see (2.8), (2.19)). Then $`𝒜_YLY𝐤^r`$ is a differential Lie algebra. In particular, $`𝒜(L)=𝒜_XL`$ is a local Lie algebra.
Note that the differential Lie algebras $`𝒜_YL`$ that we get are with “constant coefficients”: in (7.10) all $`P_{kl}^{ij},Q_{kl}^{ij}H`$.
### 7.4. Topology on the annihilation algebra
Now let us discuss the problem of defining a topology on $`𝒜(M)=X_HM`$ where $`M`$ is any finite $`H`$-module. Recall that $`X`$ has a decreasing filtration $`X=\mathrm{F}_1X\mathrm{F}_0X\mathrm{}`$ defined in Section 2. We can use this filtration to construct an induced filtration on $`𝒜(M)`$ as follows. Choose a finite-dimensional (over $`𝐤`$) subspace $`M_0`$ of $`M`$ which generates $`M`$ over $`H`$, and set:
(7.11)
$$\mathrm{F}_i𝒜(M)=\{x_Hm|x\mathrm{F}_iX,mM_0\}.$$
Note that since $`H`$ is cocommutative, its filtration satisfies (2.15), hence $`\mathrm{F}_iX=0`$. This implies:
(7.12)
$$\underset{i}{}\mathrm{F}_i𝒜(M)=0.$$
The filtration (7.11) will in general depend on the choice of $`M_0`$, but the topology induced by it will not, as any two such filtrations are equivalent by the next lemma.
###### Lemma 7.11.
Let $`M_0`$ and $`M_0^{}`$ be two finite-dimensional subspaces of $`M`$ generating it over $`H`$, and let $`\{\mathrm{F}_i𝒜(M)\}`$, $`\{\mathrm{F}_i^{}𝒜(M)\}`$ be the corresponding filtrations on $`𝒜(M)`$. Then there exist integers $`a,b`$ such that $`\mathrm{F}_{i+a}𝒜(M)\mathrm{F}_i^{}𝒜(M)\mathrm{F}_{i+b}𝒜(M)`$ for all values of $`i`$.
###### Proof.
Let us choose bases of $`M_0`$ and $`M_0^{}`$, and let us fix expressions of elements from the first basis as $`H`$-linear combinations of elements from the second basis. Denote by $`a`$ the highest degree of the coefficients of all these expressions. Using (2.22), we see that $`\mathrm{F}_i𝒜(M)\mathrm{F}_{ia}^{}𝒜(M)`$ for all $`i`$. Repeating the same reasoning after switching the roles of $`M_0`$ and $`M_0^{}`$, we get $`\mathrm{F}_i^{}𝒜(M)\mathrm{F}_{ib}𝒜(M)`$ for some $`b`$ and all $`i`$. ∎
###### Proposition 7.12.
Let $`H`$ be a cocommutative Hopf algebra which satisfies (2.16).
(i) If $`M`$ is a finite $`H`$-module, then $`𝒜(M)`$ is a linearly compact topological vector space when provided with the filtration (7.11). The action of $`H`$ on $`𝒜(M)`$ is continuous if we endow $`H`$ with the discrete topology.
(ii) If $`L`$ is a finite Lie $`H`$-pseudoalgebra, then its annihilation algebra $`𝒜(L)`$ is a linearly compact Lie $`H`$-differential algebra, i.e., it is a linearly compact topological vector space and both the Lie bracket and the action of $`H`$ are continuous.
A similar statement holds for representations and for associative pseudoalgebras as well.
###### Proof.
(i) The linear-compactness follows from Proposition 6.1, since (7.11) a filtration by finite-codimensional subspaces with trivial intersection and $`𝒜(M)`$ is complete with respect to this filtration. The continuity of the $`H`$-action follows from (2.22):
(7.13)
$$\mathrm{F}^iH\mathrm{F}_j𝒜(M)\mathrm{F}_{ji}𝒜(M)\text{for all}i,j.$$
(ii) It only remains to check that the Lie bracket of $`𝒜(L)`$ is continuous. Let $`L_0`$ be a finite-dimensional (over $`𝐤`$) subspace of $`L`$ which generates it over $`H`$. For $`a,bL_0`$, we can write
$$[ab]=_i(f_ig_i)_He_i$$
for some $`f_i,g_iH`$ and $`e_iL_0`$. Then the Lie bracket in $`𝒜(L)`$, for $`x,yX`$, is given by:
$$[x_Ha,y_Hb]=_i(xf_i)(yg_i)_He_i.$$
We can find a number $`p`$ such that all coefficients $`f_i,g_iH`$ occurring in pseudobrackets of any elements $`a,bL_0`$ belong to $`\mathrm{F}^pH`$. Then equations (2.21), (2.22) imply:
(7.14)
$$[\mathrm{F}_i𝒜(L),\mathrm{F}_j𝒜(L)]\mathrm{F}_{i+js}𝒜(L)\text{for all}i,j,$$
where $`s=2p`$. This shows that the Lie bracket is continuous. ∎
###### Lemma 7.13.
Let $`H=U(𝔡)\mathrm{}𝐤[\mathrm{\Gamma }]`$. Then for any nonzero $`h𝔡`$ and for every open subspace $`U`$ of $`𝒜(M)`$ there is some $`n`$ such that $`h^nU=𝒜(M)`$. In particular, each such $`h`$ acts surjectively on $`𝒜(M)`$.
###### Proof.
Follows immediately from Lemma 6.10. ∎
### 7.5. Growth of the annihilation algebra
Let $`M`$ be a finite $`H`$-module. Then any choice of a finite-dimensional subspace $`M_0`$ generating $`M`$ over $`H`$ provides $`=𝒜(M)`$ with a filtration $`_n:=\mathrm{F}_nX_HM_0`$.
###### Definition 7.14.
For a filtered vector space $`=_1_0\mathrm{}`$ we define its growth $`\mathrm{gw}`$ to be $`d`$ if the function $`ndim/_n`$ can be bounded from above and below by polynomials of degree $`d`$.
By Lemma 7.11, a different choice of $`M_0`$ would give a uniformly equivalent filtration of the same growth as $`\{_n\}`$. Hence, we can speak of the growth of $`𝒜(M)`$ independently of the choice of $`M_0`$.
###### Proposition 7.15.
Let $`H=U(𝔡)`$ be the universal enveloping algebra of a finite-dimensional Lie algebra $`𝔡`$, and $`M`$ be a finitely generated $`H`$-module. Then the growth of $`𝒜(M)`$ is equal to the dimension of $`𝔡`$.
###### Proof.
First of all, notice that we can assume $`M`$ is torsion-free, since by Lemma 7.4, $`𝒜(M)𝒜(M/\mathrm{Tor}M)`$ where $`\mathrm{Tor}M`$ is the torsion submodule of $`M`$. The proof of the proposition is then based on Lemma 2.2 and the following two lemmas.
###### Lemma 7.16.
The map $`𝒜(f):𝒜(M)𝒜(F)`$ induced by the inclusion $`f:MF`$ constructed in Lemma 2.2 is uniformly continuous, i.e., for every $`i`$ we have:
$$\mathrm{F}_{ia}𝒜(F)𝒜(f)(\mathrm{F}_i𝒜(M))\mathrm{F}_{i+b}𝒜(F),$$
where $`a`$ and $`b`$ are independent of $`i`$.
The same is true for $`𝒜(g):𝒜(F)𝒜(M)`$ where $`g`$ is the embedding $`g:hFM`$ from Lemma 2.2.
###### Proof.
Let us choose finite-dimensional vector subspaces $`F_0`$ of $`F`$ generating $`F`$ over $`H`$, and $`M_0`$ of $`M`$ generating $`M`$ over $`H`$ and containing $`hF_0`$. Let us also choose a second finite-dimensional vector subspace $`F_0^{}`$ of $`F`$ containing $`M_0`$ and generating $`F`$ over $`H`$. We will denote the filtrations induced by these subspaces by $`\{_i\}`$, $`\{_i\}`$, and $`\{_i^{}\}`$, respectively.
Up to identifying $`hF`$ with $`F`$, we have constructed injective maps $`F\stackrel{𝑔}{}M\stackrel{𝑓}{}F`$ such that the composition $`fg`$ is a multiplication by $`h`$. These maps induce maps $`𝒜(F)\stackrel{𝒜(g)}{}𝒜(M)\stackrel{𝒜(f)}{}𝒜(F)`$ which are surjective, as one can see by tensoring by $`X`$ and using that $`𝒜(T)=0`$ if $`T`$ is a torsion $`H`$-module (see Lemma 7.4).
The above maps are also continuous with respect to the common topology defined by any of the above constructed filtrations. In fact, by construction, one has
$$𝒜(g)(_i)_i\text{and}𝒜(f)(_i)_i^{}.$$
The second inclusion proves that $`𝒜(f)(_i)_{i+b}`$ for some $`b`$ independent of $`i`$, because the filtrations $`\{_i\}`$ and $`\{_i^{}\}`$ are uniformly equivalent by Lemma 7.11.
Applying $`𝒜(f)`$ to the first inclusion, we get $`𝒜(f)𝒜(g)(_i)𝒜(f)(_i)`$. On the other hand, $`𝒜(f)𝒜(g)=𝒜(fg)=h_H\mathrm{id}_F`$, and Lemma 6.10 implies that $`𝒜(f)𝒜(g)(_i)=_{ia}`$ where $`a`$ is such that $`h\mathrm{F}^aH`$ but $`h\mathrm{F}^{a1}H`$. Therefore $`_{ia}𝒜(f)(_i)`$ for all $`i`$.
A similar argument works for $`𝒜(g)`$. ∎
###### Lemma 7.17.
If $`\phi :𝒩`$ is a surjective uniformly continuous map of filtered modules, then $`\mathrm{gw}\mathrm{gw}𝒩`$.
###### Proof.
By assumption $`\phi (_i)𝒩_{i+b}`$ for all $`i`$ and some $`b`$ independent of $`i`$. The induced map $`/_i𝒩/𝒩_{i+b}`$ is a surjective map of finite-dimensional vector spaces. Hence $`\mathrm{gw}\mathrm{gw}𝒩`$. ∎
Using the above lemmas, now we can complete the proof of Proposition 7.15. We have constructed an embedding $`f:MF`$ of $`M`$ into a free $`H`$-module $`F`$, and we have shown that the induced map $`𝒜(f):𝒜(M)𝒜(F)`$ is surjective and uniformly continuous. This implies $`\mathrm{gw}𝒜(M)\mathrm{gw}𝒜(F)`$. Similarly, the inclusion $`hFM`$ gives us $`\mathrm{gw}𝒜(F)=\mathrm{gw}𝒜(hF)\mathrm{gw}𝒜(M)`$. Therefore, $`\mathrm{gw}𝒜(M)=\mathrm{gw}𝒜(F)`$. It remains to note that, for any (nonzero) free $`H`$-module $`F`$ of finite rank, one has $`\mathrm{gw}𝒜(F)=dim𝔡`$. This follows from the fact that $`\mathrm{gw}X=N=dim𝔡`$ because $`X𝒪_N`$. ∎
###### Remark 7.18.
The growth of a linearly compact Lie algebra $``$ satisfying the descending chain condition can be defined as follows. Take a fundamental subalgebra $`A`$, and build a filtration of $``$ by:
$$_0^A=A,_{i+1}^A=\{x_i^A|[x,]_i^A\},i0.$$
Taking $`A^{}=_k^A`$, we have: $`_i^A^{}=_{k+i}^A`$, hence replacing $`A`$ by $`A^{}`$ does not change the growth. Therefore, by the Chevalley principle \[G1\], the growth of this filtration does not depend on the choice of $`A`$. We will denote this common growth by $`\mathrm{gw}`$.
Notice that all simple linearly compact Lie algebras satisfy the descending chain condition, and therefore have a well defined growth which equals $`N`$ for $`W_N,S_N,H_N`$ and $`K_N`$, and $`0`$ for finite-dimensional Lie algebras.
## 8. Primitive Pseudoalgebras of Vector Fields
In this section, $`𝔡`$ will be a (finite-dimensional) Lie algebra and $`H=U(𝔡)`$ will be its universal enveloping algebra. As usual, we will identify $`𝔡`$ with its image in $`H`$. Then $`X:=H^{}`$ is the algebra of formal power series on $`𝔡^{}`$, which is isomorphic as a topological algebra to $`𝒪_N`$ for $`N=dim𝔡`$. In this section we are going to define $`H`$-pseudoalgebra analogues of the primitive linearly compact Lie algebras $`W_N`$, $`S_N`$, $`H_N`$, $`K_N`$, which will be called primitive pseudoalgebras of vector fields.
### 8.1. $`W(𝔡)`$
Let $`Y`$ be a commutative associative algebra on which $`𝔡`$ acts by derivations from the right (i.e., $`Y`$ is an $`H`$-differential algebra). One can define a left action of $`Y𝔡`$ on $`Y`$ using the right action of $`𝔡`$ on $`Y`$:
(8.1)
$$(xa)z=x(za),x,zY,a𝔡.$$
This will define a representation of $`Y𝔡`$ in $`Y`$ if the Lie bracket of $`𝔡`$ is extended to $`Y𝔡`$ by the formula
(8.2)
$$[xa,yb]=xy[a,b]x(ya)b+(xb)ya.$$
In particular, for $`Y=X=H^{}`$, this gives the Lie algebra of all vector fields on $`X`$, which is isomorphic to $`W_N`$ for $`N=dim𝔡`$.
Comparing (8.2) with (7.2), we are led to define the pseudoalgebra $`W(𝔡)=H𝔡`$ with pseudobracket
(8.3)
$$\begin{array}{cc}\hfill [(fa)(gb)]& =(fg)_H(1[a,b])\hfill \\ & (fga)_H(1b)+(fbg)_H(1a).\hfill \end{array}$$
It is easy to check that $`W(𝔡)`$ is indeed a Lie pseudoalgebra, and that the Lie algebra $`𝒜_YW(𝔡)`$ defined in Section 7 is isomorphic to $`Y𝔡`$ with bracket defined by (8.2). In a similar fashion, the module $`Y`$ over $`Y𝔡`$, defined by (8.1), leads to a structure of a $`W(𝔡)`$-module on $`H`$:
(8.4)
$$(fa)g=(fga)_H1.$$
### 8.2. Differential forms
We can think of $`X=H^{}`$ as the space of functions on $`𝔡`$, and of the elements of $`X𝔡`$ as vector fields. Then the space of $`n`$-forms ($`n=0,\mathrm{},dim𝔡`$) is
$$\mathrm{\Omega }_X^n:=\mathrm{Hom}_𝐤(^n𝔡,X)X^n𝔡^{}.$$
It is convenient to extend the elements $`\omega \mathrm{\Omega }_X^n`$ to functions from $`^n(X𝔡)`$ to $`X`$, polylinear over $`X`$:
$$\omega (x_1a_1\mathrm{}x_na_n)=x_1\mathrm{}x_n\omega (a_1\mathrm{}a_n),$$
so that
$$\mathrm{\Omega }_X^n=\mathrm{Hom}_X(^n(X𝔡),X).$$
We view $`X`$ as a left ($`X𝔡`$)-module using the right action of $`𝔡`$, see (8.1). There is a differential $`\mathrm{d}:\mathrm{\Omega }_X^n\mathrm{\Omega }_X^{n+1}`$ satisfying $`\mathrm{d}^2=0`$; this is just the usual differential for the cohomology of $`𝔡`$ with coefficients in $`X`$ where $`X`$ is viewed as a right $`𝔡`$-module:
$$\begin{array}{cc}\hfill (\mathrm{d}\omega )(a_1\mathrm{}a_n)& =\underset{i<j}{}(1)^{i+j}\omega ([a_i,a_j]a_1\mathrm{}\widehat{a}_i\mathrm{}\widehat{a}_j\mathrm{}a_n)\hfill \\ & +\underset{i}{}(1)^i\omega (a_1\mathrm{}\widehat{a}_i\mathrm{}a_n)a_i.\hfill \end{array}$$
The following analogue of the Poincaré Lemma is very useful.
###### Lemma 8.1.
The complex $`(\mathrm{\Omega }_X^{},\mathrm{d})`$ is acyclic, i.e., its $`n`$-th cohomology is trivial for $`n>0`$ and $`1`$-dimensional for $`n=0`$.
###### Proof.
It is well known that $`\mathrm{H}^n(𝔡,U(𝔡)^{})\mathrm{H}_n(𝔡,U(𝔡))^{}`$ is trivial for $`n>0`$ and $`1`$-dimensional for $`n=0`$; see e.g. \[Fu\]. ∎
For a vector field $`AX𝔡`$, we have the contraction operator $`\iota _A:\mathrm{\Omega }_X^n\mathrm{\Omega }_X^{n1}`$ given by:
$$(\iota _A\omega )(a_1\mathrm{}a_{n1})=\omega (Aa_1\mathrm{}a_{n1}).$$
We define the Lie derivative $`L_A:\mathrm{\Omega }_X^n\mathrm{\Omega }_X^n`$ by Cartan’s formula $`L_A=\mathrm{d}\iota _A+\iota _A\mathrm{d}`$. Explicitly, for $`xaX𝔡`$, we have:
(8.5)
$$\begin{array}{cc}\hfill (L_{xa}\omega )(& a_1\mathrm{}a_n)=x(\omega (a_1\mathrm{}a_n)a)\hfill \\ & +\underset{i}{}(1)^i(xa_i)\omega (aa_1\mathrm{}\widehat{a}_i\mathrm{}a_n)\hfill \\ & +\underset{i}{}(1)^ix\omega ([a,a_i]a_1\mathrm{}\widehat{a}_i\mathrm{}a_n).\hfill \end{array}$$
The Lie derivative provides each $`\mathrm{\Omega }_X^n`$ with the structure of a module over the Lie algebra of vector fields $`X𝔡`$.
For $`n=0`$, $`\mathrm{\Omega }_X^0=X`$ and this is the usual action (8.1) of $`X𝔡`$ on $`X`$. When $`n=N=dim𝔡`$, we have $`\mathrm{\Omega }_X^N=Xv_0`$ where $`v_0^N𝔡^{}`$, $`v_00`$ is a volume form. An easy calculation shows that
(8.6)
$$L_{xa}(yv_0)=\left((xy)(a+\mathrm{tr}\mathrm{ad}a)\right)v_0,x,yX,a𝔡.$$
### 8.3. Pseudoforms
The module $`\mathrm{\Omega }_X^n`$ over the Lie algebra $`X𝔡`$ leads to a module $`\mathrm{\Omega }^n(𝔡)`$ over the Lie pseudoalgebra $`W(𝔡)`$ which we now define. We let
$$\mathrm{\Omega }^n(𝔡)=H^n𝔡^{},\mathrm{\Omega }(𝔡)=_{n=0}^N\mathrm{\Omega }^n(𝔡)(N=dim𝔡).$$
The elements of $`\mathrm{\Omega }(𝔡)`$ are called pseudoforms.
$`\mathrm{\Omega }^n(𝔡)`$ is a free $`H`$-module, so that $`𝒜(\mathrm{\Omega }^n(𝔡))=X_H\mathrm{\Omega }^n(𝔡)X^n𝔡^{}=\mathrm{\Omega }_X^n`$. The action of $`W(𝔡)=H𝔡`$ on $`\mathrm{\Omega }^n(𝔡)`$ is obtained by comparing (7.2) with (8.5). To write an explicit formula, we identify $`\mathrm{\Omega }^n(𝔡)`$ with the space of linear maps from $`^n𝔡`$ to $`H`$, and $`(HH)_H\mathrm{\Omega }^n(𝔡)`$ with the space of linear maps from $`^n𝔡`$ to $`HH`$. Then for $`faW(𝔡)`$, $`w\mathrm{\Omega }^n(𝔡)`$, and $`a_i𝔡`$, we have:
(8.7)
$$\begin{array}{cc}\hfill ((fa)w)(& a_1\mathrm{}a_n)=fw(a_1\mathrm{}a_n)a\hfill \\ & +\underset{i}{}(1)^ifa_iw(aa_1\mathrm{}\widehat{a}_i\mathrm{}a_n)\hfill \\ & +\underset{i}{}(1)^ifw([a,a_i]a_1\mathrm{}\widehat{a}_i\mathrm{}a_n).\hfill \end{array}$$
When $`n=0`$, $`\mathrm{\Omega }^0(𝔡)=H`$, and we recover (8.4). In the other extreme case, when $`n=N:=dim𝔡`$, $`\mathrm{\Omega }^N(𝔡)=Hv_0`$ is again a free $`H`$-module of rank one, where $`v_0^N𝔡^{}`$, $`v_00`$ is any volume form on $`𝔡`$. We have (cf. (8.6)):
(8.8)
$$(fa)v_0=\left(f(a+\mathrm{tr}\mathrm{ad}a)1+fa\right)_Hv_0.$$
Define polylinear maps $`_\iota \mathrm{Lin}(\{W(𝔡),\mathrm{\Omega }^n(𝔡)\},\mathrm{\Omega }^{n1}(𝔡))`$ by
(8.9)
$$((fa)_\iota w)(a_1\mathrm{}a_{n1})=fw(aa_1\mathrm{}a_{n1}).$$
Also define a differential $`\mathrm{d}:\mathrm{\Omega }^n(𝔡)\mathrm{\Omega }^{n+1}(𝔡)`$ by
$$\begin{array}{cc}& \begin{array}{cc}\hfill (\mathrm{d}w)(a_1\mathrm{}a_{n+1})& =\underset{i<j}{}(1)^{i+j}w([a_i,a_j]a_1\mathrm{}\widehat{a}_i\mathrm{}\widehat{a}_j\mathrm{}a_{n+1})\hfill \\ & +\underset{i}{}(1)^iw(a_1\mathrm{}\widehat{a}_i\mathrm{}a_{n+1})a_i\text{if}n1,\hfill \end{array}\hfill \\ & (\mathrm{d}w)(a)=wa\text{if}n=0,\hfill \end{array}$$
so that $`\mathrm{d}`$ is $`H`$-linear and $`\mathrm{d}^2=0`$. For any pseudoform $`w`$ and $`xX`$, we have a differential form $`x_Hw`$ and the relation $`\mathrm{d}(x_Hw)=x_H\mathrm{d}w`$.
###### Remark 8.2.
The $`n`$-th cohomology of the complex $`(\mathrm{\Omega }^{}(𝔡),\mathrm{d})`$ is equal to $`\mathrm{H}^n(𝔡,U(𝔡))`$. This space is trivial for $`nN=dim𝔡`$ and $`1`$-dimensional for $`n=N`$. This follows from the Poincaré duality $`\mathrm{H}^n(𝔡,U(𝔡))\mathrm{H}^{Nn}(𝔡,U(𝔡)^{})`$; see e.g. \[Fu\].
We have the following analogue of Cartan’s formula for the action of $`W(𝔡)`$ on $`\mathrm{\Omega }(𝔡)`$:
(8.10)
$$\alpha w=((\mathrm{id}\mathrm{id})_H\mathrm{d})(\alpha _\iota w)+\alpha _\iota (\mathrm{d}w)(HH)_H\mathrm{\Omega }(𝔡).$$
This implies that the action of $`W(𝔡)`$ commutes with $`\mathrm{d}`$:
(8.11)
$$\alpha (\mathrm{d}w)=((\mathrm{id}\mathrm{id})_H\mathrm{d})(\alpha w).$$
We note that the maps $`\alpha _\iota `$ anticommute with each other:
(8.12)
$$\alpha _\iota (\beta _\iota w)+((\sigma \mathrm{id})_H\mathrm{id})\beta _\iota (\alpha _\iota w)=0$$
for $`\alpha ,\beta W(𝔡)`$, $`w\mathrm{\Omega }(𝔡)`$.
The wedge product on $`^{}𝔡^{}`$ can be extended to a pseudoproduct $``$ on $`\mathrm{\Omega }(𝔡)=H^{}𝔡^{}`$, so that it becomes a current pseudoalgebra. Then it is easy to check that for $`\alpha W(𝔡)`$, $`v^m𝔡^{}`$, $`w^n𝔡^{}`$, one has:
(8.13) $`\alpha (vw)`$ $`=(\alpha v)w+((\sigma \mathrm{id})_H\mathrm{id})v(\alpha w),`$
and similarly,
(8.14) $`\alpha _\iota (vw)`$ $`=(\alpha _\iota v)w+(1)^m((\sigma \mathrm{id})_H\mathrm{id})v(\alpha _\iota w).`$
This can be interpreted as saying that $`\alpha `$ and $`\alpha _\iota `$ are superderivations of $`\mathrm{\Omega }(𝔡)`$, see Example 10.10 below.
### 8.4. $`S(𝔡,\chi )`$
The divergence of a vector field $`_ix_ia_iX𝔡`$ is defined by $`\mathrm{div}(_ix_ia_i)=_ix_ia_iX`$. Then one easily checks
(8.15)
$$\mathrm{div}([A,B])=A\mathrm{div}(B)B\mathrm{div}(A),A,BX𝔡,$$
so that the divergence zero vector fields form a Lie subalgebra $`S_N`$ of $`W_N`$. Let $`\chi `$ be a trace form on $`𝔡`$, i.e., a linear functional from $`𝔡`$ to $`𝐤`$ which vanishes on $`[𝔡,𝔡]`$. Then we can define
$$\mathrm{div}^\chi (_ix_ia_i):=_ix_i(a_i+\chi (a_i)),$$
which still satisfies (8.15).
###### Remark 8.3.
Let $`\chi `$ be as above, and let $`\psi =\chi \mathrm{tr}\mathrm{ad}`$, which is again a trace form on $`𝔡`$. We can consider $`\psi `$ as an element of $`\mathrm{\Omega }_X^1=X𝔡^{}`$; then $`\mathrm{d}\psi =0`$ and by Lemma 8.1 we have $`\psi =\mathrm{d}z`$ for some $`zX`$. This means that $`\psi (a)=za`$ for all $`a𝔡`$. Let $`y=e^z`$; then $`ya=y\psi (a)`$ for $`a𝔡`$. Consider the volume form $`v=yv_0`$, where $`v_0^N𝔡^{}`$, $`v_00`$. Equation (8.6) gives
(8.16)
$$L_Av=\mathrm{div}^\chi (A)v\text{for}AX𝔡.$$
Therefore, the Lie algebra of vector fields $`A`$ with $`\mathrm{div}^\chi (A)=0`$ coincides with the Lie algebra $`S_N(v)`$ of vector fields annihilating the volume form $`v`$.
Using the notation $`\alpha _xx_H\alpha 𝒜(W(𝔡))X𝔡`$ for $`\alpha W(𝔡)`$, $`xX`$, we find for $`\alpha =ha`$:
$$\alpha _x=x_H(ha)=xh_H(1a)xhaX𝔡,$$
hence, $`\mathrm{div}^\chi (\alpha _x)=xh(a+\chi (a))`$. Define the divergence operator $`\mathrm{div}^\chi :W(𝔡)H`$ by the formula:
(8.17)
$$\mathrm{div}^\chi (_ih_ia_i)=_ih_i(a_i+\chi (a_i)).$$
Then we have:
(8.18)
$$\mathrm{div}^\chi (\alpha _x)=x\mathrm{div}^\chi \alpha \text{for}\alpha W(𝔡),xX.$$
Since $`\mathrm{div}^\chi `$ is $`H`$-linear, we can define
$$\mathrm{div}_2^\chi :(HH)_HW(𝔡)\stackrel{\mathrm{id}_H\mathrm{div}^\chi }{}(HH)_HH\stackrel{}{}HH.$$
Similarly to (8.15), one has:
(8.19)
$$\mathrm{div}_2^\chi ([\alpha \beta ])=(\mathrm{div}^\chi \alpha 1)\sigma (\beta )(1\mathrm{div}^\chi \beta )\alpha ,\alpha ,\beta W(𝔡),$$
where $`\sigma :HHHH`$ is the transposition.
Equation (8.19) implies that
(8.20)
$$S(𝔡,\chi ):=\{\alpha W(𝔡)|\mathrm{div}^\chi \alpha =0\}$$
is a subalgebra of the Lie pseudoalgebra $`W(𝔡)`$. By Eq. (8.18) and Remark 8.3, its annihilation algebra
(8.21)
$$𝒜(S(𝔡,\chi ))=\{AW_N|\mathrm{div}^\chi A=0\}S_N.$$
The rank of $`S(𝔡,\chi )`$ as an $`H`$-module is $`N1`$; however, it is free only for $`N=2`$.
###### Proposition 8.4.
$`S(𝔡,\chi )`$ is generated over $`H`$ by elements
(8.22)
$$e_{ab}:=(a+\chi (a))b(b+\chi (b))a1[a,b]\text{for}a,b𝔡.$$
These elements satisfy $`e_{ab}=e_{ba}`$ and the relations (for $`\chi =0`$):
(8.23)
$$ae_{bc}+be_{ca}+ce_{ab}=e_{[a,b],c}+e_{[b,c],a}+e_{[c,a],b}.$$
For $`\chi =0`$, their pseudobrackets are given by:
$`[e_{ab}e_{cd}]=(ad)_He_{bc}`$ $`+(bc)_He_{ad}(ac)_He_{bd}(bd)_He_{ac}`$
$`+(a1)_He_{b,[c,d]}(b1)_He_{a,[c,d]}`$
(8.24) $`(1c)_He_{d,[a,b]}+(1d)_He_{c,[a,b]}`$
$`(11)_He_{[a,b],[c,d]}.`$
For arbitrary $`\chi `$, replace everywhere in (8.23), (8.24) all $`h𝔡`$ with $`h+\chi (h)`$.
###### Remark 8.5.
Equation (8.24) implies that for $`\chi =0`$:
(8.25)
$$\begin{array}{cc}\hfill [e_{ab}e_{ab}]& =(baab)_He_{ab}\hfill \\ & +\left(1bb1\right)_He_{a,[a,b]}+\left(a11a\right)_He_{b,[a,b]}.\hfill \end{array}$$
(Again, for any $`\chi `$, replace $`a,b`$ with $`a+\chi (a),b+\chi (b)`$.) In particular, when the elements $`a,b`$ span a Lie algebra, $`He_{ab}`$ is a Lie pseudoalgebra.
In the proof of Proposition 8.4 we are going to use the following lemma.
###### Lemma 8.6.
Let $`H=U(𝔡)`$, and let $`\{_1,\mathrm{},_N\}`$ be a basis of $`𝔡`$. If elements $`h_i\mathrm{F}^dH`$ are such that $`_ih_i_i\mathrm{F}^dH`$, then there exist $`f_{ij}\mathrm{F}^{d1}H`$ such that
$$\underset{i}{}h_i_i=\underset{i,j}{}(f_{ij}1)(_i_j_j_i)mod\mathrm{F}^{d1}H𝔡.$$
###### Proof.
The proof is by induction on the number of $`h_i`$ not contained in $`\mathrm{F}^{d1}H`$, the basis of induction being trivial. Consider $`_{i=1}^nh_i_i\mathrm{F}^dH`$, with all $`h_i\mathrm{F}^{d1}H`$. We can write $`h_i=f_i_1+k_i`$ so that $`k_i\mathrm{F}^dH`$ is a linear combination of Poincaré–Birkhoff–Witt basis elements of $`H`$ not containing $`_1`$ in their expression, and $`f_i\mathrm{F}^{d1}H`$. Then:
$`{\displaystyle \underset{i=1}{\overset{n}{}}}h_i_i`$ $`=h_1_1+{\displaystyle \underset{i=2}{\overset{n}{}}}\left(f_i_1_i+k_i_i\right)`$
$`=\left(h_1+{\displaystyle \underset{i=2}{\overset{n}{}}}f_i_i\right)_1+{\displaystyle \underset{i=2}{\overset{n}{}}}k_i_i+{\displaystyle \underset{i=2}{\overset{n}{}}}f_i[_1,_i].`$
Since the third summand in the right-hand side belongs to $`\mathrm{F}^dH`$, it follows that the first and second summands lie in $`\mathrm{F}^dH`$ too. This implies: $`h_1+_{i=2}^nf_i_i\mathrm{F}^{d1}H`$. Hence
$$\underset{i=1}{\overset{n}{}}h_i_i=\underset{i=2}{\overset{n}{}}\left(f_i_1_if_i_i_1\right)+\underset{i=2}{\overset{n}{}}k_i_imod\mathrm{F}^{d1}H𝔡,$$
and we can apply the inductive assumption. ∎
###### Proof of Proposition 8.4.
First of all, it is easy to check that the elements (8.22) indeed belong to $`S(𝔡,\chi )`$. Equation (8.23) is easy, and the computation of the pseudobrackets is straightforward using (8.3), reformulated as
(8.26)
$$\begin{array}{cc}\hfill [(1a)(1b)]& =\left((a+\chi (a))1\right)_H(1b)\hfill \\ & \left(1(b+\chi (b))\right)_H(1a)(11)_He_{ab}.\hfill \end{array}$$
Now let us consider an element $`\alpha =_ih_i_iS(𝔡,\chi )`$, $`h_iH`$. We will prove that $`\alpha `$ can be expressed as $`H`$-linear combination of the above elements (8.22) by induction on the maximal degree $`d`$ of the $`h_i`$. Since $`\alpha S(𝔡,\chi )`$, then $`_ih_i(_i+\chi (_i))=0`$, hence $`_ih_i_i\mathrm{F}^dH`$.
By Lemma 8.6, we can find elements $`f_{ij}\mathrm{F}^{d1}H`$ such that
$$\alpha =\underset{i,j}{}(f_{ij}1)(_i_j_j_i)mod\mathrm{F}^{d1}H𝔡.$$
Therefore the difference
$$\alpha \underset{i,j}{}(f_{ij}1)\left((_i+\chi (_i))_j(_j+\chi (_j))_i1[_i,_j]\right)$$
still lies in $`S(𝔡,\chi )`$ and its first tensor factor terms have degree strictly less than $`d`$. By inductive assumption, we are done. ∎
###### Remark 8.7.
(i) Let, as before, $`\chi 𝔡^{}`$ be such that $`\chi ([𝔡,𝔡])=0`$. For any $`\lambda 𝐤`$, let $`V_{\lambda ,\chi }=Hv`$ be a free $`H`$-module of rank $`1`$ with the following action of $`W(𝔡)`$ on it:
(8.27)
$$\alpha v=(\lambda \mathrm{div}^\chi \alpha 1\alpha )_Hv.$$
Using (8.19), it is easy to check that this indeed defines a representation of $`W(𝔡)`$.
For $`\lambda =0`$ we get the action (8.4), while for $`\lambda =1`$, $`\chi =\mathrm{tr}\mathrm{ad}`$ we get (8.8). One can show that all representation of $`W(𝔡)`$ on a free $`H`$-module of rank $`1`$ are given by
(8.28)
$$(1a)v=\left((\lambda a+\chi ^{}(a))11a\right)_Hv,$$
where $`a𝔡`$, $`\lambda 𝐤`$ and $`\chi ^{}`$ is a trace form on $`𝔡`$. This can be rewritten as in (8.27), for $`\chi =\chi ^{}/\lambda `$ whenever $`\lambda 0`$.
(ii) More generally, let $`M`$ be any $`W(𝔡)`$-module, equipped with a compatible action of $`H=\mathrm{Cur}𝐤`$. Here $`H=\mathrm{Cur}𝐤`$ is the associative pseudoalgebra with a pseudoproduct $`fg=(fg)_H1`$, and compatibility of the actions of $`W(𝔡)`$ and $`H`$ means that
(8.29)
$$\alpha (hm)((\sigma \mathrm{id})_H\mathrm{id})h(\alpha m)=(\alpha h)m$$
for $`\alpha W(𝔡)`$, $`hH`$, $`mM`$, where $`\alpha h=(1h)\alpha _H1`$ is the action (8.4) of $`W(𝔡)`$ on $`H`$.
Then, for any $`\lambda ,\chi `$ as above,
(8.30)
$$\alpha _{\lambda ,\chi }m=\lambda (\mathrm{div}^\chi \alpha )m+\alpha m$$
is an action of $`W(𝔡)`$ on $`M`$.
### 8.5. Pseudoalgebras of rank $`1`$
All Lie pseudoalgebras that are free of rank one over $`H`$ were described by Proposition 4.1 and Lemma 4.2. The next lemma implies that all of them are subalgebras of $`W(𝔡)`$.
###### Lemma 8.8.
Let $`\alpha HH`$ be a solution of equations (4.1), (4.2). Write $`\alpha =r+s11s`$ with a skew-symmetric $`r𝔡𝔡`$ and $`s𝔡`$, as in Lemma 4.2. Consider $`e=r+1sH𝔡`$ as an element of $`W(𝔡)`$. Then $`[ee]=\alpha _He`$ in $`W(𝔡)`$.
###### Proof.
Straightforward computation, using the definition (8.3) and equations (4.3), (4.4). ∎
Let us study equations (4.3, 4.4) in more detail. We can write
(8.31)
$$r=_i(a_ib_ib_ia_i)$$
for some linearly independent $`a_i,b_i𝔡`$. Denote by $`𝔡_1`$ their linear span, and let $`𝔡_0=𝔡_1+𝐤s`$.
###### Lemma 8.9.
$`𝔡_0`$ is a Lie subalgebra of $`𝔡`$, and $`𝔡_1`$ is $`\mathrm{ad}s`$-invariant. Moreover, $`[a_i,a_j],[b_i,b_j],[a_i,b_j]`$ and $`[a_i,b_i]+s`$ belong to $`𝔡_1`$ for $`ij`$.
###### Proof.
Similar to that of Proposition 2.2.6 in \[CP\]. If $`𝔡_1=𝔡`$, there is nothing to prove. Let $`\{c_j\}`$ be elements that complement $`\{a_i,b_i\}`$ to a basis of $`𝔡`$. If $`s`$ is not in $`𝔡_1`$ we take it to be one of the $`c_j`$’s.
Write out (4.3) as
$$_i\left([a_i,s]b_i[b_i,s]a_i+a_i[b_i,s]b_i[a_i,s]\right)=0.$$
Now, if $`[a_i,s]`$ involves some $`c_j`$’s, there is no way to cancel out the terms $`c_jb_i`$. This proves that $`[s,𝔡_1]𝔡_1`$.
Similarly, (4.4) reads
$`_{i,j}`$ $`\left([a_i,a_j]b_ib_j[b_i,a_j]a_ib_j+[b_i,b_j]a_ia_j[a_i,b_j]b_ia_j+\text{cyclic}\right)`$
$`+_i\left(a_ib_isb_ia_is+\text{cyclic}\right)=0.`$
If, for example, $`[a_i,a_j]`$ involves $`c_k`$’s, then the terms $`c_kb_ib_j`$ cannot be cancelled. Therefore $`[a_i,a_j]𝔡_1`$. If $`[a_i,b_j]`$ involves $`c_k`$’s, then the terms $`c_kb_ia_j`$ can be cancelled only with terms coming from $`sr`$. This shows that $`[a_i,b_j]+\delta _{ij}s𝔡_1`$. ∎
The universal enveloping algebra $`H_0=U(𝔡_0)`$ is a Hopf subalgebra of $`H=U(𝔡)`$. Since $`\alpha H_0H_0`$, we can consider the Lie pseudoalgebra $`H_0e`$ with pseudobracket $`[ee]=\alpha _{H_0}e`$. Then our pseudoalgebra $`He`$ is a current pseudoalgebra over $`H_0e`$.
Clearly, $`𝔡_1`$ is even dimensional. There are two cases which are treated in detail in the next two subsections: when $`𝔡_0=𝔡_1`$ and when $`𝔡_0=𝔡_1𝐤s`$. They give rise to Lie pseudoalgebras $`H(𝔡,\chi ,\omega ),K(𝔡,\theta )`$ whose annihilation Lie algebras are of hamiltonian and contact type, respectively. The following theorem summarizes some of the results of Sections 4.3 and 8.58.7.
###### Theorem 8.10.
Any Lie pseudoalgebra which is free of rank one is either abelian or isomorphic to a current pseudoalgebra over one of the Lie pseudoalgebras $`H(𝔡,\chi ,\omega )`$, $`K(𝔡,\theta )`$ defined in Sections 8.6, 8.7, respectively.
### 8.6. $`H(𝔡,\chi ,\omega )`$
This is defined as a Lie $`H`$-pseudoalgebra of rank $`1`$ (see Section 8.5) corresponding to a solution $`(r,s)`$ of equations (4.3), (4.4) with a nondegenerate $`r𝔡𝔡`$ (i.e., $`𝔡_1=𝔡`$), in which case $`N=dim𝔡`$ is even. The parameters $`\chi `$ and $`\omega `$ are defined as follows.
Since $`r`$ is nondegenerate, the linear map $`𝔡^{}𝔡`$ induced by it is invertible; its inverse gives rise to a $`2`$-form $`\omega ^2𝔡^{}`$. Explicitly, if $`r=r^{ij}_i_j`$ where $`\{_i\}`$ is a basis of $`𝔡`$, then $`\omega (_i_j)=\omega _{ij}`$ is the matrix inverse to $`r^{ij}`$. We also define a $`1`$-form $`\chi :=\iota _s\omega 𝔡^{}`$.
Conversely, given a nondegenerate skew-symmetric $`2`$-form $`\omega `$ and a $`1`$-form $`\chi `$, we can define uniquely $`r𝔡𝔡`$ as the dual to $`\omega `$ and $`s𝔡`$ so that $`\chi =\iota _s\omega `$.
###### Lemma 8.11.
When $`r𝔡𝔡`$ is nondegenerate, equations (4.3), (4.4) are equivalent to the following identities for the above-defined $`\omega ,\chi :`$
(8.32) $`\mathrm{d}\omega +\chi \omega `$ $`=0,`$
(8.33) $`\mathrm{d}\chi `$ $`=0,`$
which simply mean that $`\omega `$ is a $`2`$-cocycle for $`𝔡`$ in the $`1`$-dimensional $`𝔡`$-module defined by $`\chi `$. This establishes a one-to-one correspondence between solutions $`(r,s)`$ of (4.3), (4.4) with nondegenerate $`r`$ and solutions $`(\omega ,\chi )`$ of (8.32), (8.33) with nondegenerate $`\omega `$.
###### Proof.
Let us write $`[_i,_j]=c_{ij}^k_k`$ and $`s=s^k_k`$ (summation over repeated indices). Then (4.4) is equivalent to
(8.34)
$$\left(r^{ij}r^{kl}c_{ik}^m+r^{mj}s^l\right)+\text{cyclic}=0,$$
where “cyclic” means summing over cyclic permutations of the indices $`m,j,l`$. Multiply this equation by $`\omega _{jn}\omega _{lp}\omega _{mq}`$ and sum over $`m,j,l`$. Using that $`r^{ij}\omega _{jn}=\delta _n^i`$, we get
(8.35)
$$\left(c_{np}^m\omega _{mq}+s^l\omega _{lp}\omega _{nq}\right)+\text{cyclic}=0,$$
where now the cyclic permutations are over $`n,p,q`$. This is exactly Eq. (8.32). Conversely, multiplying (8.35) by $`r^{in}r^{jp}r^{kq}`$ and summing over $`n,p,q`$, we get (8.34).
Similarly, since $`[s,𝔡_1]𝔡_1`$, we can write $`[s,_i]=_kc_i^k_k`$. Then (4.3) is equivalent to
(8.36)
$$r^{ij}c_i^k+r^{kl}c_l^j=0,$$
which after multiplying by $`\omega _{jm}\omega _{kn}`$ and summing over $`j,k`$ becomes
(8.37)
$$c_m^k\omega _{kn}+c_n^j\omega _{mj}=0,$$
or $`L_s\omega =0`$. Conversely, (8.37) gives (8.36) after multiplying by $`r^{pm}r^{qn}`$ and summing over $`m,n`$.
Now start with a solution $`(r,s)`$ of (4.3, 4.4). Above we have deduced (8.32) and $`L_s\omega =0`$. On the other hand, since $`\iota _s\chi =0`$, we have $`\iota _s(\chi \omega )=0`$, and (8.32) implies $`\iota _s\mathrm{d}\omega =0`$. Together with $`L_s\omega =0`$ this gives $`\mathrm{d}\iota _s\omega =0`$, which is (8.33).
If we start with a solution $`(\omega ,\chi )`$ of (8.32, 8.33), the above arguments can be inverted to show that $`L_s\omega =0`$, and we get (4.3, 4.4). ∎
In the basis $`\{a_i,b_i\}`$ of $`𝔡`$ we have (8.31) and $`\omega (a_ib_i)=\omega (b_ia_i)=1`$, all other values of $`\omega `$ are zero. For $`e=r+1s`$ and any $`xX`$, the element $`e_x:=x_He`$ of the annihilation algebra $`𝒜(W(𝔡))X𝔡`$ is equal to $`(xa_ib_ixb_ia_i)+xs`$, and it is easy to check that
(8.38)
$$\omega (e_xa)=x(a+\chi (a)),a𝔡.$$
Since $`\mathrm{d}\chi =0`$, Lemma 8.1 implies that $`\chi =\mathrm{d}y`$ for some $`y\mathrm{\Omega }_X^0=X`$, i.e., $`\chi (a)=ya`$. Then $`\stackrel{~}{\omega }:=e^y\omega `$ satisfies $`\stackrel{~}{\omega }(e_xa)=(xe^y)a`$ for any $`xX`$, $`a𝔡`$. This is equivalent to $`\iota _{e_x}\stackrel{~}{\omega }=\mathrm{d}(xe^y)`$. Moreover, (8.32) implies $`\mathrm{d}\stackrel{~}{\omega }=0`$. Therefore, $`L_{e_x}\stackrel{~}{\omega }=0`$, and we have the following proposition.
###### Proposition 8.12.
Let $`H(𝔡,\chi ,\omega ):=He`$ be a Lie $`H`$-pseudoalgebra of rank $`1`$ corresponding to a solution $`(r,s)`$ of equations (4.3), (4.4) with a nondegenerate $`r𝔡𝔡`$. Define the $`2`$-form $`\stackrel{~}{\omega }`$ as above. Then $`\stackrel{~}{\omega }`$ is a symplectic form, and the subalgebra $`X_HH(𝔡,\chi ,\omega )`$ of $`X_HW(𝔡)X𝔡`$ is the Lie algebra $`H_N(\stackrel{~}{\omega })`$ of vector fields annihilating $`\stackrel{~}{\omega }`$ $`(`$which is isomorphic to $`H_N)`$.
###### Proof.
It remains to show that, conversely, any vector field that preserves the form $`\stackrel{~}{\omega }`$ is equal to $`e_x`$ for some $`xX`$. Indeed, let $`AX𝔡`$ be such that $`L_A\stackrel{~}{\omega }=0`$. Since $`\mathrm{d}\stackrel{~}{\omega }=0`$ and $`\stackrel{~}{\omega }=e^y\omega `$, this is equivalent to $`\mathrm{d}(e^y\iota _A\omega )=0`$ which implies $`e^y\iota _A\omega =\mathrm{d}z`$ for some $`zX`$. In other words, $`e^y\omega (Aa)=za`$ for any $`a𝔡`$. Using $`\chi (a)=ya`$, we get $`\omega (Aa)=x(a+\chi (a))`$ for $`x=e^yz`$. This, together with (8.38), implies $`A=e_x`$ since the $`2`$-form $`\omega `$ is nondegenerate. ∎
###### Remark 8.13.
Let $`r𝔡𝔡`$ be given by (8.31), and let $`x=_i[a_i,b_i]`$, $`\varphi =\chi +\iota _x\omega =\iota _{xs}\omega `$. Then it is easy to check that $`\mathrm{div}^\varphi (r+1s)=0`$, so we have: $`H(𝔡,\chi ,\omega )S(𝔡,\varphi )`$.
###### Example 8.14.
Let the Lie algebra $`𝔡`$ be $`2`$-dimensional with basis $`\{a,b\}`$ and commutation relations $`[a,b]=\lambda b`$. Then up to multiplication by a scalar, all nondegenerate solutions $`(r,s)`$ of (4.3) are given by: $`r=abba`$, any $`s`$ in case $`\lambda =0`$, and by the same $`r`$, and $`s𝐤b`$ when $`\lambda 0`$. It is immediate to see that in both cases $`s`$ can be written as $`\varphi (a)b+\varphi (b)a+[a,b]`$ for some trace form $`\varphi 𝔡^{}`$. Then $`r1s=e_{ab}`$ is a free generator of $`S(𝔡,\varphi )`$, since $`dim𝔡=2`$ (see Proposition 8.4). This shows that the above pairs $`(r,s)`$ also satisfy (4.4). We have: $`H(𝔡,\chi ,\omega )=S(𝔡,\varphi )`$, where $`\chi =\iota _s\omega =\varphi +\mathrm{tr}\mathrm{ad}`$. (Note that $`\mathrm{tr}\mathrm{ad}=\iota _x\omega `$ for $`x=[a,b]=\lambda b`$.)
###### Example 8.15.
When $`𝔡`$ is abelian of dimension $`N=2n>2`$, then (8.32) and (8.33) become $`\chi \omega =0`$, hence $`\chi =0`$ and $`\omega `$ is any nondegenerate skew-symmetric $`2`$-form. In this case all solutions of (4.3), (4.4) are: $`s=0`$ and $`r`$ given by (8.31) in some basis $`\{a_i,b_i\}`$ of $`𝔡`$.
###### Example 8.16.
When the Lie algebra $`𝔡`$ is simple, there are no solutions $`(\omega ,\chi )`$ of (8.32), (8.33) with a nondegenerate $`\omega `$. Indeed, since $`[𝔡,𝔡]=𝔡`$, we have $`\chi =0`$, and $`\omega `$ is a $`2`$-cocycle: $`\mathrm{d}\omega =0`$. Any $`2`$-cocycle $`\omega ^2𝔡^{}`$ for a simple Lie algebra $`𝔡`$ is degenerate, since $`\omega =\mathrm{d}\alpha `$ for some $`\alpha 𝔡^{}`$ and the stabilizer $`𝔡_\alpha `$ of $`\alpha `$ is always non-zero.
### 8.7. $`K(𝔡,\theta )`$
This is defined as a Lie $`H`$-pseudoalgebra of rank $`1`$ (see Section 8.5) corresponding to a solution $`(r,s)`$ of equations (4.3), (4.4) with $`𝔡=𝔡_1𝐤s`$ and nondegenerate $`r𝔡_1𝔡_1`$; in this case $`N=dim𝔡`$ is odd. The parameter $`\theta `$ is defined below.
Let $`\{_i\}`$ be a basis of $`𝔡_1`$, and $`r=r^{ij}_i_j`$. As before, we define a $`2`$-form $`\omega `$ on $`𝔡_1`$ by $`\omega (_i_j)=\omega _{ij}`$, where $`(\omega _{ij})`$ is the matrix inverse to $`(r^{ij})`$. Let us write $`[_i,_j]=c_{ij}^k_k+c_{ij}s`$ and $`[s,_j]=c_j^k_k`$. Then we have:
###### Lemma 8.17.
With the above notation, equations (4.3), (4.4) are equivalent to the following identities:
(8.39) $`\mathrm{d}\omega `$ $`=0\text{on }^3𝔡_1,`$
(8.40) $`c_{ij}`$ $`=\omega _{ij},`$
(8.41) $`L_s\omega `$ $`=0.`$
If we extend $`\omega `$ to a $`2`$-form on $`𝔡`$ by defining $`\iota _s\omega =0`$, then $`\omega `$ is closed: $`\mathrm{d}\omega =0`$.
###### Proof.
The proof is very similar to that of Lemma 8.11. There we showed that $`L_s\omega =0`$ is equivalent to (4.3), and the same argument applies here. Similarly, (4.4) is equivalent to (8.39, 8.40). Now if $`\iota _s\omega =0`$, then $`L_s\omega =0`$ implies $`\iota _s\mathrm{d}\omega =0`$, which together with (8.39) leads to $`\mathrm{d}\omega =0`$. ∎
Let $`\omega `$ be extended to a $`2`$-form on $`𝔡`$ by defining $`\iota _s\omega =0`$, so that $`\mathrm{d}\omega =0`$. We define a $`1`$-form $`\theta 𝔡^{}`$ by $`\theta (s):=1`$, $`\theta |_{𝔡_1}:=0`$. Then we have $`\mathrm{d}\theta =\omega `$; indeed:
$`(\mathrm{d}\theta )(_i_j)`$ $`=\theta ([_i,_j])=c_{ij}=\omega _{ij}=\omega (_i_j),`$
$`(\mathrm{d}\theta )(s_j)`$ $`=\theta ([s,_j])=0=\omega (s_j),`$
using (8.40) and the fact that $`[s,𝔡_1]𝔡_1`$.
###### Lemma 8.18.
There is a one-to-one correspondence between contact forms $`\theta `$, i.e. $`1`$-forms $`\theta 𝔡^{}`$ such that $`\theta (\mathrm{d}\theta )^{(N1)/2}0`$ $`(N=dim𝔡)`$, and solutions $`(r,s)`$ of (4.3), (4.4) with $`𝔡=𝔡_1𝐤s`$ and nondegenerate $`r𝔡_1𝔡_1`$.
###### Proof.
Given $`(r,s)`$, above we have defined the $`1`$-form $`\theta `$ such that $`\theta (s)=1`$, $`\theta |_{𝔡_1}=0`$ and $`\mathrm{d}\theta =\omega `$. Since $`\omega ^2𝔡_1^{}`$ is nondegenerate, we have $`\theta \omega ^{(N1)/2}0`$. Conversely, starting with a contact $`1`$-form $`\theta 𝔡^{}`$, we can define $`s`$ and $`\omega `$ satisfying (8.39)–(8.41). ∎
###### Example 8.19.
When $`𝔡`$ is the Heisenberg Lie algebra with a basis $`\{a_i,b_i,c\}`$ and the only nonzero commutation relations $`[a_i,b_i]=c`$ ($`1in`$, $`N=2n+1`$), then
$$r=_{i=1}^n(a_ib_ib_ia_i),s=c$$
is a solution of (4.3), (4.4).
###### Example 8.20.
When $`𝔡`$ is abelian and $`dim𝔡=2n+1>1`$, then equations (4.3), (4.4) have no solutions $`(r,s)`$ with $`𝔡=𝔡_1𝐤s`$ and a nondegenerate $`r𝔡_1𝔡_1`$, because $`\mathrm{d}\theta =0`$ and therefore there are no contact forms.
###### Example 8.21.
When the Lie algebra $`𝔡`$ is simple, a solution $`(r,s)`$ of (4.3), (4.4) with $`𝔡=𝔡_1𝐤s`$ and a nondegenerate $`r𝔡_1𝔡_1`$ exists iff $`𝔡=𝔰𝔩_2`$, and it is as follows:
$$r=ef:=effe,s=h.$$
Only $`𝔡=𝔰𝔩_2`$ is possible since the dimension of the stabilizer of $`\theta `$ equals $`1`$.
Now let us compute $`L_{e_x}\theta `$. Recall that, as in Section 8.6, for any $`xX`$ we identify $`e_x:=x_He`$ with $`(xa_ib_ixb_ia_i)+xs`$. Similarly to (8.38), it is easy to see that $`\omega (e_xa)=xa`$ for $`a𝔡_1`$ (in this case $`\chi =\iota _s\omega :=0`$). On the other hand, $`\iota _{e_x}\theta =\theta (e_x)=x`$, and hence $`(\mathrm{d}\iota _{e_x}\theta )(a)=(\mathrm{d}x)(a)=xa`$ for any $`a𝔡`$. Therefore $`(L_{e_x}\theta )(a)=0`$ for $`a𝔡_1`$, and $`(L_{e_x}\theta )(s)=xs`$. In other words,
$$L_{e_x}\theta =(xs)\theta ,$$
and we have the following proposition.
###### Proposition 8.22.
Let $`K(𝔡,\theta ):=He`$ be a Lie $`H`$-pseudoalgebra of rank $`1`$ corresponding to a solution $`(r,s)`$ of equations (4.3), (4.4) with $`𝔡=𝔡_1𝐤s`$ and a nondegenerate $`r𝔡_1𝔡_1`$, where the $`1`$-form $`\theta 𝔡^{}`$ is defined by $`\theta (s)=1`$, $`\theta |_{𝔡_1}=0`$. Then $`\theta `$ is a contact form, and the subalgebra $`X_HK(𝔡,\theta )`$ of $`X_HW(𝔡)X𝔡`$ is the Lie algebra $`K_N(\theta )`$ of vector fields that preserve $`\theta `$ up to a multiplication by a function $`(`$which is isomorphic to $`K_N)`$.
###### Proof.
It remains to show that, conversely, any vector field from $`K_N(\theta )`$ is equal to $`e_x`$ for some $`xX`$. Indeed, let $`AX𝔡`$ be such that $`L_A\theta =f\theta `$ for some $`fX`$. Let us write $`A=_i(x_ia_i+y_ib_i)+xs`$ for some $`x_i,y_i,xX`$. Then $`\omega (Aa_i)=y_i`$ and $`\omega (Ab_i)=x_i`$, while $`\theta (A)=x`$. Therefore $`(L_A\theta )(a)=\omega (Aa)+xa`$, which implies $`y_i+xa_i=0`$, $`x_i+xb_i=0`$, and $`xs=f`$. ∎
###### Remark 8.23.
To any $`H`$-type Lie pseudoalgebra, i.e., to any triple $`(𝔡,\omega ,\chi )`$ where $`𝔡`$ is a finite-dimensional Lie algebra, $`\omega ^2𝔡^{}`$ is a non-degenerate $`2`$-form and $`\chi 𝔡^{}`$ satisfying (8.32) and (8.33), we can associate a $`K`$-type Lie pseudoalgebra as follows. Set on the vector space $`𝔡^{}=𝔡𝐤c`$ the Lie bracket $`[,]^{}`$ defined as:
$$[g,h]^{}=[g,h]+\omega (g,h)c,[g,c]^{}=\chi (g)c,$$
for $`g,h𝔡`$. Then $`c+s𝔡^{}`$ stabilizes $`𝔡`$, where $`s𝔡`$ is the unique element such that $`\chi =\iota _s\omega `$; indeed,
$$[g,s+c]^{}=[g,s]+\omega (g,s)c+\chi (g)c=[g,s]𝔡.$$
Define $`\theta (𝔡^{})^{}`$ as the unique element restricting to $`0`$ on $`𝔡`$ such that $`\theta (c)=1`$.
Note that not all $`K`$-type data are obtained in this way, since the Lie algebra $`𝔡^{}`$ just constructed always has a one dimensional ideal $`𝐤c`$, and this fails in Example 8.21.
### 8.8. Annihilation algebras of pseudoalgebras of vector fields
To conclude this section, we determine the annihilation algebras of the primitive pseudoalgebras of vector fields defined above, and of current pseudoalgebras over them.
###### Theorem 8.24.
(i) If $`L`$ is one of the Lie $`H=U(𝔡)`$-pseudoalgebras $`W(𝔡)`$, $`S(𝔡,\chi )`$, $`H(𝔡,\chi ,\omega )`$ or $`K(𝔡,\theta )`$, then its annihilation algebra $`𝒜(L)`$ is isomorphic to $`W_N`$, $`S_N`$, $`P_N`$ or $`K_N`$, respectively.
(ii) If $`L=\mathrm{Cur}L^{}`$ is a current pseudoalgebra over the Lie $`H^{}`$-pseudoalgebra $`L^{}`$, then its annihilation algebra $`𝒜(L)`$ is isomorphic to a current Lie algebra $`𝒪_r\widehat{}𝒜(L^{})`$ over $`𝒜(L^{})`$, where $`H^{}=U(𝔡^{})`$ and $`𝔡^{}`$ is a codimension $`r`$ subalgebra of $`𝔡`$.
###### Proof.
(i) We have seen in Section 8.1 that $`𝒜(W(𝔡))W_N`$. Let $`L`$ be one of the pseudoalgebras $`S(𝔡,\chi )`$, $`H(𝔡,\chi ,\omega )`$ or $`K(𝔡,\theta )`$ and $`i`$ be its natural embedding in $`W(𝔡)`$. We have shown in Sections 8.4, 8.6 and 8.7 that in this case the image of $`𝒜(L)`$ in $`W_N`$ under $`𝒜(i)`$ is $`S_N`$, $`H_N`$ or $`K_N`$, respectively.
Lemma 11.14 below implies that $`𝒜(L)`$ is a central extension of its image in $`W_N`$. Moreover, since $`L`$ is simple, it is equal to its derived subalgebra, and therefore $`𝒜(L)`$ is equal to its derived subalgebra (see Section 13.1). Hence, $`𝒜(L)`$ is an irreducible central extension of its image in $`W_N`$.
Now Proposition 6.12(iii) implies that $`𝒜(L)S_N`$, $`K_N`$ in the cases $`L=S(𝔡,\chi )`$, $`K(𝔡,\theta )`$ respectively, and $`𝒜(L)`$ is a quotient of $`P_N`$ in the case $`L=H(𝔡,\chi ,\omega )`$. However, since $`L=H(𝔡,\chi ,\omega )`$ is a free $`H`$-module of rank one, $`𝒜(L)`$ is isomorphic to $`X`$ as a topological $`H`$-module. Therefore, $`𝒜(L)P_N`$.
(ii) Note that $`X=H^{}`$ maps surjectively to $`X^{}=(H^{})^{}`$ with kernel isomorphic to $`𝒪_r`$. Moreover $`X𝒪_N`$, $`X^{}𝒪_N^{}`$ $`(N^{}=Nr)`$, and hence $`X𝒪_r\widehat{}X^{}`$. We have: $`𝒜(L^{}):=X^{}_H^{}L^{}`$ and $`𝒜(L):=X_HL=X_H(H_H^{}L^{})X_H^{}L^{}(𝒪_r\widehat{}X^{})_H^{}L^{}𝒪_r\widehat{}(X^{}_H^{}L^{})`$. ∎
###### Remark 8.25.
Let us assume that the base field $`𝐤=`$, and let $`L`$ be as in Theorem 8.24(i). Then the action of $`𝔡`$ on $`𝒜(L)`$ can be constructed via the embedding of $`𝔡`$ in $`W_N`$ as follows.
(i) Any $`N`$-dimensional Lie algebra $`𝔡`$ can be embedded in $`W_N`$: every $`a𝔡`$ defines a left-invariant vector field on the connected simply-connected Lie group $`D`$ with Lie algebra $`𝔡`$, and we take the corresponding formal vector field in the formal neighborhood of the identity element. (See also Proposition 6.9.)
(ii) If we have a homomorphism of Lie algebras $`\chi :𝔡`$, it defines a homomorphism $`\stackrel{~}{\chi }`$ of $`D`$ to $`^\times `$. Consider a volume form $`v`$ on $`D`$ defined, up to a constant factor, by the property $`gv_0=\stackrel{~}{\chi }(g)v_0`$, $`gD`$, where $`v_0`$ is the value of $`v`$ at the identity element. Then we get an embedding of $`𝔡`$ in $`CS_N(v)=\mathrm{Der}S_N(v)\mathrm{Der}S_N`$.
(iii) Given $`\chi `$ and $`\omega ^2𝔡^{}`$ such that $`\mathrm{d}\omega +\chi \omega =0`$, consider a $`2`$-form $`s`$ on $`D`$ whose value at the identity element is $`\omega `$ and such that $`gs=\stackrel{~}{\chi }(g)s`$, $`gD`$. Then $`s`$ is a symplectic form on $`D`$, and we get an embedding of $`𝔡`$ in $`CH_N(s)=\mathrm{Der}H_N(s)\mathrm{Der}P_N`$.
(iv) Given a contact form $`\theta 𝔡^{}`$, consider the left-invariant $`1`$-form $`c`$ on $`D`$ with the value $`\theta `$ at the identity element. Then $`c`$ is a contact form on $`D`$, and we get an embedding of $`𝔡`$ in $`K_N(c)K_N`$.
## 9. $`H`$-Conformal Algebras
The goal of this section is to reformulate the definition of a Lie (or associative) $`H`$-pseudoalgebra in terms of the properties of the $`x`$-brackets (or products) introduced in Section 7.2. The resulting notion, equivalent to that of an $`H`$-pseudoalgebra, will be called an $`H`$-conformal algebra.
Let us start by deriving explicit formulas for the $`x`$-brackets. We will use the notation of Section 7.2. Let $`(L,\beta )`$ be a Lie $`H`$-pseudoalgebra with a pseudobracket
(9.1)
$$[ab]\beta (ab)=_i(f_ig_i)_He_i.$$
Then for $`xX`$, $`hH`$ we have $`\eta (xh)=x,h_{(1)}h_{(2)}`$ (see (7.5)), and
$`(\eta _H\beta )\left((x_Ha)(h_Hb)\right)`$ $`=_i\eta (xf_ihg_i)_He_i`$
$`=_ixf_i,(hg_i)_{(1)}(hg_i)_{(2)}_He_i.`$
Taking $`h=1`$, we get the following expression for the $`x`$-bracket in $`L`$:
(9.2)
$$[a_xb]=_iS(x),f_ig_{i}^{}{}_{(1)}{}^{}g_{i}^{}{}_{(2)}{}^{}e_i,\text{if}[ab]=_i(f_ig_i)_He_i.$$
Here we can recognize the Fourier transform $``$, defined by (2.33):
$$(fg)=fg_{(1)}g_{(2)}.$$
The identity (2.35):
$$fg=(fg_{(1)}1)\mathrm{\Delta }(g_{(2)}),$$
implies
(9.3)
$$[ab]=_i(f_ig_{i}^{}{}_{(1)}{}^{}1)_Hg_{i}^{}{}_{(2)}{}^{}e_i.$$
Hence $`[ab]`$ can be written uniquely in the form $`_i(h_i1)_Hc_i`$, where $`\{h_i\}`$ is a fixed $`𝐤`$-basis of $`H`$ (cf. Lemma 2.5).
We introduce another bracket $`[a,b]HL`$ as the Fourier transform of $`[ab]`$:
(9.4)
$$[a,b]=_i(f_ig_i)(1e_i)=_if_ig_{i}^{}{}_{(1)}{}^{}g_{i}^{}{}_{(2)}{}^{}e_i.$$
In other words,
(9.5)
$$[a,b]=_ih_ic_i\text{if}[ab]=_i(h_i1)_Hc_i.$$
Then we have:
(9.6)
$$[a_xb]=(S(x),\mathrm{id})[a,b]=_iS(x),h_ic_i.$$
Using properties (2.38)–(2.41) of the Fourier transform, it is straightforward to derive the properties of the bracket (9.5). Then the definition of a Lie pseudoalgebra can be equivalently reformulated as follows.
###### Definition 9.1.
A Lie $`H`$-conformal algebra is a left $`H`$-module $`L`$ equipped with a bracket $`[,]:LLHL`$, satisfying the following properties ($`a,b,cL`$, $`hH`$):
(9.7) $`[ha,b]`$ $`=(h1)[a,b],`$
(9.8) $`[a,hb]`$ $`=(1h_{(2)})[a,b](h_{(1)}1).`$
If $`[a,b]`$ is given by (9.5), then
(9.9)
$$[b,a]=_ih_{i}^{}{}_{(1)}{}^{}h_{i}^{}{}_{(2)}{}^{}c_i.$$
(9.10)
$$[a,[b,c]](\sigma \mathrm{id})[b,[a,c]]=(^1\mathrm{id})[[a,b],c]$$
in $`HHL`$, where $`\sigma :HHHH`$ is the permutation $`\sigma (fg)=gf`$, and
(9.11) $`[a,[b,c]]`$ $`=(\sigma \mathrm{id})(\mathrm{id}[a,])[b,c],`$
(9.12) $`[[a,b],c]`$ $`=(\mathrm{id}[,c])[a,b].`$
###### Examples 9.2.
(i) For the current Lie pseudoalgebra $`\mathrm{Cur}𝔤=H𝔤`$ with the pseudobracket (4.2), the bracket (9.5) is given by:
$$[fa,gb]=fg_{(1)}(g_{(2)}[a,b]).$$
(ii) For the Lie pseudoalgebra $`W(𝔡)=H𝔡`$ with pseudobracket defined by (8.3), the bracket (9.5) is given by:
$$[1a,1b]=1(1[a,b])+a(1b)+b(1a)1(ab).$$
One can also reformulate Definition 9.1 in terms of the $`x`$-brackets (9.6).
###### Definition 9.3.
A Lie $`H`$-conformal algebra is a left $`H`$-module $`L`$ equipped with $`x`$-brackets $`[a_xb]L`$ for $`a,bL`$, $`xX`$, satisfying the following properties:
(9.13)
$$\mathrm{codim}\{xX|[a_xb]=0\}<\mathrm{}\text{for any }a,bL.$$
Equivalently, for any basis $`\{x_i\}`$ of $`X`$,
(9.14)
$$[a_{x_i}b]0\text{for only a finite number of }i\text{.}$$
(9.15) $`[ha_xb]`$ $`=[a_{xh}b],`$
(9.16) $`[a_xhb]`$ $`=h_{(2)}[a_{h_{(1)}x}b].`$
Choose dual bases $`\{h_i\}`$, $`\{x_i\}`$ in $`H`$ and $`X`$. Then:
(9.17)
$$[a_xb]=_ix,h_{i}^{}{}_{(1)}{}^{}h_{i}^{}{}_{(2)}{}^{}[b_{x_i}a].$$
(9.18)
$$[a_x[b_yc]][b_y[a_xc]]=[[a_{x_{(2)}}b]_{yx_{(1)}}c].$$
Lemma 2.4 implies that (9.18) can be rewritten as follows:
(9.19)
$$[a_x[b_yc]][b_y[a_xc]]=_i[[a_{x_i}b]_{y(xS(h_i))}c].$$
In particular, the right-hand side of (9.18) is well defined: the sum is finite because of (9.14).
The definitions of an associative $`H`$-conformal algebra or of representations of $`H`$-conformal algebras are obvious modifications of the above. For example, in terms of $`x`$-products, the associativity looks as follows (cf. (9.18)):
(9.20)
$$a_x(b_yc)=(a_{x_{(2)}}b)_{yx_{(1)}}c.$$
The same argument as the one used for $``$ shows that the map $`xyx_{(2)}yx_{(1)}`$ has an inverse given by $`xyx_{(2)}yx_{(1)}`$. Therefore, (9.20) is equivalent to the following equation:
(9.21)
$$a_{x_{(2)}}(b_{yx_{(1)}}c)=(a_xb)_yc.$$
Note that when considering associative $`H`$-conformal algebras, $`H`$ need not be cocommutative, so $`X`$ may be noncommutative.
We also note that there is a simple relationship between the $`x`$-bracket (9.6) of a Lie $`H`$-conformal algebra (or, equivalently, pseudoalgebra) $`L`$ and the commutator in its annihilation Lie algebra $`𝒜(L)`$ defined in Section 7. Let $`\{h_i\}`$, $`\{x_i\}`$ again be dual bases in $`H`$ and $`X`$. Then in (9.5) one has $`c_i=[a_{S^1(x_i)}b]`$; therefore
(9.22)
$$[a,b]=_iS(h_i)[a_{x_i}b]\text{and}[ab]=(S(h_i)1)_H[a_{x_i}b].$$
Recall that we denote the element $`x_Ha`$ of $`𝒜(L):=X_HL`$ by $`a_x`$. Then the definition (7.2) and (9.22) imply
(9.23)
$$[a_x,b_y]=_i[a_{x_i}b]_{(xS(h_i))y}=[a_{x_{(2)}}b]_{x_{(1)}y},$$
using (2.32). This is also equivalent to:
(9.24)
$$[a_xb]_y=[a_{x_{(2)}},b_{x_{(1)}y}]=_i[a_{x_i},b_{(h_iS(x))y}].$$
Comparing these formulas with Eq. (9.18), we obtain the following important result.
###### Proposition 9.4.
Any module $`M`$ over a Lie pseudoalgebra $`L`$ has a natural structure of an $`𝒜(L)`$-module, given by $`(x_Ha)m=a_xm`$, where
(9.25)
$$a_xm=_iS(x),f_ig_{i}^{}{}_{(1)}{}^{}g_{i}^{}{}_{(2)}{}^{}v_i,\text{if}am=_i(f_ig_i)_Hv_i$$
for $`aL,xX,mM`$. This action is compatible with the action of $`H`$ $`(`$see (7.4)$`)`$ and satisfies the locality condition:
(9.26)
$$\mathrm{codim}\{xX|a_xm=0\}<\mathrm{},aL,mM,$$
or equivalently, for any basis $`\{x_i\}`$ of $`X`$,
(9.27)
$$a_{x_i}m0\text{for only a finite number of }i\text{.}$$
(The above conditions on $`M`$ mean that, when endowed with the discrete topology, $`M`$ is a topological $`𝒜(L)`$-module in the category $`^l(H)`$.)
Conversely, any $`𝒜(L)`$-module $`M`$ satisfying the above conditions has a natural structure of an $`L`$-module, given by:
(9.28)
$$am=_i(S(h_i)1)_Ha_{x_i}m,$$
where $`\{h_i\}`$, $`\{x_i\}`$ are dual bases in $`H`$ and $`X`$, and we use the notation $`a_xx_Ha`$.
This proposition provides the main tool for constructing modules over Lie pseudoalgebras. Of course, there is an analogous result in the case of associative algebras as well.
Finally, let us give two more expressions for the bracket in $`𝒜(L)`$ which will be useful later. Recall that, by Proposition 9.4, we have an action of $`𝒜(L)`$ on $`L`$ given by $`a_xb=[a_xb]`$. Recall also that the action of $`H`$ on $`𝒜(L)`$ is defined by $`h(a_x)=a_{hx}`$. Let $`\alpha 𝒜(L)`$, $`bL`$, $`yX`$. Then:
(9.29) $`(\alpha b)_y`$ $`=_i[h_i\alpha ,b_{x_iy}],`$
(9.30) $`[\alpha ,b_y]`$ $`=_i\left((S(h_i)\alpha )b\right)_{x_iy}.`$
Note that the infinite sums on the right-hand sides make sense since they are convergent in the complete topology of $`𝒜(L)`$. It is enough to prove both statements for $`\alpha `$ of the form $`a_x=x_Ha`$ since such elements span $`𝒜(L)`$. Equation (9.30) then follows from (9.23) and (2.32). Analogously, (9.29) derives from (9.24) by noticing that $`x_{(1)}x_{(2)}=_ix_ih_ix`$.
## 10. $`H`$-Pseudolinear Algebra
The definition of a module over a pseudoalgebra motivates the following definition of a pseudolinear map.
###### Definition 10.1.
Let $`V`$ and $`W`$ be two $`H`$-modules. An $`H`$-pseudolinear map from $`V`$ to $`W`$ is a $`𝐤`$-linear map $`\varphi :V(HH)_HW`$ such that
(10.1) $`\varphi (hv)`$ $`=((1h)_H1)\varphi (v),hH,vV.`$
We denote the space of all such $`\varphi `$ by $`\mathrm{Chom}(V,W)`$. We define a left action of $`H`$ on $`\mathrm{Chom}(V,W)`$ by:
(10.2) $`(h\varphi )(v)`$ $`=((h1)_H1)\varphi (v).`$
When $`V=W`$, we set $`\mathrm{Cend}V=\mathrm{Chom}(V,V)`$.
###### Example 10.2.
Let $`A`$ be an $`H`$-pseudoalgebra and $`V`$ be an $`A`$-module. Then for any $`aA`$ the map $`m_a:V(HH)_HV`$ defined by $`m_a(v)=av`$ is an $`H`$-pseudolinear map. Moreover, we have $`hm_a=m_{ha}`$ for $`hH`$.
Consider the map $`\rho :\mathrm{Chom}(V,W)V(HH)_HW`$ given by $`\rho (\varphi v)=\varphi (v)`$. By definition it is $`H`$-bilinear, so it is a polylinear map in $`^{}(H)`$. We will also use the notation $`\varphi v:=\varphi (v)`$ and consider this as a pseudoproduct (or rather action, see Proposition 10.5 below).
The corresponding $`x`$-products are called Fourier coefficients of $`\varphi `$ and are given by a formula analogous to (9.2):
(10.3)
$$\varphi _xv=_iS(x),f_ig_{i}^{}{}_{(1)}{}^{}g_{i}^{}{}_{(2)}{}^{}w_i,\text{if}\varphi (v)=_i(f_ig_i)_Hw_i.$$
They satisfy a locality relation and an $`H`$-sesqui-linearity relation similar to (9.13) and (9.16):
(10.4)
$$\mathrm{codim}\{xX|\varphi _xv=0\}<\mathrm{}\text{for any}vV,$$
(10.5)
$$\varphi _x(hv)=h_{(2)}(\varphi _{h_{(1)}x}v).$$
Conversely, any collection of maps $`\varphi _x\mathrm{Hom}(V,W)`$, $`xX`$, satisfying relations (10.4), (10.5) comes from an $`H`$-pseudolinear map $`\varphi \mathrm{Chom}(V,W)`$. Explicitly (cf. (9.28)):
(10.6)
$$\varphi (v)=_i(S(h_i)1)_H\varphi _{x_i}v,$$
where $`\{h_i\}`$, $`\{x_i\}`$ are dual bases in $`H`$ and $`X`$.
###### Remark 10.3.
It follows from (10.5) that for $`\varphi \mathrm{Chom}(V,W)`$, the map $`\varphi _1:VW`$ is $`H`$-linear, where $`1X`$ is the unit element. This establishes an isomorphism $`\mathrm{Hom}_H(V,W)𝐤_H\mathrm{Chom}(V,W)\mathrm{Chom}(V,W)/H_+\mathrm{Chom}(V,W)`$, where $`H_+=\{hH|\epsilon (h)=0\}`$ is the augmentation ideal.
###### Lemma 10.4.
Let $`U,V,W`$ be three $`H`$-modules, and assume that $`U`$ is finite. Then there is a unique polylinear map
$$\mu \mathrm{Lin}(\{\mathrm{Chom}(V,W),\mathrm{Chom}(U,V)\},\mathrm{Chom}(U,W))$$
in $`^{}(H)`$, denoted as $`\mu (\varphi \psi )=\varphi \psi `$, such that
(10.7)
$$(\varphi \psi )u=\varphi (\psi u)$$
in $`H^3_HW`$ for $`\varphi \mathrm{Chom}(V,W)`$, $`\psi \mathrm{Chom}(U,V)`$, $`uU`$.
###### Proof.
We define $`\varphi \psi `$ in terms of its Fourier coefficients — the $`x`$-products $`\varphi _x\psi `$. We have already seen, when we discussed associativity, that (10.7) is equivalent to the following equation (cf. (9.20)):
$$\varphi _x(\psi _yu)=(\varphi _{x_{(2)}}\psi )_{yx_{(1)}}u.$$
This can be inverted to find (cf. (9.21)):
$$(\varphi _x\psi )_yu=\varphi _{x_{(2)}}(\psi _{yx_{(1)}}u)=_i\varphi _{x_i}(\psi _{y(h_iS(x))}u).$$
The $`H`$-sesqui-linearity properties of $`(\varphi _x\psi )_yu`$ with respect to $`x`$ and $`y`$ are easy to check by a direct calculation. By properties (2.21), (2.28), (2.29) of the filtration $`\{\mathrm{F}_nX\}`$, and locality of $`\psi `$, it follows that for each fixed $`xX`$, $`uU`$ there is an $`n`$ such that $`(\varphi _x\psi )_yu=0`$ for $`y\mathrm{F}_nX`$. Therefore, for each $`xX`$ we have defined $`\varphi _x\psi \mathrm{Chom}(U,W)`$.
In order that $`\varphi \psi `$ be well defined, we need to check that $`\varphi _x\psi `$ satisfies locality, i.e., that $`\varphi _x\psi =0`$ when $`x\mathrm{F}_nX`$ with $`n0`$. By the locality of $`\varphi `$ and $`\psi `$, for each $`uU`$ there is an $`n`$ such that $`(\varphi _x\psi )_yu=0`$ for $`x\mathrm{F}_nX`$ and all $`yX`$. Since $`U`$ is finite, we can choose an $`n`$ that works for all $`u`$ belonging to a set of generators of $`U`$ over $`H`$. Now the $`H`$-sesqui-linearity of $`(\varphi _x\psi )_yu`$ with respect to $`y`$ (for fixed $`x`$) implies that $`(\varphi _x\psi )_yu=0`$ for all $`y`$ and $`u`$. Hence $`\varphi _x\psi =0`$ for $`x\mathrm{F}_nX`$. ∎
Specifying to the case $`U=V=W`$, we obtain a pseudoproduct $`\mu `$ on $`\mathrm{Cend}V`$, and an action $`\rho `$ of $`\mathrm{Cend}V`$ on $`V`$.
###### Proposition 10.5.
(i) For any finite $`H`$-module $`V`$, the above pseudoproduct provides $`\mathrm{Cend}V`$ with the structure of an associative $`H`$-pseudoalgebra. $`V`$ has a natural structure of a $`\mathrm{Cend}V`$-module given by $`\varphi v\varphi (v)`$.
(ii) For an associative $`H`$-pseudoalgebra $`A`$, giving a structure of an $`A`$-module on $`V`$ is equivalent to giving a homomorphism of associative $`H`$-pseudoalgebras from $`A`$ to $`\mathrm{Cend}V`$.
###### Proof.
Part (i) is an immediate consequence of Lemma 10.4. Indeed, the only thing that remains to be checked is the associativity of $`\mathrm{Cend}V`$, and it follows from (10.7):
$`\left((\varphi \psi )\chi \right)v`$ $`=(\varphi \psi )(\chi v)=\varphi (\psi (\chi v))`$
$`=\varphi ((\psi \chi )v)=\left(\varphi (\psi \chi )\right)v.`$
To prove part (ii), we associate with each $`aA`$ the $`H`$-pseudolinear map $`m_a\mathrm{Cend}V`$ given by $`m_a(v)=av`$. Then
$$(m_am_b)v=m_a(m_bv)=a(bv)=(ab)v=m_{ab}v,$$
which shows that $`m_am_b=m_{ab}`$. ∎
We denote by $`\mathrm{gc}V`$ the Lie $`H`$-pseudoalgebra obtained from the associative one $`\mathrm{Cend}V`$ by the construction of Proposition 3.15. Then $`V`$ is a $`\mathrm{gc}V`$-module, and one has a statement analogous to part (ii) above.
###### Proposition 10.6.
Let $`V`$ be a finite $`H`$-module. Then, for a Lie $`H`$-pseudoalgebra $`L`$, giving a structure of an $`L`$-module on $`V`$ is equivalent to giving a homomorphism of Lie $`H`$-pseudoalgebras from $`L`$ to $`\mathrm{gc}V`$.
###### Remark 10.7.
Let $`L`$ be a Lie $`H`$-pseudoalgebra, and $`U,V`$ be finite $`L`$-modules. Then the formula ($`aL`$, $`uU`$, $`\varphi \mathrm{Chom}(U,V)`$)
(10.8)
$$(a\varphi )(u)=a(\varphi u)((\sigma \mathrm{id})_H\mathrm{id})\varphi (au)$$
provides $`\mathrm{Chom}(U,V)`$ with the structure of an $`L`$-module.
###### Definition 10.8.
(i) Let $`A`$ be an associative $`H`$-pseudoalgebra. A derivation of $`A`$ is an $`H`$-pseudolinear map $`\varphi \mathrm{gc}A`$ which satisfies
(10.9)
$$\varphi (ab)=(\varphi a)b+((\sigma \mathrm{id})_H\mathrm{id})a(\varphi b),a,bA.$$
We denote the space of all such $`\varphi `$ by $`\mathrm{Der}A`$.
(ii) Similarly, for a Lie $`H`$-pseudoalgebra $`L`$, let $`\mathrm{Der}L`$ be the space of all $`\varphi \mathrm{gc}L`$ satisfying
(10.10)
$$\varphi [ab]=[(\varphi a)b]+((\sigma \mathrm{id})_H\mathrm{id})[a(\varphi b)],a,bL.$$
The next result follows easily from definitions.
###### Lemma 10.9.
(i) For any $`H`$-pseudoalgebra $`A`$, $`\mathrm{Der}A`$ is a subalgebra of $`\mathrm{gc}A`$.
(ii) When $`A`$ is associative (respectively Lie), we have a homomorphism of pseudoalgebras $`i:A\mathrm{Der}A`$ given by $`i(a)(b)=ab(\sigma _H\mathrm{id})ba`$ $`(`$respectively $`i(a)(b)=[ab])`$, whose kernel is the center of $`A`$.
(iii) For any $`xX`$ and $`\varphi \mathrm{Der}A`$, $`\varphi _x`$ is a derivation of the annihilation algebra of $`A`$. In other words, we have: $`𝒜(\mathrm{Der}A)\mathrm{Der}𝒜(A)`$.
(iv) Let $`A`$ be an associative $`H`$-pseudoalgebra and $`L`$ be the corresponding Lie pseudoalgebra with pseudobracket given by commutator. Then $`\mathrm{Der}A\mathrm{Der}L`$.
###### Example 10.10.
Recall that in Section 8.3 we defined the $`W(𝔡)`$-module of pseudoforms $`\mathrm{\Omega }(𝔡)=H^{}𝔡^{}`$. Since $`^{}𝔡^{}`$ is an associative algebra with respect to the wedge product, we can consider $`\mathrm{\Omega }(𝔡)`$ as an associative pseudoalgebra: the current pseudoalgebra over $`^{}𝔡^{}`$. Then, as in the case of usual differential forms, for any $`\alpha W(𝔡)`$, $`\alpha `$ and $`\alpha _\iota `$ are superderivations of $`\mathrm{\Omega }(𝔡)`$, see (8.13, 8.14).
In the case when $`V`$ is a free $`H`$-module of finite rank, one can give an explicit description of $`\mathrm{Cend}V`$, and hence of $`\mathrm{gc}V`$, as follows.
###### Proposition 10.11.
Let $`V=HV_0`$, where $`H`$ acts trivially on $`V_0`$ and $`dimV_0<\mathrm{}`$. Then $`\mathrm{Cend}V`$ is isomorphic to $`HH\mathrm{End}V_0`$, with $`H`$ acting by a left multiplication on the first factor, and with the following pseudoproduct:
(10.11)
$$(faA)(gbB)=(fga_{(1)})_H(1ba_{(2)}AB).$$
The action of $`\mathrm{Cend}V`$ on $`V=HV_0`$ is given by:
(10.12)
$$(faA)(hv)=(fha)_H(1Av).$$
The pseudobracket in $`\mathrm{gc}V`$ is given by:
(10.13)
$$\begin{array}{cc}\hfill [(faA)(gbB)]& =(fga_{(1)})_H(1ba_{(2)}AB)\hfill \\ & (fb_{(1)}g)_H(1ab_{(2)}BA).\hfill \end{array}$$
###### Proof.
Since $`(HH)_HVHHV_0`$, we can identify $`\mathrm{Cend}V`$ with $`HH\mathrm{End}V_0`$ so that its action on $`V`$ is given by (10.12). To prove (10.11), we use (10.7) and the explicit definition of associativity from Section 3. Due to $`H`$-bilinearity, we can assume that $`f=g=h=1`$. Then:
$`(1aA)\left((1bB)(1v)\right)`$ $`=(1aA)\left((1b)_H(1Bv)\right)`$
$`=(1a_{(1)}ba_{(2)})_H(1ABv).`$
On the other hand, we have:
$$\left((1a_{(1)})_H(1ba_{(2)}AB)\right)(1v)=(1a_{(1)}ba_{(2)})_H(1ABv).$$
Now (10.11) follows from the uniqueness from Lemma 10.4. ∎
###### Remark 10.12.
Let $`V=HV_0`$, where $`H`$ acts trivially on $`V_0`$ and $`dimV_0<\mathrm{}`$. Then $`\mathrm{Cur}\mathrm{End}V_0`$ can be identified with $`H1\mathrm{End}V_0\mathrm{Cend}V`$. Similarly, $`\mathrm{Cur}𝔤𝔩V_0`$ is a subalgebra of $`\mathrm{gc}V`$.
When $`V=H𝐤^n`$, we will denote $`\mathrm{Cend}V`$ by $`\mathrm{Cend}_n`$, and $`\mathrm{gc}V`$ by $`\mathrm{gc}_n`$. Of course, the essential case is when $`V=H`$ is of rank one. Let us describe the associative algebra $`𝒜_Y\mathrm{Cend}_1`$, where $`𝒜_Y`$ is as in Section 7. As an $`H`$-module it is isomorphic to $`Y_H\mathrm{Cend}_1YH`$ with $`H`$ acting on the first factor. We have $`a_x=x_H(1a)xa`$ for $`xY`$, $`aH`$. Comparing (7.2) with (10.11), we see that the product in $`YH`$ is given by:
(10.14)
$$(xa)(yb)=x(ya_{(1)})ba_{(2)}.$$
Hence $`𝒜_Y\mathrm{Cend}_1`$ is isomorphic to the smash product $`Y\mathrm{}H`$ (see Section 2). The annihilation algebra $`𝒜(\mathrm{Cend}_1)𝒜_X\mathrm{Cend}_1X\mathrm{}H`$ is isomorphic as an associative algebra to the Drinfeld double of $`H`$ (see \[D\]). For $`H=U(𝔡)`$, $`𝒜(\mathrm{Cend}_1)`$ can be identified with the associative algebra of all differential operators on $`X`$, while $`𝒜(\mathrm{gc}_1)`$ with the corresponding Lie algebra.
###### Example 10.13.
Let $`H=U(𝔡)`$ be the universal enveloping algebra of a Lie algebra $`𝔡`$. We identify $`𝔡`$ with its image in $`H`$, so that $`\mathrm{gc}_1=HH`$ contains $`H𝔡`$. We claim that $`fafa`$ ($`fH`$, $`a𝔡`$) is an embedding of Lie pseudoalgebras $`W(𝔡)\mathrm{gc}_1`$, compatible with their actions on $`H`$. This is immediate by comparing (10.13) with (8.3) and (10.12) with (8.4).
Consider $`H`$ as $`\mathrm{Cur}𝐤`$, i.e., as an associative $`H`$-pseudoalgebra with a pseudoproduct $`fg=(fg)_H1`$. Then $`W(𝔡)=\mathrm{Der}H\mathrm{gc}_1`$.
###### Example 10.14.
Let $`H=𝐤[\mathrm{\Gamma }]`$ be the group algebra of a group $`\mathrm{\Gamma }`$. Then for $`V=H`$ and $`f,g,a,b\mathrm{\Gamma }`$, (10.11) takes the form:
$$(fa)(gb)=(fga)_H(1ba).$$
We end this section with two lemmas that will be useful in representation theory.
###### Lemma 10.15.
For $`\varphi \mathrm{Chom}(V,W)`$, let
$$\mathrm{ker}_n\varphi =\{vV|\varphi _xv=0x\mathrm{F}_nX\},$$
so that, for example, $`\mathrm{ker}_1\varphi =\mathrm{ker}\varphi `$. If $`V`$ is a finite $`H`$-module and $`\mathrm{F}^nH`$ is finite dimensional, then $`\mathrm{ker}_n\varphi /\mathrm{ker}\varphi `$ is finite dimensional.
###### Proof.
Since $`\mathrm{ker}\varphi `$ is an $`H`$-submodule of $`V`$, after replacing $`V`$ with $`V/\mathrm{ker}\varphi `$, we can assume that $`\mathrm{ker}\varphi =0`$.
By definition, $`\mathrm{ker}_n\varphi =\varphi ^1\left((\mathrm{F}^nH𝐤)_HW\right)`$. Since, by Lemma 2.37, $`(\mathrm{F}^nH𝐤)_HW=(𝐤\mathrm{F}^nH)_HW`$, we have $`\varphi (\mathrm{ker}_n\varphi )(𝐤\mathrm{F}^nH)_HW`$.
On the other hand, since $`V`$ is finite over $`H`$ and $`\varphi `$ satisfies (10.1), there exists a finite-dimensional subspace $`W^{}`$ of $`W`$ such that $`\varphi (\mathrm{ker}_n\varphi )(𝐤H)_HW^{}`$. It follows that $`\varphi (\mathrm{ker}_n\varphi )(𝐤\mathrm{F}^nH)_HW^{}`$, which is finite dimensional. Since $`\varphi `$ is injective, $`\mathrm{ker}_n\varphi `$ is finite dimensional. ∎
###### Lemma 10.16.
Let $`\varphi \mathrm{Chom}(V,W)`$ and $`hH`$. If $`h`$ is not a divisor of zero, then:
(i) $`h\varphi =0`$ implies $`\varphi =0;`$
(ii) $`hv\mathrm{ker}\varphi `$ implies $`v\mathrm{ker}\varphi `$.
###### Proof.
Part (i) follows from Eq. (10.2): if $`\varphi (v)=_i(f_i1)_Hw_i`$ with linearly independent $`w_i`$, then $`(h\varphi )(v)=_i(hf_i1)_Hw_i`$ can be zero only if all $`hf_i=0`$, which implies $`f_i=0`$.
Similarly, part (ii) follows from (10.1), since we can write $`\varphi (v)`$ uniquely in the form $`_i(1g_i)_Hw_i`$. ∎
###### Corollary 10.17.
Let $`L`$ be a pseudoalgebra, and $`M`$ be an $`L`$-module. Then any torsion element from $`L`$ acts trivially on $`M`$, and any torsion element from $`M`$ is acted on trivially by $`L`$. In particular, the torsion of a Lie pseudoalgebra is central.
## 11. Reconstruction of an $`H`$-Pseudoalgebra from an $`H`$-Differential Algebra
### 11.1. The reconstruction functor $`𝒞`$
Let, as before, $`H`$ be a cocommutative Hopf algebra and $`X=H^{}`$. Given a topological left $`H`$-module $``$ (where $`H`$ is endowed with the discrete topology), let
(11.1)
$$𝒞()=\mathrm{Hom}_H^{\mathrm{cont}}(X,)$$
be the space of continuous $`H`$-homomorphisms. We define a structure of a left $`H`$-module on $`𝒞()`$ by
(11.2)
$$(h\alpha )(x)=\alpha (xh).$$
Then $`𝒞`$ is a covariant functor from the category of topological $`H`$-modules to the category of $`H`$-modules.
###### Lemma 11.1.
(i) The functor $`𝒞`$ is left exact: $`𝒞(i)`$ is injective if $`i`$ is injective.
(ii) The functor $`𝒞`$ preserves direct sums: $`𝒞(_1_2)=𝒞(_1)𝒞(_2)`$.
(iii) Assume that the Hopf algebra $`H`$ contains nonzero primitive elements. If $``$ is finite dimensional over $`𝐤`$ with discrete topology, then $`𝒞()=0`$.
(iv) If $`H=U(𝔡)`$, then $`𝒞()`$ has no torsion as an $`H`$-module.
###### Proof.
Parts (i) and (ii) are obvious.
(iii) By Kostant’s Theorem 2.1, $`H=U(𝔡)\mathrm{}𝐤[\mathrm{\Gamma }]`$ with $`𝔡0`$. If $``$ is finite dimensional, any continuos homomorphism $`\alpha :X`$ must contain some $`\mathrm{F}_nX`$ in its kernel. Let $`h\mathrm{F}^{n1}U(𝔡)`$ but $`h\mathrm{F}^{n2}U(𝔡)`$. Then, by Lemma 6.10, $`h\mathrm{F}_nX=X`$ so that for each $`xX`$, $`x=hy`$ for some $`y\mathrm{F}_nX`$. This implies $`\alpha (x)=\alpha (hy)=h(\alpha (y))=0`$, since $`\alpha (y)=0`$, proving part (iii).
Similarly, part (iv) follows from the fact that $`Xh=X`$ for any nonzero $`hU(𝔡)`$. ∎
If, in addition, $``$ is a topological Lie $`H`$-differential algebra, we define $`x`$-brackets in $`𝒞()`$ by the formula (cf. (9.24)):
(11.3) $`[\alpha _x\beta ](y)`$ $`=[\alpha (x_{(2)}),\beta (yx_{(1)})]=_i[\alpha (x_i),\beta (y(h_iS(x)))].`$
This is well defined because the infinite sum in the right-hand side converges in $``$. Equation (11.3) is also equivalent to (cf. (9.23)):
(11.4) $`[\alpha (x),\beta (y)]`$ $`=[\alpha _{x_{(2)}}\beta ](yx_{(1)})=_i[\alpha _{x_i}\beta ](y(xS(h_i))).`$
###### Proposition 11.2.
For any topological Lie $`H`$-differential algebra $``$, $`𝒞()`$ satisfies properties (9.15)–(9.18).
###### Proof.
This can be verified by straightforward but rather tedious computations. To illustrate them, let us check (9.15). By definition, we have:
$$[h\alpha _x\beta ](y)=_i[(h\alpha )(x_i),\beta (y(h_iS(x)))]=_i[\alpha (x_ih),\beta (y(h_iS(x)))],$$
while
$$[\alpha _{xh}\beta ](y)=_i[\alpha (x_i),\beta (y(h_iS(xh)))].$$
Hence (9.15) is a consequence of the following identity:
(11.5)
$$_ix_ihh_i=_ix_ih_iS(h),$$
which can be checked by pairing both sides with $`fzHX`$. Indeed,
$$_ix_ih,fh_i,z=zh,f=z,fS(h)=_ix_i,fh_iS(h),z.$$
A more conceptual proof can be given by noticing that formula (11.3) is the same as the formula for the commutator of $`H`$-pseudolinear maps. For $`\alpha 𝒞()`$ consider the family $`\mathrm{ad}\alpha (x)\mathrm{Hom}(,)`$ indexed by $`xX`$. It is easy to see that it satisfies (10.5). So, if it also satisfies (10.4), it would give an $`H`$-pseudolinear map from $``$ to itself. Although this is not true in general, the argument still works because all infinite series that appear will be convergent. (In other words, we embed $`𝒞()`$ in a certain completion of $`\mathrm{gc}`$.) ∎
In order to have the locality (9.13), one has to impose some restrictions on $``$. In particular, the condition that $``$ be a linearly compact topological Lie algebra will often suffice to guarantee locality of $`𝒞()`$.
In what follows, we explain how to reconstruct an $`H`$-pseudoalgebra $`L`$ which is finitely generated over $`H`$ from its annihilation Lie algebra $`𝒜(L)`$. Recall that $`𝒜(L)`$ is a linearly compact topological Lie algebra (Proposition 7.12). In many of the proofs we never exploit the algebra structure on $`𝒜(L)`$, so the corresponding statements hold for finite $`H`$-modules in general.
We let $`\widehat{L}=𝒞𝒜(L):=\mathrm{Hom}_H^{\mathrm{cont}}(X,X_HL)`$. There is an obvious map
(11.6)
$$\mathrm{\Phi }:L\widehat{L},a\alpha (x)=x_Ha.$$
It is clear by definitions that $`\mathrm{\Phi }`$ is a homomorphism of $`H`$-pseudoalgebras (or $`H`$-modules if $`L`$ is only an $`H`$-module).
### 11.2. The case of free modules
Let $`L`$ be a Lie $`H`$-pseudoalgebra which is free as an $`H`$-module: $`L=HL_0`$ with the trivial action of $`H`$ on $`L_0`$. Then $`=𝒜(L):=X_H(HL_0)XL_0`$ as an $`H`$-module.
###### Proposition 11.3.
When $`L`$ is a Lie $`H`$-pseudoalgebra that is a free $`H`$-module, the map $`\mathrm{\Phi }`$ defined by (11.6) is an isomorphism of Lie $`H`$-pseudoalgebras.
###### Proof.
To construct the inverse of $`\mathrm{\Phi }`$, identify $``$ with $`XL_0`$ and consider
$$\mathrm{\Psi }:\widehat{L}L,\alpha _iS(h_i)(\epsilon \mathrm{id})\alpha (x_i).$$
Here, as before, $`\{h_i\}`$, $`\{x_i\}`$ are dual bases in $`H`$ and $`X`$, and $`\epsilon (x)=1,x`$ for $`xX`$. This is well defined, i.e., the sum is finite, because $`\alpha (x_i)\mathrm{F}_1XL_0`$ for all but a finite number of $`x_i`$ and $`\epsilon (\mathrm{F}_1X)=0`$. Using identity (11.5), it is easy to see that $`\mathrm{\Psi }`$ is $`H`$-linear. Next, we have for $`aL_0`$:
$$\mathrm{\Psi }\mathrm{\Phi }(1a)=_iS(h_i)1,x_ia=S(1)a=1a,$$
showing that $`\mathrm{\Psi }\mathrm{\Phi }=\mathrm{id}`$. In particular, $`\mathrm{\Psi }`$ is surjective.
Assume that $`\mathrm{\Psi }(\alpha )=0`$ for some $`\alpha \widehat{L}`$. This means that $`(1,\mathrm{id})\alpha (x)=0`$ for any $`xX`$. But then for any $`hH`$, we have:
$`\left(S(h),\mathrm{id}\right)\alpha (x)`$ $`=\left(1,\mathrm{id}\right)\left((h1)\alpha (x)\right)`$
$`=\left(1,\mathrm{id}\right)\left(h(\alpha (x))\right)=\left(1,\mathrm{id}\right)\alpha (hx)=0,`$
which implies $`\alpha =0`$. Hence $`\mathrm{\Psi }`$ is injective. ∎
###### Remark 11.4.
If $`L`$ is only a free $`H`$-module, then $`\mathrm{\Phi }`$ is an isomorphism of $`H`$-modules. Analogous results hold for representations, or for associative $`H`$-pseudoalgebras.
### 11.3. Reconstruction of a non-free module
Throughout this subsection $`L`$ will be a (possibly non-free) finitely generated $`H`$-module, and $`H`$ will be the universal enveloping algebra $`U(𝔡)`$ of a finite-dimensional Lie algebra $`𝔡`$.
The natural map $`\mathrm{\Phi }:L\widehat{L}`$ (see (11.6)) is in general neither injective nor surjective. However, the induced map $`\phi =𝒜(\mathrm{\Phi }):𝒜(L)𝒜(\widehat{L})`$ has a left inverse $`\psi :x_H\alpha \alpha (x)`$. This shows that $`𝒜(\mathrm{\Phi })`$ is injective, and that $`\psi `$ is surjective.
We want to figure out to what extent injectivity and surjectivity of $`\mathrm{\Phi }`$ fail. First of all let us remark that, by Lemma 7.4, every torsion element $`aL`$ has all zero Fourier coefficients, i.e., it belongs to the kernel of $`\mathrm{\Phi }`$. In fact, we have:
###### Proposition 11.5.
For any finite $`H`$-module $`L`$, the kernel of $`\mathrm{\Phi }:L\widehat{L}`$ equals the torsion of $`L`$.
###### Proof.
It remains to show that a non-torsion element $`aL`$ does not lie in the kernel of $`\mathrm{\Phi }`$. Consider the map $`i:LF`$ constructed in Lemma 2.2. Then $`i(a)0`$. The map $`𝒜(i)`$ induced by $`i`$ maps the Fourier coefficient $`x_Ha`$ of $`a`$ to the corresponding Fourier coefficient $`x_Hi(a)`$ of a nonzero element in the free $`H`$-module $`F`$. Now, it is clear from Proposition 11.3 that $`x_Hi(a)0`$ for some $`xX`$, hence $`x_Ha`$ must be nonzero too. ∎
###### Corollary 11.6.
A finite $`H`$-module $`L`$ is torsion iff $`X_HL=0`$.
###### Corollary 11.7.
Let $`M,N`$ be finite $`H`$-modules, $`f:MN`$ be an $`H`$-linear map, and assume $`N`$ to be torsionless. Then $`𝒜(f)=0`$ if and only if $`f=0`$.
###### Proof.
$`𝒜(f)=0`$ means that $`X_Hf(M)=0`$, hence $`f(M)N`$ is torsion. ∎
###### Remark 11.8.
By Corollary 10.17, the torsion of a Lie $`H`$-pseudoalgebra $`L`$ is always central, hence the map $`\mathrm{\Phi }`$ is injective if $`L`$ is centerless.
###### Remark 11.9.
If $`\mathrm{\Phi }`$ is an isomorphism, then $`\mathrm{\Phi }^1`$ induces $`\psi `$, i.e., $`\psi =𝒜(\mathrm{\Phi }^1)`$. Corollary 11.7 tells us that if $`L`$ is torsionless and $`\psi `$ is induced by some map $`\mathrm{\Psi }`$, then $`\mathrm{\Phi }`$ is an isomorphism and $`\mathrm{\Psi }=\mathrm{\Phi }^1`$.
###### Proposition 11.10.
For any map of finite $`H`$-modules $`f:MN`$, the following conditions are equivalent:
1. $`N/f(M)`$ is torsion.
2. $`𝒜(f):𝒜(M)𝒜(N)`$ is surjective.
3. $`\mathrm{gw}\mathrm{coker}𝒜(f)<dim𝔡`$.
###### Proof.
To show the equivalence of (1) and (2), it is enough to tensor the exact sequence $`M\stackrel{𝑓}{}NN/f(M)0`$ with $`X`$, getting the exact sequence $`𝒜(M)𝒜(N)X_HN/f(M)0`$, and to apply Corollary 11.6.
Assume that (3) holds but $`N/f(M)`$ is not torsion. Then it contains a nonzero element $`a`$ which generates a free $`H`$-module. Then $`X_HaX`$ has growth $`dim𝔡`$, which is a contradiction. ∎
### 11.4. Reconstruction of a Lie pseudoalgebra
Now let $`L`$ be a Lie $`H`$-pseudoalgebra which is finite as an $`H`$-module. Again, $`H=U(𝔡)`$ will be the universal enveloping algebra of a finite-dimensional Lie algebra $`𝔡`$. Let $`=𝒜(L)`$ be the annihilation Lie algebra of $`L`$, and $`\widehat{L}=𝒞()`$, as before.
By Proposition 11.2, $`\widehat{L}`$ satisfies all the properties of a Lie $`H`$-pseudoalgebra except the locality (9.13). An indirect way to establish the locality property for $`\widehat{L}`$ is by embedding it in the bigger (local) Lie pseudoalgebra $`\mathrm{gc}L`$. In order to map $`\widehat{L}`$ to $`\mathrm{gc}L`$, we need to assign to each element of $`\widehat{L}`$ a pseudolinear map from $`L`$ to itself.
This can be done as follows. Recall that $``$ acts on $`L`$ by $`(x_Ha)b=[a_xb]`$ for $`a,bL`$, $`xX`$ (see Proposition 9.4). Now in terms of $`x`$-products the action of $`\widehat{L}`$ on $`L`$ is given by $`\alpha _xb=\alpha (x)b`$. The locality condition $`\alpha _xb=0`$ for $`x\mathrm{F}_nX`$, $`n0`$ is satisfied because $`\alpha `$ is continuous and $`L`$ is a discrete topological $``$-module (see Proposition 9.4). All the other axioms of a Lie pseudoalgebra representation follow easily from definitions.
We now need to find conditions for the above-defined $`j:\widehat{L}\mathrm{gc}L`$ to be injective. Then $`\widehat{L}`$ would embed into $`\mathrm{gc}L`$, which will show locality.
###### Lemma 11.11.
If $`L`$ is torsionless, the kernel of the above-defined $`j:\widehat{L}\mathrm{gc}L`$ consists of all elements $`\alpha `$ such that $`\alpha (X)`$ is contained in the center of $``$.
###### Proof.
Since $`L`$ is torsionless, $`\mathrm{\Phi }`$ is injective by Proposition 11.5. Hence, for $`a,bL`$, $`xX`$, one has $`[a_xb]=0`$ iff $`[a_xb]_y=0`$ for all $`yX`$. By (9.23), (9.24), this is equivalent to $`[a_x,b_y]=0`$. Hence, for $`l`$, $`lb=0`$ for all $`b`$ iff $`l`$ lies in the center of $``$. Now $`\alpha \widehat{L}`$ is in the kernel of $`j`$ iff $`\alpha (x)b=0`$ for all $`x`$ and $`b`$, which means that $`\alpha (x)`$ is central for all $`x`$. ∎
###### Lemma 11.12.
If $`L`$ is torsionless and $`=𝒜(L)`$ has a finite-dimensional center, then $`j:\widehat{L}\mathrm{gc}L`$ is injective.
###### Proof.
Let $`\alpha \widehat{L}`$ be in the kernel of $`j`$; then by the previous lemma $`\alpha (X)`$ is contained in the center of $``$. The latter is finite dimensional by assumption, so the kernel $`N`$ of $`\alpha `$ is of finite codimension in $`X`$. This implies that $`N`$ is open in $`X`$, and it contains $`\mathrm{F}_iX`$ for some $`i`$. Let $`h\mathrm{F}^{i+1}H`$ but $`h\mathrm{F}^iH`$; then by Lemma 6.10, $`h\mathrm{F}_iX=\mathrm{F}_1X=X`$. Since $`\alpha `$ is $`H`$-linear, $`N`$ is an $`H`$-submodule of $`X`$. Then $`X=h\mathrm{F}_iXhNN`$, therefore $`N=X`$ and $`\alpha =0`$. ∎
###### Proposition 11.13.
Let $`L`$ be a Lie $`H`$-pseudoalgebra which is finite and torsionless as an $`H`$-module. If its annihilation Lie algebra $`=𝒜(L)`$ has a finite-dimensional center, then $`\widehat{L}=𝒞()`$ is a Lie $`H`$-pseudoalgebra containing $`L`$ as an ideal.
###### Proof.
The only thing that remains to be checked is the locality property for $`\widehat{L}`$. It follows from that of $`\mathrm{gc}L`$, since in this case $`j:\widehat{L}\mathrm{gc}L`$ is injective. ∎
In the proof of Lemma 11.11 we have shown that, if $`L`$ is finite and torsionless, the kernel of the action of $`𝒜(L)`$ on $`L`$ is exactly the center of $`𝒜(L)`$. This implies the following result which was used in the proof of Theorem 8.24.
###### Lemma 11.14.
Let $`i:LL_1`$ be an injective map of Lie $`H`$-pseudoalgebras, and assume that $`L`$ is finite and torsionless. Then the kernel of the induced map $`𝒜(i):𝒜(L)𝒜(L_1)`$ is contained in the center of $`𝒜(L)`$.
###### Proof.
The kernel of $`𝒜(i)`$ acts trivially on $`L_1`$ and hence on $`L`$. ∎
In Section 13.4 we will need the following lemma.
###### Lemma 11.15.
Let $`L`$ be a subalgebra of $`\mathrm{gc}V`$ for some finite $`H`$-module $`V`$ $`(L`$ may be infinite$`)`$. Then the map $`\mathrm{\Phi }:L\widehat{L}`$ is injective.
###### Proof.
Assume that $`a`$ belongs to $`\mathrm{ker}\mathrm{\Phi }`$; then all $`x_Ha=0`$, $`xX`$. This implies that all Fourier coefficients $`a_x\mathrm{End}V`$ of $`a\mathrm{gc}V`$ are zero, hence $`a=0`$. ∎
For any topological Lie $`H`$-differential algebra $``$, we have a natural homomorphism $`\psi :𝒜𝒞()`$, given by $`x_Haa(x)`$ for $`a𝒞()`$, $`xX`$. The map $`\psi `$ does not need to be surjective, but we have a good control on injectivity, which can sometime prove useful.
###### Lemma 11.16.
The kernel of $`\psi :𝒜𝒞()`$ lies in the center of $`𝒜𝒞()`$.
###### Proof.
Follows easily from (9.29) and (9.30). Say that $`\alpha `$ lies in the kernel of $`\psi `$. Since $`\psi `$ is a homomorphism of Lie $`H`$-differential algebras, its kernel is an $`H`$-stable ideal of $`𝒜𝒞()`$. Then by (9.29), $`(\alpha b)_y\mathrm{ker}\psi `$ for all $`b𝒞()`$, $`yX`$, because in the right-hand side all elements $`h_i\alpha `$ lie in $`\mathrm{ker}\psi `$. This means that $`(\alpha b)(y)=0`$ for all $`yX`$, hence $`\alpha b=0`$ for every $`b𝒞()`$. Now, use this in (9.30) to obtain that $`\alpha `$ is central. ∎
## 12. Reconstruction of Pseudoalgebras of Vector Fields
In this section, we show that the reconstruction procedure of Section 11, when applied to the primitive Lie algebras of vector fields (or current algebras over them), gives the primitive pseudoalgebras of vector fields defined in Section 8 (or current pseudoalgebras over them).
As before, $`𝔡`$ will be an $`N`$-dimensional Lie algebra, and $`H=U(𝔡)`$ its universal enveloping algebra. $``$ will be a Lie algebra provided with an action of $`𝔡`$ and a filtration by subspaces $`=_1_0\mathrm{}`$. When $``$ is a subalgebra of $`W_N`$, it will always be considered with the filtration induced by the canonical filtration of $`W_N`$.
The Lie algebra $`\mathrm{Der}`$ of derivations of $``$ has the induced filtration:
$$(\mathrm{Der})_i:=\{d\mathrm{Der}|d(_j)_{i+j}j\}.$$
The action of $`𝔡`$ is called transitive if the composition of the homomorphism $`𝔡\mathrm{Der}`$ and the projection $`\mathrm{Der}\mathrm{Gr}_1(\mathrm{Der}):=\mathrm{Der}/(\mathrm{Der})_0`$ is a linear isomorphism. This is equivalent to the following two conditions: $`𝔡\mathrm{Der}`$ intersects $`(\mathrm{Der})_0`$ trivially and $`dim\mathrm{Gr}_1(\mathrm{Der})=N`$.
###### Lemma 12.1.
Let $`L`$ be a current Lie $`H`$-pseudoalgebra over a finite-dimensional simple Lie algebra or over one of the primitive pseudoalgebras of vector fields. Then the action of $`𝔡`$ on its annihilation Lie algebra $`=𝒜(L)`$ is transitive.
###### Proof.
By Theorem 8.24, $`=𝒪_r\widehat{}^{}`$ is a current Lie algebra over $`^{}`$, where $`^{}`$ is either a finite-dimensional simple Lie algebra $`𝔤`$ (for $`r=N=dim𝔡`$), or one of the Lie algebras of vector fields $`W_N^{}`$, $`S_N^{}`$, $`P_N^{}`$ or $`K_N^{}`$ ($`N^{}=Nr`$). In particular, we know that $`dim\mathrm{Gr}_1(\mathrm{Der})=N`$. By Lemma 7.13, a sufficiently high power of any nonzero element $`a𝔡`$ maps any given open subspace of $``$ surjectively onto $``$. This cannot hold if $`a`$ belongs to $`(\mathrm{Der})_0`$, therefore $`𝔡\mathrm{Der}`$ is injective and the image of $`𝔡`$ intersects $`(\mathrm{Der})_0`$ trivially. Comparing the dimensions, we get that $`𝔡\mathrm{Gr}_1(\mathrm{Der})`$ is an isomorphism. ∎
The main results of this section can be summarized by Theorem 12.2 below. Its proof follows from Sections 12.112.7.
###### Theorem 12.2.
Let $`=𝒪_r\widehat{}^{}`$ be a current Lie algebra over $`^{}`$, where $`^{}`$ is a simple linearly compact Lie algebra of growth $`N^{}=Nr`$. Assume that $`𝔡`$ acts transitively on $``$. Then there is a codimension $`r`$ subalgebra $`𝔡^{}`$ of $`𝔡`$, acting transitively on $`^{}`$, such that the $`H`$-pseudoalgebra $`𝒞()`$ is isomorphic to a current pseudoalgebra over the $`H^{}`$-pseudoalgebra $`𝒞(^{})`$, where $`H^{}=U(𝔡^{})`$. Moreover, $`𝒞(^{})`$ is either a finite-dimensional simple Lie algebra (and $`𝔡^{}=0`$) or one of the primitive $`H^{}`$-pseudoalgebras of vector fields $`W(𝔡^{})`$, $`S(𝔡^{},\chi ^{})`$, $`H(𝔡^{},\chi ^{},\omega ^{})`$ or $`K(𝔡^{},\theta ^{})`$.
### 12.1. Reconstruction from $`W_N`$
Recall that $`X=H^{}`$ can be identified with $`𝒪_N=𝐤[[t_1,\mathrm{},t_N]]`$. The action of $`𝔡`$ on $`X`$ gives an action on $`𝒪_N`$ in terms of linear differential operators, i.e., an embedding $`𝔡W_N=\mathrm{Der}𝒪_N`$ which we call the canonical embedding of $`𝔡`$ in $`W_N`$. Note that this embedding is transitive, i.e., $`𝔡W_N`$ is complementary to $`\mathrm{F}_0W_N`$.
A structure of an $`H`$-differential algebra on $`W_N`$ is equivalent to a transitive action of $`𝔡`$ on $`W_N`$ by derivations. Since $`\mathrm{Der}W_N=W_N`$, this is the same as a transitive embedding of $`𝔡`$ in $`W_N`$. By Proposition 6.9, any two such embeddings are equivalent, i.e., conjugate by an automorphism of $`W_N`$. With the canonical action of $`𝔡`$, $`W_N`$ becomes isomorphic to the annihilation algebra of the Lie pseudoalgebra $`W(𝔡)`$ defined in Section 8.1. Since $`W(𝔡)`$ is a free $`H`$-module, Proposition 11.3 shows that the reconstruction of $`W_N`$ is $`W(𝔡)`$, i.e., $`𝒞(W_N)=W(𝔡)`$.
### 12.2. Reconstruction from subalgebras of $`W_N`$
Let $``$ be a linearly compact Lie subalgebra of $`W_N`$, with the induced filtration and with a transitive action of $`𝔡`$ on it. After an automorphism of $`W_N`$, we can assume that the action of $`𝔡`$ is the canonical one. Then $`𝒞()`$ is a subalgebra of $`W(𝔡)=𝒞(W_N)`$, because the functor $`𝒞`$ is left exact. Below we will be concerned with the case when $``$ is the subalgebra consisting of vector fields annihilating some differential form.
Let $`w\mathrm{\Omega }^n(𝔡)`$ be a pseudoform, and $`IH`$ be a right ideal. We denote by $`W(𝔡,w,I)`$ the set of all elements $`\alpha W(𝔡)=H𝔡`$ such that
(12.1)
$$\alpha w(HI)_H\mathrm{\Omega }^n(𝔡).$$
It is easy to check that $`W(𝔡,w,I)`$ is a subalgebra of $`W(𝔡)`$.
###### Lemma 12.3.
Let $`\omega \mathrm{\Omega }_X^n`$ be a differential form, and $`W_N(\omega )`$ be the Lie subalgebra of $`W_N`$ consisting of vector fields annihilating $`\omega `$. If $`\omega =y_Hw`$ for some $`yX`$, $`w\mathrm{\Omega }^n(𝔡)`$, then $`𝒞(W_N(\omega ))`$ is isomorphic to the Lie pseudoalgebra $`W(𝔡,w,I)`$ where $`I=\{hH|yh=0\}`$.
###### Proof.
As was already remarked, $`𝒞(W_N(\omega ))`$ is a subalgebra of $`W(𝔡)`$. Since $`\mathrm{\Omega }^n(𝔡)=H^n𝔡^{}`$ is a free $`H`$-module, we have $`(HH)_H\mathrm{\Omega }^n(𝔡)HH^n𝔡^{}`$. For $`\alpha W(𝔡)`$, write $`\alpha w=_i(f_ig_i)_Hw_i`$ with $`f_i,g_iH`$ and linearly independent $`w_i^n𝔡^{}`$. Then for any $`xX`$ we have (cf. (7.2)):
$$L_{x_H\alpha }\omega =_i(xf_i)(yg_i)w_i.$$
This is zero for any $`x`$ iff $`yg_i=0`$ for all $`i`$, which means $`g_iI`$. ∎
### 12.3. Reconstruction from current algebras over $`W_N^{}`$
Let now $`=𝒪_r\widehat{}W_N^{}`$ be a current algebra over $`W_N^{}`$, and $`𝔡`$ be an $`N=N^{}+r`$ dimensional Lie algebra acting transitively on $``$. Then, by Proposition 6.12, $`𝔡\mathrm{Der}=W_r1+𝒪_r\widehat{}W_N^{}W_N`$. The Lie algebra $``$ is described as the subalgebra of $`W_N`$ consisting of vector fields annihilating the functions $`t_{N^{}+1},\mathrm{},t_N`$, hence it is an intersection of algebras of the form $`W_N(f)`$ ($`f\mathrm{\Omega }_X^0=X`$), see Section 12.2.
After an automorphism of $`W_N`$, we can assume that the action of $`𝔡`$ on it is the canonical one. Then $``$ becomes the intersection of $`W_N(f_i)`$ ($`i=N^{}+1,\mathrm{},N`$) where $`f_iX`$ is the image of $`t_i`$. Now Lemma 12.3 implies that $`𝒞()=W(𝔡,1,I)`$ where $`1\mathrm{\Omega }_0(𝔡)=H`$ and $`I=\{hH|f_ih=0(i=N^{}+1,\mathrm{},N)\}`$.
Recall that for $`\alpha W(𝔡)=H𝔡`$, its action on $`1H`$ is given by $`\alpha 1=\alpha _H1\alpha `$. Therefore $`\alpha W(𝔡,1,I)`$ iff $`\alpha `$ belongs to $`(H𝔡)(HI)=H𝔡^{}`$, where the intersection $`𝔡^{}=𝔡I`$ is a Lie subalgebra of $`𝔡`$. Then $`H^{}=U(𝔡^{})`$ is a Hopf subalgebra of $`H`$, and $`H𝔡^{}H_H^{}(H^{}𝔡^{})`$ is a current pseudoalgebra over $`H^{}𝔡^{}=W(𝔡^{})`$. We have thus proved the following lemma.
###### Lemma 12.4.
The reconstruction of a current Lie algebra over $`W_N^{}`$, provided with a transitive action of a Lie algebra $`𝔡`$, is a current Lie $`H`$-pseudoalgebra over $`W(𝔡^{})`$ where $`𝔡^{}`$ is an $`N^{}`$-dimensional Lie subalgebra of $`𝔡`$.
This result is a special case of Lemma 12.8 below.
### 12.4. Solving compatible systems of linear differential equations
Let $`A`$ be any associative $`𝐤`$-algebra, and let $`𝒪_r=𝐤[[t_1,\mathrm{},t_r]]`$, $`W_r=\mathrm{Der}𝒪_r`$, as before. For fixed $`n0`$, let $`f_i(t)A[[t_1,\mathrm{},t_r]]`$ $`(i=1,\mathrm{},r+n)`$ be formal power series with coefficients in $`A`$, where $`t=(t_1,\mathrm{},t_r)`$. Note that $`W_r`$ acts on $`A[[t_1,\mathrm{},t_r]]`$ by derivations.
Given $`r+n`$ linear differential operators $`D_1,\mathrm{},D_{r+n}W_r`$, consider the following system of differential equations for an unknown $`y(t)A[[t_1,\mathrm{},t_r]]`$:
(12.2)
$$D_i(y(t))=y(t)f_i(t),i=1,\mathrm{},r+n.$$
We assume that the operators $`D_i`$ satisfy
(12.3)
$$[D_i,D_j]=_kc_{ij}^k(t)D_k\text{with}c_{ij}^k(t)𝒪_r;$$
in other words, the space of all operators of the form $`_ip_i(t)D_i`$ with $`p_i(t)𝒪_r`$ is a Lie algebra.
Suppose we have found a solution to the system (12.2). Combining equations (12.2) and (12.3), we get:
$`[D_i,D_j](y)`$ $`=D_iD_j(y)D_jD_i(y)=D_i(yf_j)D_j(yf_i)`$
$`=yf_if_j+yD_i(f_j)yf_jf_iyD_j(f_i),`$
and
$`[D_i,D_j](y)`$ $`=_kc_{ij}^kD_k(y)=_kc_{ij}^kyf_k.`$
The system (12.2) is called compatible if
(12.4)
$$[f_i(t),f_j(t)]+D_i(f_j(t))D_j(f_i(t))=_kc_{ij}^k(t)f_k(t)\text{for all}i,j.$$
When $`y(t)`$ is not a divisor of zero in $`A[[t_1,\mathrm{},t_r]]`$ the compatibility of the system is a necessary condition for having a solution. The compatibility (12.4) is equivalent to saying that $`_ip_i(t)D_i_ip_i(t)(D_i+f_i(t))`$ is a homomorphism of Lie algebras.
We will be interested in solving a more general system of equations than (12.2). Before formulating it, let us note that the above remarks have obvious analogues for systems of the form
(12.5)
$$D_i(z(t))=h_i(t)z(t),i=1,\mathrm{},r+n$$
with $`z(t),h_i(t)A[[t_1,\mathrm{},t_r]]`$. The compatibility of (12.5) is equivalent to (12.4) with $`f_i`$ replaced by $`h_i`$.
Now consider the system
(12.6)
$$D_i(g(t))=g(t)f_i(t)h_i(t)g(t),i=1,\mathrm{},r+n$$
for an unknown $`g(t)A[[t_1,\mathrm{},t_r]]`$. We will show it has a solution, provided that both (12.2) and (12.5) are compatible and some initial conditions at $`t=0`$ are satisfied. (The compatibility of (12.2) and (12.5) implies the compatibility of (12.6).)
###### Proposition 12.5.
In the above notation, let the operators $`D_iW_r`$ satisfy (12.3) and
(12.7)
$$D_i|_{t=0}=\{\begin{array}{cc}_{t_i},\hfill & 1ir,\hfill \\ 0,\hfill & r+1ir+n.\hfill \end{array}$$
Assume that the systems (12.2) and (12.5) are compatible (cf. (12.4)), and that
(12.8)
$$f_i(0)=h_i(0),r+1ir+n.$$
Then the system (12.6) has a unique solution $`g(t)A[[t_1,\mathrm{},t_r]]`$ for any given initial condition $`g(0)A`$ which commutes with $`f_i(0)`$ $`(r+1ir+n)`$.
###### Proof.
For $`r=0`$, both sides of (12.6) are trivial. For $`r1`$, we will proceed by induction on $`r`$.
First of all, note that the compatibility or solvability of the systems (12.2), (12.5) or (12.6) does not change when we apply an automorphism of $`𝒪_r`$. The same is true when we make an elementary transformation: multiply one equation by a function (an element of $`𝒪_r`$) and add it to another equation. For example, we can replace all $`D_i`$ ($`ir`$) by $`D_ip_i(t)D_r`$, and correspondingly $`f_i(t)`$ by $`f_i(t)p_i(t)f_r(t)`$ and $`h_i(t)`$ by $`h_i(t)p_i(t)h_r(t)`$, as long as we do not violate (12.7, 12.8).
Any vector field $`D_rW_r`$ satisfying $`D_r|_{t=0}=_{t_r}`$ can be brought to $`_{t_r}`$ after an automorphism of $`𝒪_r`$, so we will assume that $`D_r=_{t_r}`$. Replacing $`D_i`$ ($`ir`$) by $`D_iD_i(t_r)D_r`$, we can assume in addition that $`D_i(t_r)=0`$ for $`ir`$.
Now it makes sense to put $`t_r=0`$ in the equations with $`ir`$ in (12.6). Let us denote $`\overline{D}_i=D_i|_{t_r=0}`$, $`\overline{f}_i(\overline{t})=f_i(t_1,\mathrm{},t_{r1},0)`$, $`\overline{h}_i(\overline{t})=h_i(t_1,\mathrm{},t_{r1},0)`$, $`\overline{t}=(t_1,\mathrm{},t_{r1})`$. Consider the reduced system
(12.9)
$$\overline{D}_i(\overline{g}(\overline{t}))=\overline{g}(\overline{t})\overline{f}_i(\overline{t})\overline{h}_i(\overline{t})\overline{g}(\overline{t}),i=1,\mathrm{},r1,r+1,\mathrm{},r+n$$
for an unknown $`\overline{g}(\overline{t})A[[t_1,\mathrm{},t_{r1}]]`$. Note that, since $`D_i(t_r)=\delta _{ir}`$, we have: $`[D_i,D_j](t_r)=0`$ for any $`i,j`$, hence $`[D_i,D_j]`$ does not contain $`D_r`$. In particular, putting $`t_r=0`$ we see that the operators $`\overline{D}_i`$ satisfy (12.3). The other assumptions of the proposition are also easy to check, so by induction the system (12.9) has a solution $`\overline{g}(\overline{t})`$.
The equation
(12.10)
$$_{t_r}g(t)=g(t)f_r(t)h_r(t)g(t)$$
has a unique solution $`g(t)`$ satisfying the initial condition
(12.11)
$$g(t_1,\mathrm{},t_{r1},0)=\overline{g}(\overline{t}).$$
We claim that this $`g(t)`$ is then a solution of the system (12.6). Indeed, it satisfies (12.6) for $`t_r=0`$. Next, we compute for $`ir`$ (using (12.9), (12.10), and the compatibility of (12.2), (12.5)):
$`D_r`$ $`D_i(g)|_{t_r=0}=[D_r,D_i](g)|_{t_r=0}+D_iD_r(g)|_{t_r=0}`$
$`=_j\overline{c}_{ri}^j\overline{D}_j\overline{g}+\overline{D}_i(\overline{g}\overline{f}_r\overline{h}_r\overline{g})`$
$`=_j\overline{c}_{ri}^j(\overline{g}\overline{f}_j\overline{h}_j\overline{g})+(\overline{g}\overline{f}_i\overline{h}_i\overline{g})\overline{f}_r+\overline{g}\overline{D}_i(\overline{f}_r)\overline{D}_i(\overline{h}_r)\overline{g}\overline{h}_r(\overline{g}\overline{f}_i\overline{h}_i\overline{g})`$
$`=g\left(D_r(f_i)+f_rf_i\right)|_{t_r=0}\left(D_r(h_i)h_ih_r\right)g|_{t_r=0}\overline{h}_i\overline{g}\overline{f}_r\overline{h}_r\overline{g}\overline{f}_i`$
$`=D_r(gf_ih_ig)|_{t_r=0}.`$
This shows that $`_{t_r}(D_i(g)gf_i+h_ig)|_{t_r=0}=0`$. We can apply the same argument with $`[D_r,D_i]`$ instead of $`D_i`$, and so on, to show that all derivatives with respect to $`t_r`$ vanish at $`t_r=0`$. ∎
###### Remark 12.6.
Any solution $`g(t)`$ of the system (12.6), such that $`g(0)`$ is invertible in $`A`$, is invertible in $`A[[t_1,\mathrm{},t_r]]`$. Its inverse $`g(t)^1`$ satisfies (12.6) with $`f_ih_i`$.
### 12.5. Reconstruction from a current Lie algebra
Let $`^{}`$ be a simple linearly compact Lie algebra, and let $`=𝒪_r\widehat{}^{}`$ be a current algebra over $`^{}`$. The filtration by subspaces $`^{}=_1^{}_0^{}\mathrm{}`$ and the canonical filtration of $`𝒪_r`$ give rise to the product filtration of $``$. Assume that $`𝔡`$ acts on $``$ transitively by derivations. By Proposition 6.12, we have $`\mathrm{Der}=W_r1+𝒪_r\widehat{}\mathrm{Der}^{}`$.
Denote by $`j`$ the embedding $`𝔡\mathrm{Der}`$, and by $`p`$ the projection $`\mathrm{Der}W_r`$. The preimage $`𝔡^{}:=(pj)^1(\mathrm{F}_0W_r)`$ is a Lie subalgebra of $`𝔡`$ of codimension $`r`$. We have $`𝔡^{}\mathrm{F}_0W_r1+𝒪_r\widehat{}\mathrm{Der}^{}`$. The latter contains $`\mathrm{F}_0W_r1+\mathrm{F}_0𝒪_r\widehat{}\mathrm{Der}^{}`$ as an ideal, hence we get a Lie algebra homomorphism $`j^{}:𝔡^{}\mathrm{Der}^{}`$. It leads to a transitive action of $`𝔡^{}`$ on $`^{}`$, because the action of $`𝔡`$ on $``$ is transitive.
###### Lemma 12.7.
Any two transitive embeddings $`j:𝔡\mathrm{Der}`$, that induce the same subalgebra $`𝔡^{}`$ and the same $`j^{}:𝔡^{}\mathrm{Der}^{}`$, are equivalent up to an automorphism of $`\mathrm{Der}`$.
###### Proof.
Let us choose a basis $`\{_i\}`$ of $`𝔡`$ and write $`j(_i)=D_i+f_i(t)`$ ($`i=1,\mathrm{},N=r+n=dim𝔡`$) where $`D_iW_r`$ and $`f_i(t)𝒪_r\widehat{}\mathrm{Der}^{}`$, $`t=(t_1,\mathrm{},t_r)`$. Note that $`D_i=(pj)(_i)`$ and $`pj:𝔡W_r`$ is a Lie algebra homomorphism. We can choose the basis $`\{_i\}`$ in such a way that $`D_i|_{t=0}=_{t_i}`$ for $`1ir`$, and $`D_i|_{t=0}=0`$ for $`r+1ir+n`$. Then $`\{_i\}_{i=r+1,\mathrm{},r+n}`$ is a basis of $`𝔡^{}`$. Moreover, note that $`j^{}:𝔡^{}\mathrm{Der}^{}`$ is given by $`j^{}(_i)=f_i(0)`$.
Let $`\stackrel{~}{j}`$ be another transitive embedding of $`𝔡`$ into $`\mathrm{Der}`$. Since, by Proposition 6.9, the homomorphism $`pj`$ is uniquely determined by the choice of $`𝔡^{}`$, we can assume that $`(p\stackrel{~}{j})(_i)=D_i`$. Then $`\stackrel{~}{j}(_i)=D_i+h_i(t)`$ for some $`h_i(t)𝒪_r\widehat{}\mathrm{Der}^{}`$. By assumption, $`j^{}=\stackrel{~}{j}^{}`$, hence $`f_i(0)=h_i(0)`$ for $`r+1ir+n`$.
Now we want to find an automorphism $`g(t)𝒪_r\widehat{}\mathrm{Aut}^{}`$ such that $`g(0)=\mathrm{id}`$ and $`g(t)(D_i+f_i(t))=(D_i+h_i(t))g(t)`$. This equation is equivalent to (12.6), and it is easy to see that all conditions of Proposition 12.5 are satisfied: for example, the system (12.2) is compatible because $`j`$ is a homomorphism. This completes the proof. ∎
Now given the embedding $`j^{}:𝔡^{}\mathrm{Der}^{}`$ we can consider the reconstruction $`L^{}:=𝒞(^{})`$ which is a Lie $`H^{}`$-pseudoalgebra, where $`H^{}=U(𝔡^{})`$. Given $`L^{}`$ we can take the current $`H`$-pseudoalgebra $`L:=\mathrm{Cur}L^{}=H_H^{}L^{}`$. Since its annihilation Lie algebra $`𝒜(L)`$ is isomorphic to $``$, we get an embedding $`\stackrel{~}{j}:𝔡\mathrm{Der}`$. It induces the same embedding $`j^{}`$ as our initial $`j`$, so by the previous lemma $`j`$ and $`\stackrel{~}{j}`$ are equivalent. But then the reconstruction $`𝒞()`$ of $``$ provided with $`j`$ is isomorphic to the reconstruction of $``$ provided with $`\stackrel{~}{j}`$, which is $`L`$. This can be summarized as follows.
###### Lemma 12.8.
The reconstruction $`𝒞()`$ of a current Lie algebra $`=𝒪_r\widehat{}^{}`$ over a simple linearly compact Lie algebra $`^{}`$, provided with a transitive action of a Lie algebra $`𝔡`$, is a current Lie $`H`$-pseudoalgebra over the $`H^{}`$-pseudoalgebra $`𝒞(^{})`$, where $`H^{}=U(𝔡^{})`$ and $`𝔡^{}`$ is a Lie subalgebra of $`𝔡`$ of codimension $`r`$.
### 12.6. Reconstruction from $`S_N`$
Now consider $`S_N`$ with a transitive action of $`𝔡`$ on it. Since $`\mathrm{Der}S_N=CS_NW_N`$, we have $`𝔡W_N`$. After an automorphism of $`W_N`$, we can assume $`𝔡W_N`$ is the canonical embedding, while $`S_N`$ becomes $`W_N(\omega )(S_N(\omega ))`$ where $`\omega \mathrm{\Omega }_X^N`$ is a volume form. We can write $`\omega =y_Hv`$ with $`yX`$ and $`v^N𝔡^{}`$. Then, by Lemma 12.3, the reconstruction of $`W_N(\omega )`$ is $`W(𝔡,v,I)`$ where $`I=\{hH|yh=0\}`$ is as before.
The action of $`W(𝔡)`$ on $`v`$ is given by (8.8). In the notation of Section 8.4, we have for $`\alpha W(𝔡)`$:
$$\alpha v=(\mathrm{div}^{\mathrm{tr}\mathrm{ad}}(\alpha )1+\alpha )_Hv.$$
This shows that $`\alpha W(𝔡,v,I)`$ iff $`\mathrm{div}^{\mathrm{tr}\mathrm{ad}}(\alpha )1+\alpha HI`$.
Note that, since $`\omega 0`$, we have $`I𝐤=0`$. The intersection $`I(𝔡+𝐤)`$ is a Lie algebra. The projection $`\pi :(𝔡+𝐤)𝔡`$ is a Lie algebra homomorphism, which maps $`I(𝔡+𝐤)`$ isomorphically onto a subalgebra $`𝔡^{}`$ of $`𝔡`$. The inverse isomorphism $`𝔡^{}I(𝔡+𝐤)`$ is given by $`aa+\chi (a)`$ for some linear functional $`\chi :𝔡^{}𝐤`$ which vanishes on $`[𝔡^{},𝔡^{}]`$. Conversely, any such $`\chi `$ gives rise to an isomorphism as above.
For $`\beta H(𝔡+𝐤)`$, the equation $`\beta HI`$ is equivalent to the following two conditions: $`(\mathrm{id}\pi )(\beta )H𝔡^{}`$ and $`(\mathrm{id}\pi +\mathrm{id}\chi \pi )(\beta )=\beta `$. Applying this for $`\beta =\mathrm{div}^{\mathrm{tr}\mathrm{ad}}(\alpha )1+\alpha `$, we get $`(\mathrm{id}\pi )(\beta )=\alpha H𝔡^{}`$ and $`(\mathrm{id}\pi +\mathrm{id}\chi \pi )(\beta )=\alpha +(\mathrm{id}\chi )(\alpha )=\beta `$. The latter equation is equivalent to $`(\mathrm{id}\chi )(\alpha )=\mathrm{div}^{\mathrm{tr}\mathrm{ad}}(\alpha )1`$, i.e. to $`\mathrm{div}^{\mathrm{tr}\mathrm{ad}\chi }(\alpha )=0`$. We have proved:
###### Lemma 12.9.
The reconstruction of the Lie algebra $`S_N`$, provided with a transitive action of a Lie algebra $`𝔡`$, is a current Lie $`H`$-pseudoalgebra over $`S(𝔡^{},\chi ^{})`$ where $`𝔡^{}`$ is a Lie subalgebra of $`𝔡`$ and $`\chi ^{}`$ is a linear functional $`𝔡^{}𝐤`$ which vanishes on $`[𝔡^{},𝔡^{}]`$.
In fact, one can show that in this case $`𝔡^{}=𝔡`$, i.e., $`dimI(𝔡+𝐤)=N`$, but the above statement is sufficient for our purposes.
### 12.7. Reconstruction from $`K_N`$ and $`H_N`$
Now let $``$ be one of the Lie algebras $`K_N`$ or $`P_N`$, together with a transitive action of $`𝔡`$ on it. We know from Section 6 that as a topological vector space $``$ is homeomorphic to $`X`$. Since, by Proposition 6.9, all transitive actions of $`𝔡`$ on $`X`$ are equivalent, $``$ is isomorphic to the “canonical” $`H`$-module $`X`$, i.e., we may assume that the embedding $`𝔡W_N=\mathrm{Der}X`$ is the canonical one (Section 12.1). Then, by Proposition 11.3, the reconstruction of $``$ is isomorphic to $`H`$ as an $`H`$-module. In other words, $`𝒞()`$ is a free $`H`$-module of rank one.
###### Lemma 12.10.
The reconstruction of the Lie algebras $`K_N`$ and $`H_N`$, provided with a transitive action of a Lie algebra $`𝔡`$, is a free $`H`$-module of rank one.
###### Proof.
It is enough to show that the reconstruction functor $`𝒞`$ gives the same result on the topological $`H`$-modules $`X=P_N`$ and $`X/𝐤=H_N`$. In order to do so, we must show that every $`H`$-linear continuous homomorphism of $`X`$ to $`X/𝐤`$ can be obtained from a unique $`H`$-linear continuous homomorphism of $`X`$ to itself by composing with the canonical projection $`XX/𝐤`$.
Since $`X/𝐤`$ is linearly compact, by Remark 6.2 there is a bijection between $`\mathrm{Hom}_H^{\mathrm{cont}}(X,X/𝐤)`$ and $`\mathrm{Hom}_H((X/𝐤)^{},H)`$. The space $`(X/𝐤)^{}`$ is nothing but the augmentation ideal $`H_+=\mathrm{ker}\epsilon H`$. Therefore we are reduced to show that every $`H`$-linear map $`\varphi :H_+H`$ is a restriction of a unique $`H`$-linear map $`HH`$.
An $`H`$-linear $`\varphi :H_+H`$ is determined by its value on $`𝔡H_+`$. If $`a,b𝔡`$, then $`abba=[a,b]`$, hence $`a\varphi (b)b\varphi (a)=\varphi ([a,b])`$. Let $`d`$ be the maximal degree of $`\varphi (a)`$ for $`a𝔡`$. Then $`a\varphi (b)=b\varphi (a)`$ modulo $`\mathrm{F}^dH`$. This means that there exists some $`\alpha \mathrm{F}^{d1}H`$ such that for every $`a𝔡`$, $`\varphi (a)=a\alpha `$ modulo $`\mathrm{F}^{d1}H`$. Then the difference between $`\varphi `$ and right multiplication by $`\alpha `$ is still $`H`$-linear, and its maximal degree on elements from $`𝔡`$ is strictly less than $`d`$. The proof now follows by induction. ∎
All Lie pseudoalgebras that are free $`H`$-module of rank one are classified in Theorem 8.10: they are isomorphic to current pseudoalgebras over $`K(𝔡^{},\theta ^{})`$ or $`H(𝔡^{},\chi ^{},\omega ^{})`$.
## 13. Structure Theory of Lie Pseudoalgebras
### 13.1. Structural correspondence between a Lie pseudoalgebra and its annihilation algebra
Recall that a Lie $`H`$-pseudoalgebra $`L`$ is called finite if it is finitely generated as an $`H`$-module. If $`H`$ is Noetherian (e.g., $`H=U(𝔡)`$ for a finite-dimensional Lie algebra $`𝔡`$) and $`L`$ is finite, then $`L`$ is a Noetherian $`H`$-module, i.e., every increasing sequence of $`H`$-submodules of $`L`$ stabilizes.
For any two subspaces $`A`$ and $`B`$ of $`L`$, let
(13.1)
$$[A,B]=\{[a_xb]|aA,bB,xX\}.$$
Define the derived series of $`L`$ by $`L^{(0)}=L`$, $`L^{(1)}=[L,L]`$, $`L^{(n+1)}=[L^{(n)},L^{(n)}]`$. A Lie pseudoalgebra $`L`$ is called solvable if $`L^{(n)}=0`$ for some $`n`$. Similarly, define the central series of $`L`$ by $`L^0=L`$, $`L^1=[L,L]`$, $`L^{n+1}=[L^n,L]`$. The Lie pseudoalgebra $`L`$ is called nilpotent if $`L^n=0`$ for some $`n`$. As usual, $`L`$ is called abelian if $`[L,L]=0`$, i.e., if $`[ab]=0`$ for all $`a,bL`$.
A Lie pseudoalgebra $`L`$ is called simple if it contains no nontrivial ideals and is not abelian. Note that $`[L,L]`$ is an ideal of $`L`$, so in particular, $`[L,L]=L`$ if $`L`$ is simple. $`L`$ is called semisimple if it contains no nonzero abelian ideals.
We will show that, as in the Lie algebra case, $`L`$ is semisimple if and only if its radical is zero. Provided that it exists, we define the radical of $`L`$, $`\mathrm{Rad}L`$, to be its maximal solvable ideal. When $`H`$ is Noetherian and $`L`$ is finite, $`\mathrm{Rad}L`$ exists because of the Noetherianity of $`L`$ and part (ii) of the next lemma.
###### Lemma 13.1.
(i) If $`S`$ is a solvable ideal in $`L`$ and $`L/S`$ is solvable, then $`L`$ is solvable.
(ii) If $`S_1,S_2`$ are solvable ideals in $`L`$, then their sum $`S_1+S_2`$ is a solvable ideal.
(iii) $`L/\mathrm{Rad}L`$ is semisimple. $`L`$ is semisimple iff $`\mathrm{Rad}L=0`$.
###### Proof.
(i) is standard.
(ii) follows from (i) and the fact that $`(S_1+S_2)/S_1S_2/(S_1S_2)`$.
(iii) If $`L/\mathrm{Rad}L`$ has an abelian ideal $`I`$, then the preimage of $`I`$ under the natural projection $`LL/\mathrm{Rad}L`$ must be solvable and strictly bigger than $`\mathrm{Rad}L`$, which is a contradiction. ∎
It is easy to see, using (9.23, 9.24), that for any two subspaces $`A,BL`$, we have:
(13.2)
$$[X_HA,X_HB]=X_H[A,B]$$
as subspaces of $`𝒜(L)=X_HL`$. In particular, if $`I`$ is an ideal of $`L`$, then $`X_HI`$ is an ideal of $`𝒜(L)`$. We will call an ideal of $`𝒜(L)`$ regular if it is of the form $`X_HI`$ for some ideal $`I`$ of $`L`$.
###### Lemma 13.2.
Let $`L`$ be a Lie $`H`$-pseudoalgebra and $`IL`$ be an ideal. Then:
(i) $`X_HI=0`$ only if $`I`$ is central.
(ii) $`X_HI=𝒜(L)`$ only if $`[L,L]I`$.
###### Proof.
(i) has already been proved, when $`L`$ is finite, in Corollary 11.6 and Remark 11.8. In the general case, it can be deduced from Proposition 9.4. Let $`aI`$, then $`a_xx_Ha=0`$ for any $`xX`$. Hence the action of $`a_x`$ on $`L`$ is trivial, and by (9.22), $`[ab]=0`$ for any $`bL`$.
In order to prove (ii), notice that $`X_HL/I=0`$. Then build a Lie $`H`$-pseudoalgebra structure on $`\stackrel{~}{L}=LL/I`$ by letting $`L`$ act on the abelian ideal $`L/I`$ via the adjoint action. Then by part (i), $`L/I`$ is central in $`\stackrel{~}{L}`$, hence $`L`$ acts trivially on $`L/I`$. This means $`[L,L]I`$. ∎
Using this lemma and (13.2), it is easy to prove the next two results.
###### Proposition 13.3.
A Lie pseudoalgebra $`L`$ is solvable (respectively nilpotent) if and only if its annihilation Lie algebra $`𝒜(L)`$ is.
###### Proposition 13.4.
Let $`L`$ be a centerless Lie $`H`$-pseudoalgebra which is equal to its derived subalgebra $`[L,L]`$. Then $`L`$ is simple if either of the following conditions holds:
(i) $`𝒜(L)`$ has no nontrivial $`H`$-invariant ideals.
(ii) $`L`$ is finite or free, and $`𝒜(L)`$ has no non-central $`H`$-invariant ideals.
###### Proof.
(i) is immediate from Lemma 13.2.
Assume that (ii) holds but $`L`$ is not simple. Then $`𝒜(L)`$ has a nontrivial central regular ideal. If $`a_xx_Ha`$ is central in $`𝒜(L)`$ for every $`xX`$, then by (9.24) $`[a_xb]_y=0`$ for every $`bL`$, $`x,yX`$. When $`L`$ is either finite or free, $`l_y=0`$ for all $`yX`$ if and only if $`l=0`$ (cf. Corollary 11.6). Therefore $`[a_xb]=0`$ for all $`bL`$, $`xX`$, and by (9.22) we get $`[ab]=0`$ for any $`bL`$. Hence $`a=0`$. ∎
As an immediate consequence we obtain:
###### Corollary 13.5.
Let $`L`$ be a current Lie $`H`$-pseudoalgebra over a finite-dimensional simple Lie algebra or over one of the primitive pseudoalgebras of vector fields. Then $`L`$ is simple.
###### Proof.
It is easy to check that $`L`$ satisfies the assumptions of Proposition 13.4 (see Theorem 8.24). ∎
The following proposition will play an important role in the classification of finite simple Lie pseudoalgebras.
###### Proposition 13.6.
For any Lie $`H`$-pseudoalgebra $`L`$, any non-central $`H`$-invariant ideal $`J`$ of $`𝒜(L)`$ contains a nonzero regular ideal.
###### Proof.
Let $`\alpha J`$ be non-central. Assume that $`X_H\alpha l=0`$ for all $`lL`$. Note that by Proposition 9.4, we have: $`h(\alpha l)=(h_{(1)}\alpha )(h_{(2)}l)`$ for $`hH`$. This implies: $`(h\alpha )l=h_{(1)}\left(\alpha (h_{(2)}l)\right)`$, which gives $`X_H(h\alpha )l=0`$ for any $`hH`$, $`lL`$. Then we can use (9.30) to show that $`\alpha `$ is central in $`𝒜(L)`$, which is a contradiction.
Therefore, there is some $`lL`$ such that $`\alpha l=a`$ has a nonzero Fourier coefficient, i.e., $`X_Ha0`$. Since $`a_y=(\alpha l)_y=_i[h_i\alpha ,l_{yx_i}]`$, and $`J`$ is $`H`$-stable, we see that all Fourier coefficients of $`a`$ lie in $`J`$. Then, due to (9.24), all elements in the ideal $`(a)`$ of $`L`$ generated by $`a`$ have all of their Fourier coefficients in $`J`$, i.e., $`0X_H(a)J`$. ∎
### 13.2. Annihilation algebras of finite simple Lie $`U(𝔡)`$-pseudoalgebras
We will now approach the problem of classification of all finite simple Lie $`H`$-pseudoalgebras. In view of Kostant’s Theorem 2.1 and the results of Section 5, we will first restrict ourselves to the case when $`H`$ is the universal enveloping algebra of a Lie algebra $`𝔡`$. Moreover, we will assume that $`𝔡`$ is finite dimensional; in this case $`H=U(𝔡)`$ is filtered by finite-dimensional subspaces. The classification is done in two steps: the first one (done in this subsection) is classifying all Lie algebras that can arise as $`𝒜(L)`$ for some finite simple Lie $`H`$-pseudoalgebra $`L`$, the second step (done in the next subsection) involves a reconstruction of $`L`$ from its annihilation Lie algebra $`𝒜(L)`$ and the action of $`H`$ on it.
###### Theorem 13.7.
If $`L`$ is a finite simple Lie $`H=U(𝔡)`$-pseudoalgebra, then its annihilation Lie algebra $`𝒜(L)`$ is isomorphic (as a topological Lie algebra) to an irreducible central extension of a current Lie algebra $`𝒪_r\widehat{}𝔰`$ where $`𝔰`$ is a simple linearly compact Lie algebra of growth $`\mathrm{gw}𝔰=dim𝔡r`$.
###### Proof.
First of all, we observe that $`=𝒜(L)`$ is a linearly compact Lie algebra with respect to the topology defined in Section 7.4, see Proposition 7.12(ii). Consider the extended annihilation algebra $`^e=𝔡`$, obtained by letting $`𝔡`$ act on $`=𝒜(L)`$ according to its $`H=U(𝔡)`$-module structure.
###### Lemma 13.8.
$`^e=𝔡`$ is a linearly compact Lie algebra possessing a fundamental subalgebra, i.e., an open subalgebra containing no ideals of $`^e`$.
###### Proof.
Indeed, if $`L_0`$ is a finite-dimensional subspace of $`L`$ generating it over $`H`$, then because of (7.14), $`_i=\mathrm{F}_iX_HL_0`$ is a subalgebra of $``$ for $`is`$. None of $`_i`$ contains ideals of $`𝔡`$, since every such ideal is stable under the action of $`H`$ and $`H\mathrm{F}_iX=X`$, which implies $`H_i=`$. ∎
The center $`Z`$ of $``$ is an $`H`$-stable closed ideal. The quotient $`^e/Z=𝔡(/Z)`$ is a linearly compact Lie algebra possessing a fundamental subalgebra $`_s/(Z_s)`$. Theorem 13.7 will be deduced from Proposition 6.11 applied for $`\overline{}^e:=^e/Z`$.
By Proposition 13.6, any nonzero $`H`$-stable ideal of $`\overline{}:=/Z`$ contains the image of a nonzero regular ideal of $``$. Since $`L`$ is simple, this means that the only nonzero $`H`$-stable ideal of $`\overline{}`$ is the whole $`\overline{}`$. Then every nonzero ideal of $`\overline{}^e`$ contained in $`\overline{}`$ must equal $`\overline{}`$. Hence $`\overline{}`$ is a minimal closed ideal of a linearly compact Lie algebra satisfying the assumptions of Proposition 6.11(i), and is therefore of the form stated in part (ii) of this proposition.
Therefore, $``$ is a central extension of a current Lie algebra over a simple linearly compact Lie algebra. Moreover, $``$ equals its derived subalgebra (otherwise we would have a proper nonabelian subideal of $``$). Hence it is an irreducible central extension.
Consider the canonical filtration $`\mathrm{F}_n(𝒪_r\widehat{}𝔰):=_i\mathrm{F}_{ni}𝒪_r\widehat{}\mathrm{F}_i𝔰`$, where $`\mathrm{F}_i𝔰`$ is the canonical filtration of $`𝔰`$ defined in Section 6 (if $`dim𝔰<\mathrm{}`$ we put $`\mathrm{F}_i𝔰=0`$ for $`i0`$). Then the growth of $`𝒪_r\widehat{}𝔰`$ (with respect to this filtration) equals $`\mathrm{gw}𝒪_r+\mathrm{gw}𝔰=r+\mathrm{gw}𝔰`$. It is clear from Proposition 6.12 that any irreducible central extension of $`𝒪_r\widehat{}𝔰`$ has the same growth. On the other hand, with respect to the filtration defined by (7.11), the growth of $``$ is equal to $`N=dim𝔡`$ (see Proposition 7.15). We have to show that the two different filtrations give the same growth.
Recall that by Lemma 7.13, a sufficiently high power of any nonzero element $`a𝔡`$ maps any given open subspace of $``$ surjectively onto $``$. Then the same argument as in the proof of Lemma 12.1 shows that $`𝔡\mathrm{Der}`$ intersects $`\mathrm{F}_0(\mathrm{Der})`$ trivially, where $`\mathrm{F}_0(\mathrm{Der})`$ is induced by the canonical filtration on $`𝒪_r\widehat{}𝔰`$. This implies $`Nr+\mathrm{gw}𝔰`$.
To show the inverse inequality, note that since $`\mathrm{F}_0(𝒪_r\widehat{}𝔰)`$ is open in $`\overline{}=/Z𝒪_r\widehat{}𝔰`$, it contains some $`\overline{}_m:=_m/(Z_m)`$. Now (7.14) implies $`[\overline{}_i,\overline{}]\overline{}_{is1}`$, which together with (6.1) leads to $`\overline{}_{m+n(s+1)}\mathrm{F}_n(𝒪_r\widehat{}𝔰)`$ for all $`n0`$. This implies $`Nr+\mathrm{gw}𝔰`$.
This completes the proof of Theorem 13.7. ∎
In fact, the above arguments can be used to prove a stronger statement than Theorem 13.7.
###### Corollary 13.9.
Let $`L`$ be a finite Lie $`H`$-pseudoalgebra and $`M`$ be a minimal nonabelian ideal of $`L`$. Then the annihilation algebra of $`M`$ is one of the Lie algebras described in Theorem 13.7.
###### Proof.
The only place in the proof of Theorem 13.7 where we used the simplicity of $`L`$ was where we deduced that any nonzero regular ideal of $`𝒜(L)`$ must equal the whole $`𝒜(L)`$. This argument is modified as follows. Let $`J=X_HI`$ be a nonabelian regular ideal of $`𝒜(L)`$ contained in $`𝒜(M)`$. Then the minimality of $`M`$ implies that $`I=M`$ and $`J=𝒜(M)`$. The proof is concluded again by applying Proposition 6.11. ∎
### 13.3. Classification of finite simple Lie $`U(𝔡)`$-pseudoalgebras
We will call a pseudoalgebra of vector fields any subalgebra of the Lie pseudoalgebra $`W(𝔡)`$. As in Section 8, a pseudoalgebra of vector fields is called primitive if it is one of the following: $`W(𝔡)`$, $`S(𝔡,\chi )`$, $`H(𝔡,\chi ,\omega )`$ or $`K(𝔡,\theta )`$ (then its annihilation algebra $`𝒜(L)`$ is isomorphic to one of the primitive Lie algebras $`W_N`$, $`S_N`$, $`P_N`$ or $`K_N`$).
The following is the main theorem of this section.
###### Theorem 13.10.
Let $`H=U(𝔡)`$ be the universal enveloping algebra of a finite-dimensional Lie algebra $`𝔡`$. Then any finite simple Lie $`H`$-pseudoalgebra $`L`$ is isomorphic to a current pseudoalgebra over a finite-dimensional simple Lie algebra or over one of the primitive pseudoalgebras of vector fields.
Explicitly, $`L\mathrm{Cur}_H^{}^HL^{}`$, where $`H^{}=U(𝔡^{})`$, $`𝔡^{}`$ is a subalgebra of $`𝔡`$, and $`L^{}`$ is one of the following:
(a) $`L^{}`$ is a finite-dimensional simple Lie algebra and $`𝔡^{}=0`$;
(b) $`L^{}=W(𝔡^{})`$, $`𝔡^{}`$ is arbitrary;
(c) $`L^{}=S(𝔡^{},\chi ^{})`$, where $`𝔡^{}`$ is arbitrary and $`\chi ^{}(𝔡^{})^{}`$ is such that $`\chi ^{}([𝔡^{},𝔡^{}])=0`$;
(d) $`L^{}=H(𝔡^{},\chi ^{},\omega ^{})`$, where $`N^{}=dim𝔡^{}`$ is even, $`\chi ^{}`$ is as in (c), and $`\omega ^{}^2(𝔡^{})^{}`$ is such that $`(\omega ^{})^{N^{}/2}0`$ and $`\mathrm{d}\omega ^{}+\chi ^{}\omega ^{}=0`$;
(e) $`L^{}=K(𝔡^{},\theta ^{})`$, where $`N^{}=dim𝔡^{}`$ is odd and $`\theta ^{}(𝔡^{})^{}`$ is such that $`\theta ^{}(\mathrm{d}\theta ^{})^{(N^{}1)/2}0`$.
###### Proof.
By Theorem 13.7, the annihilation algebra $``$ of $`L`$ is an irreducible central extension of a current Lie algebra $`\overline{}=𝒪_r\widehat{}\overline{𝔰}`$, where $`\overline{𝔰}`$ is a simple linearly compact Lie algebra of growth $`N^{}=Nr`$. We have surjective maps
(13.3)
$$𝒪_r\widehat{}𝔰𝒪_r\widehat{}\overline{𝔰},$$
where $`𝔰`$ is the universal central extension of $`\overline{𝔰}`$. By Theorem 6.8, $`\overline{𝔰}`$ is either finite dimensional (when $`N^{}=0`$) or one of the Lie algebras $`W_N^{}`$, $`S_N^{}`$, $`H_N^{}`$ or $`K_N^{}`$. By Proposition 6.12, we have $`𝔰=\overline{𝔰}`$ in all cases, except $`\overline{𝔰}=H_N^{}`$ in which case the center of $`𝔰=P_N^{}`$ is 1-dimensional.
Note that $`\mathrm{Der}𝔰=\mathrm{Der}\overline{𝔰}`$, and therefore, by Proposition 6.12, we have $`\mathrm{Der}(𝒪_r\widehat{}𝔰)=\mathrm{Der}(𝒪_r\widehat{}\overline{𝔰})`$. This implies $`\mathrm{Der}(𝒪_r\widehat{}𝔰)=\mathrm{Der}=\mathrm{Der}(𝒪_r\widehat{}\overline{𝔰})`$. Then the action of $`𝔡`$ on $``$ induces actions on $`𝒪_r\widehat{}𝔰`$ and $`𝒪_r\widehat{}\overline{𝔰}`$. The argument from the proof of Lemma 12.1 shows that these actions are transitive.
Now, let us apply the reconstruction functor $`𝒞`$ to the maps in (13.3). By Theorem 12.2, $`𝒞(𝒪_r\widehat{}𝔰)\mathrm{Cur}_H^{}^H𝒞(𝔰)`$, and $`S:=𝒞(𝔰)`$ is one of the Lie pseudoalgebras described in (a)–(e) above. Moreover, by Lemma 12.10, we have $`𝒞(\overline{𝔰})𝒞(𝔰)=S`$, and hence $`𝒞(𝒪_r\widehat{}\overline{𝔰})\mathrm{Cur}_H^{}^HS`$. We therefore obtain $`H`$-linear maps $`\mathrm{Cur}_H^{}^HS\widehat{L}\mathrm{Cur}_H^{}^HS`$ whose composition is the identity. Hence $`\widehat{L}:=𝒞()`$ is isomorphic to $`\mathrm{Cur}_H^{}^HS`$, which is a simple Lie pseudoalgebra (Corollary 13.5).
The homomorphism $`\mathrm{\Phi }:L\widehat{L}`$ given by (11.6) is injective because $`L`$ is centerless (Remark 11.8). The action of $`\widehat{L}`$ on $`L`$ built in Section 11.4 shows that the image of $`\mathrm{\Phi }`$ is an ideal of $`\widehat{L}`$. Since $`\widehat{L}`$ is simple, it follows that $`\mathrm{\Phi }`$ is an isomorphism. ∎
Corollary 13.9 and the above proof imply the following result.
###### Corollary 13.11.
Let $`L`$ be a finite Lie pseudoalgebra and $`M`$ be a minimal nonabelian ideal of $`L`$. Then $`M`$ is a simple Lie pseudoalgebra.
###### Lemma 13.12.
If $`L`$ is a centerless Lie pseudoalgebra, then any nonzero finite ideal of $`L`$ contains a nonzero minimal ideal.
###### Proof.
By Zorn’s Lemma, it is enough to show that $`I_\alpha 0`$ for any collection of finite ideals $`\{I_\alpha \}_{\alpha A}`$ of $`L`$ such that $`I_\alpha I_\beta `$ for $`\alpha <\beta `$, where $`A`$ is a totally ordered index set. Assume that $`I_\alpha =0`$. Then there is a chain of ideals $`\{I_\alpha \}_{\alpha A^{}}`$ $`(A^{}A)`$ of the same rank whose intersection is zero. Fix some $`\alpha _0A^{}`$. Then for any $`\beta A^{}`$, $`\beta <\alpha _0`$, the module $`I_{\alpha _0}/I_\beta `$ is torsion, so by Corollary 10.17, $`L`$ acts trivially on it. This implies $`[L,I_{\alpha _0}]I_\beta `$ for each such $`\beta `$, hence $`I_{\alpha _0}`$ is central. ∎
### 13.4. Derivations of finite simple Lie $`U(𝔡)`$-pseudoalgebras
We will determine all derivations of a finite simple Lie $`H=U(𝔡)`$-pseudoalgebra $`L`$ (see Definition 10.8).
First let us consider the case when $`L=\mathrm{Cur}𝔤:=H𝔤`$ is a current pseudoalgebra over a finite-dimensional Lie algebra $`𝔤`$. The Lie pseudoalgebra $`W(𝔡)`$ acts on $`L`$ by just acting on the first factor in $`H𝔤`$ (cf. (8.4)):
(13.4)
$$(fa)(gb)=(fga)_H(1b),f,gH,a𝔡,b𝔤.$$
We also have an embedding $`\mathrm{Cur}\mathrm{Der}𝔤\mathrm{Der}L`$. The image of $`\mathrm{Cur}\mathrm{Der}𝔤`$ in $`\mathrm{Der}L`$ is normalized by that of $`W(𝔡)`$, and the two form a semidirect sum $`W(𝔡)\mathrm{Cur}\mathrm{Der}𝔤`$ which as an $`H`$-module is isomorphic to $`H(𝔡\mathrm{Der}𝔤)`$.
###### Proposition 13.13.
For any simple finite-dimensional Lie algebra $`𝔤`$, we have $`\mathrm{Der}\mathrm{Cur}𝔤=W(𝔡)\mathrm{Cur}𝔤`$.
###### Proof.
By Lemma 10.9(iii), the annihilation algebra $`𝒜(\mathrm{Der}\mathrm{Cur}𝔤)\mathrm{Der}𝒜(\mathrm{Cur}𝔤)=\mathrm{Der}(X𝔤)`$. By Proposition 6.12(ii), the latter is isomorphic to $`W_N1+𝒪_N𝔤`$, since $`X𝒪_N`$. Then $`𝒞𝒜(\mathrm{Der}\mathrm{Cur}𝔤)𝒞(\mathrm{Der}(X𝔤))=W(𝔡)+\mathrm{Cur}𝔤`$ (see Theorem 12.2). Now by Lemma 11.15, $`\mathrm{Der}\mathrm{Cur}𝔤𝒞𝒜(\mathrm{Der}\mathrm{Cur}𝔤)W(𝔡)+\mathrm{Cur}𝔤`$. ∎
A similar argument as in the proof of the proposition shows that $`\mathrm{Der}L=L`$ when $`L`$ is one of the primitive pseudoalgebras of vector fields $`W(𝔡)`$, $`S(𝔡,\chi )`$, $`H(𝔡,\chi ,\omega )`$ or $`K(𝔡,\theta )`$. In fact, the same holds when $`L`$ is a current pseudoalgebra over one of them.
###### Proposition 13.14.
Let $`L`$ be a simple pseudoalgebra of vector fields $`(`$i.e., $`L`$ is a current pseudoalgebra over one of the primitive ones$`)`$. Then $`\mathrm{Der}L=L`$.
###### Proof.
Let $`L`$ be a current pseudoalgebra over $`L^{}`$, and $`L^{}W(𝔡^{})`$ be one of the primitive pseudoalgebras of vector fields, where $`𝔡^{}`$ is a Lie subalgebra of $`𝔡`$. Then, by Theorem 8.24, the annihilation algebra $`=𝒜(L)`$ is a current Lie algebra over $`^{}=𝒜(L^{})`$: $`=𝒪_r\widehat{}^{}`$, and $`^{}`$ is isomorphic to $`W_N^{}`$, $`S_N^{}`$, $`P_N^{}`$ or $`K_N^{}`$, where $`N^{}=dim𝔡^{}=Nr`$, $`N=dim𝔡`$.
As in the proof of Proposition 13.13, we have $`\mathrm{Der}L𝒞𝒜(\mathrm{Der}L)𝒞(\mathrm{Der})`$. By Proposition 6.12, we have: $`\mathrm{Der}=W_r1+𝒪_r\widehat{}\mathrm{Der}^{}`$, and $`\mathrm{Der}^{}=`$ $`W_N^{}`$, $`CS_N^{}`$, $`CH_N^{}`$ or $`K_N^{}`$ is a Lie subalgebra of $`W_N^{}`$. In particular, we see that $`\mathrm{Der}W_N`$, and hence $`𝒞(\mathrm{Der})`$ is a subalgebra of $`W(𝔡)`$.
So, we have: $`\mathrm{Der}LW(𝔡)`$, $`L=\mathrm{Cur}L^{}\mathrm{Cur}W(𝔡^{})H𝔡^{}W(𝔡)`$. Take any two nonzero elements $`aW(𝔡)`$, $`bH𝔡^{}`$. Then we claim that $`[ab]H𝔡^{}`$ implies $`aH𝔡^{}`$. This follows easily from the definition (8.3), using that $`𝔡^{}`$ is a subalgebra of $`𝔡`$ (see the proof of Proposition 13.16 below for a similar argument).
Therefore $`\mathrm{Der}L\mathrm{Cur}W(𝔡^{})`$. For $`a\mathrm{Cur}W(𝔡^{})`$, we can write $`a=h_i_H^{}a_i`$ for some $`a_iW(𝔡^{})`$ and $`h_iH`$ such that the classes $`h_iH^{}`$ are linearly independent in $`H/H^{}`$. Then if $`a\mathrm{Der}L`$, it is easy to see that all $`a_i`$ must belong to $`\mathrm{Der}L^{}`$. Hence $`\mathrm{Der}L=\mathrm{Cur}\mathrm{Der}L^{}`$. But $`\mathrm{Der}L^{}𝒞(\mathrm{Der}^{})=L^{}`$, so $`\mathrm{Der}L=\mathrm{Cur}L^{}=L`$. ∎
### 13.5. Finite semisimple Lie $`U(𝔡)`$-pseudoalgebras
Recall that a Lie $`H`$-pseudoalgebra $`L`$ is called semisimple if it contains no nonzero abelian ideals. Let $`H=U(𝔡)`$, for a finite-dimensional Lie algebra $`𝔡`$.
If $`𝔤`$ is a simple finite-dimensional Lie algebra, then by Proposition 13.13, we have $`\mathrm{Der}\mathrm{Cur}𝔤=W(𝔡)\mathrm{Cur}𝔤`$. It is easy to see that for any subalgebra $`A`$ of the Lie pseudoalgebra $`W(𝔡)`$, the Lie pseudoalgebra $`A\mathrm{Cur}𝔤`$ is semisimple. Indeed, assume that $`IA\mathrm{Cur}𝔤`$ is an abelian ideal. Then $`I\mathrm{Cur}𝔤`$ is an abelian ideal in $`\mathrm{Cur}𝔤`$, hence $`I\mathrm{Cur}𝔤=0`$. But this is impossible unless $`I=0`$ because the pseudobracket of any element from $`(W(𝔡)+\mathrm{Cur}𝔤)\mathrm{Cur}𝔤`$ with elements from $`\mathrm{Cur}𝔤`$ gives nonzero elements from $`\mathrm{Cur}𝔤`$ (see (13.4)). Note that this argument implies that any nonzero ideal of $`A\mathrm{Cur}𝔤`$ contains $`\mathrm{Cur}𝔤`$.
Now we can classify all finite semisimple Lie $`U(𝔡)`$-pseudoalgebras.
###### Theorem 13.15.
Any finite semisimple Lie $`U(𝔡)`$-pseudoalgebra $`L`$ is a direct sum of finite simple Lie pseudoalgebras $`(`$described by Theorem 13.10$`)`$ and of pseudoalgebras of the form $`A\mathrm{Cur}𝔤`$, where $`A`$ is a subalgebra of $`W(𝔡)`$ and $`𝔤`$ is a simple finite-dimensional Lie algebra.
###### Proof.
Consider the set $`\{M_i\}`$ of all minimal nonzero ideals of $`L`$. This set is nonempty by Lemma 13.12, and finite because $`L`$ is a Noetherian $`H`$-module. The adjoint action of $`L`$ on $`M_i`$ gives a homomorphism of Lie pseudoalgebras $`L\mathrm{Der}M_i`$, cf. Lemma 10.9(ii).
We claim that the direct sum of these homomorphisms is an injective map. Indeed, let $`NL`$ be the set of all elements that act trivially on all $`M_i`$. This set is an ideal of $`L`$. If it is nonzero it must contain some minimal ideal $`M_i`$. But then this $`M_i`$ is abelian, which contradicts the semisimplicity of $`L`$.
Therefore we have embeddings $`M_iL\mathrm{Der}M_i`$. By Corollary 13.11 all $`M_i`$ are simple Lie pseudoalgebras. If $`M_i`$ is not a current pseudoalgebra over a finite-dimensional Lie algebra, then by Proposition 13.14, $`\mathrm{Der}M_i=M_i`$. For $`M_i=\mathrm{Cur}𝔤`$, we have $`\mathrm{Der}\mathrm{Cur}𝔤=W(𝔡)\mathrm{Cur}\mathrm{Der}𝔤`$. Any subalgebra of $`W(𝔡)\mathrm{Cur}𝔤`$ containing $`\mathrm{Cur}𝔤`$ is of the form $`A\mathrm{Cur}𝔤`$, where $`A`$ is a subalgebra of $`W(𝔡)`$. ∎
Recall that a pseudoalgebra of vector fields is any subalgebra of the Lie pseudoalgebra $`W(𝔡)`$.
###### Proposition 13.16.
For any two nonzero elements $`a,bW(𝔡)`$, we have $`[ab]0`$. In particular, $`W(𝔡)`$ does not contain nonzero abelian elements, i.e., elements $`a`$ such that $`[aa]=0`$.
###### Proof.
Let us write
$$a=_ih_i_i,b=_jk_j_j,$$
where $`h_i,k_jH`$ and $`\{_i\}`$ is a basis of $`𝔡`$. Denote by $`m`$ (respectively $`n`$) the maximal degree of the $`h_i`$ (respectively $`k_j`$). We have (cf. (8.3)):
$`[ab]`$ $`=_{i,j}(h_ik_j)_H(1[_i,_j])`$
$`_{i,j}(h_ik_j_i)_H(1_j)+_{i,j}(h_i_jk_j)_H(1_i).`$
Assume that $`[ab]=0`$. Notice that only the third summand contains coefficients from $`HH`$ of degree $`(m+1,n)`$, hence it must be zero modulo $`\mathrm{F}^mH\mathrm{F}^nH`$. Since the $`_i`$ are linearly independent, the same is true for each term $`_jh_i_jk_j`$. If we choose $`i`$ such that $`h_i`$ is of degree exactly $`m`$, we get a contradiction. ∎
###### Corollary 13.17.
A finite Lie $`U(𝔡)`$-pseudoalgebra $`L`$ contains no nonzero abelian elements iff it is a direct sum of pseudoalgebras of vector fields.
###### Proof.
Assume that $`L`$ is not a direct sum of pseudoalgebras of vector fields. If $`L`$ is not semisimple, then $`\mathrm{Rad}L`$ contains nonzero abelian elements. If $`L`$ is semisimple, Theorem 13.15 implies that $`L`$ contains a subalgebra of the form $`A\mathrm{Cur}𝔤`$ with $`𝔤0`$, and therefore contains nonzero abelian elements (for example $`1a`$ for any $`a𝔤`$).
The converse statement follows from Proposition 13.16. ∎
###### Theorem 13.18.
Any pseudoalgebra of vector fields is simple.
###### Proof.
By Proposition 13.16, a pseudoalgebra $`L`$ of vector fields does not contain nonzero abelian elements, and hence is semisimple. Then, by Theorem 13.15, $`L`$ is a direct sum of finite simple Lie pseudoalgebras and of pseudoalgebras of the form $`A\mathrm{Cur}𝔤`$. In fact, there is only one summand, as $`[ab]0`$ for any two nonzero elements $`a,bW(𝔡)`$. Furthermore, $`L`$ cannot be of the form $`A\mathrm{Cur}𝔤`$ with $`𝔤0`$, because $`\mathrm{Cur}𝔤`$ contains nonzero abelian elements. ∎
###### Corollary 13.19.
Any finite semisimple Lie $`U(𝔡)`$-pseudoalgebra $`L`$ is a direct sum of pseudoalgebras of the form $`A\mathrm{Cur}𝔤`$, where $`A`$ is either $`0`$ or one of the simple pseudoalgebras of vector fields $`(`$described by Theorem 13.10$`)`$, and $`𝔤`$ is either $`0`$ or a simple finite-dimensional Lie algebra.
We can also describe all ideals of a finite semisimple Lie pseudoalgebra $`L`$. By the above corollary, it is enough to consider the case $`L=A\mathrm{Cur}𝔤`$ with $`A0`$, $`𝔤0`$.
###### Proposition 13.20.
Let $`L=A\mathrm{Cur}𝔤`$ where $`A`$ is a pseudoalgebra of vector fields and $`𝔤`$ is a simple finite-dimensional Lie algebra. Then the only nonzero proper ideal of $`L`$ is $`\mathrm{Cur}𝔤`$.
###### Proof.
We have already noticed (see the paragraph before Theorem 13.15) that any nonzero ideal $`I`$ of $`L`$ contains $`\mathrm{Cur}𝔤`$. Then $`I/\mathrm{Cur}𝔤`$ is an ideal of $`L/\mathrm{Cur}𝔤A`$, but $`A`$ is simple by Theorem 13.18. ∎
### 13.6. Homomorphisms between finite simple Lie $`U(𝔡)`$-pseudoalgebras
In this subsection, $`H=U(𝔡)`$ is again the universal enveloping algebra of a finite-dimensional Lie algebra $`𝔡`$.
###### Theorem 13.21.
For any finite-dimensional Lie algebra $`𝔤`$ and any pseudoalgebra of vector fields $`L`$, there are no nontrivial homomorphisms between $`L`$ and $`\mathrm{Cur}𝔤`$.
###### Proof.
Any homomorphism $`\mathrm{Cur}𝔤L`$ leads to abelian elements in $`L`$, and therefore is zero (see Proposition 13.16).
Let $`f`$ be a homomorphism from $`L`$ to $`\mathrm{Cur}𝔤`$. Then $`f`$ induces a homomorphism of Lie algebras $`𝒜(f):𝒜(L)𝒜(\mathrm{Cur}𝔤)`$. By Theorem 13.18, $`L`$ is simple, so $`L=\mathrm{Cur}_H^{}^HL^{}`$ where $`L^{}`$ is a primitive $`H^{}`$-pseudoalgebra of vector fields ($`H^{}=U(𝔡^{})`$ and $`𝔡^{}`$ is a subalgebra of $`𝔡`$). By Theorem 8.24, the annihilation algebra $`=𝒜(L)`$ is isomorphic to a current Lie algebra $`𝒪_r\widehat{}^{}`$ over $`^{}=𝒜(L^{})`$. Moreover, the quotient of $`^{}`$ by its center is infinite dimensional and simple. On the other hand, the annihilation algebra $`𝒜(\mathrm{Cur}𝔤)X𝔤`$ is a current Lie algebra over $`𝔤`$, which is a projective limit of finite-dimensional Lie algebras $`(X/\mathrm{F}_nX)𝔤`$. Hence the adjoint action of $`^{}1^{}`$ on $``$ maps trivially to each of them via $`𝒜(f)`$. But since $`[^{},]=`$, this implies that each $`(X/\mathrm{F}_nX)𝔤`$ is trivial. Therefore $`𝒜(f)=0`$, and by Corollary 11.7, we get $`f=0`$. ∎
###### Theorem 13.22.
Let $`𝔤`$ and $`𝔥`$ be finite-dimensional simple Lie algebras. Then any isomorphism $`f:\mathrm{Cur}𝔤\stackrel{}{}\mathrm{Cur}𝔥`$ maps $`1𝔤`$ onto $`1𝔥`$, and thus is induced by some isomorphism of Lie algebras $`𝔤\stackrel{}{}𝔥`$. In particular, $`\mathrm{Aut}\mathrm{Cur}𝔤\mathrm{Aut}𝔤`$.
Recall that $`𝒜(\mathrm{Cur}𝔤)X𝔤`$ is a current Lie algebra. In the proof of the theorem we are going to use the following lemma.
###### Lemma 13.23.
Let $`𝔤`$ be a finite-dimensional simple Lie algebra, and $`R`$ be a commutative associative algebra. Then all ideals of $`R𝔤`$ are of the form $`I𝔤`$ where $`I`$ is an ideal of $`R`$.
###### Proof.
As a $`𝔤`$-module, $`R𝔤`$ is isomorphic to a direct sum of several copies of $`𝔤`$. Any ideal $`J`$ of $`R𝔤`$ is in particular a $`𝔤`$-module, hence it is spanned over $`𝐤`$ by elements of the form $`raJ`$ where $`rR`$ and $`a`$ is a root vector in $`𝔤`$. If $`r0`$ is such that $`raJ`$ for some nonzero $`a𝔤`$, then $`r𝔤J`$, since $`\{a𝔤|raJ\}`$ is an ideal of $`𝔤`$ and $`𝔤`$ is simple. Setting $`I=\{rR|r𝔤J\}`$, we see that $`I`$ is an ideal of $`R`$ and $`J=I𝔤`$. ∎
###### Proof of Theorem 13.22.
Define a map $`\rho :𝒜(\mathrm{Cur}𝔥)𝔥`$ by the formula:
(13.5)
$$\rho (x_H(1a))=x,1a,xX,a𝔥.$$
Then for $`a=_ih_ia_i\mathrm{Cur}𝔥=H𝔥`$, we have:
(13.6)
$$\rho (x_Ha)=_iS(x),h_ia_i.$$
It is easy to see that $`\rho `$ is a surjective Lie algebra homomorphism.
Any isomorphism $`f:\mathrm{Cur}𝔤\stackrel{}{}\mathrm{Cur}𝔥`$ induces an isomorphism of Lie algebras $`\phi =𝒜(f):𝒜(\mathrm{Cur}𝔤)\stackrel{}{}𝒜(\mathrm{Cur}𝔥)`$. By Lemma 13.23, $`\mathrm{ker}\rho \phi =I𝔤`$ for some proper ideal $`I`$ of $`X`$. Recall that $`X`$ is isomorphic as a topological algebra to $`𝒪_N=𝐤[[t_1,\mathrm{},t_N]]`$ ($`N=dim𝔡`$), and $`𝒪_N`$ has a unique maximal ideal $`M_0=(t_1,\mathrm{},t_N)`$. Noting that $`M_0`$ corresponds to $`\mathrm{F}_0X:=\{xX|x,1=0\}`$ via the isomorphism $`X𝒪_N`$, we deduce that $`I\mathrm{F}_0X`$. If $`I\mathrm{F}_0X`$, then $`(\mathrm{F}_0X/I)𝔤`$ is a nontrivial ideal of $`(X/I)𝔤𝔥`$, which is impossible because $`𝔥`$ is simple. It follows that $`\mathrm{ker}\rho \phi =\mathrm{F}_0X𝔤`$.
Now fix $`a𝔤`$ and write $`f(a)=_ih_ia_i`$ for some $`h_iH`$ and linearly independent $`a_i𝔥`$. Assume that, say, $`h_1𝐤=\mathrm{F}^0H`$. Then we can find $`x\mathrm{F}_0X`$ such that $`S(x),h_10`$. Then, by (13.6), the element $`xa\mathrm{F}_0X𝔤`$ is mapped by $`\rho \phi `$ to $`_iS(x),h_ia_i0`$, which is a contradiction. This shows that $`f(a)1𝔥`$, completing the proof. ∎
We turn now to the description of subalgebras of $`W(𝔡)`$. Recall that the Lie pseudoalgebra $`W(𝔡)`$ acts on $`H=U(𝔡)`$ by (8.4). Hence any homomorphism of Lie pseudoalgebras $`LW(𝔡)`$ gives rise to a structure of an $`L`$-module on $`H`$.
Let us first consider the case when $`L`$ is a free $`H`$-module of rank one: $`L=He`$ with a pseudobracket $`[ee]=\alpha _He`$, $`\alpha HH`$. Let $`M=Hm`$ be an $`L`$-module, with action $`em=\beta _Hm`$, $`\beta HH`$. We already know (Lemma 4.2) that $`\alpha `$ must be of the form $`\alpha =r+s11s`$ where $`r𝔡𝔡,s𝔡`$. Moreover $`r`$ and $`s`$ satisfy equations (4.3) and (4.4). Furthermore, $`\beta `$ defines a representation of $`L`$ if and only if it satisfies the following equation in $`HHH`$ (cf. Proposition 4.1):
(13.7)
$$(1\beta )(\mathrm{id}\mathrm{\Delta })(\beta )(\sigma \mathrm{id})\left((1\beta )(\mathrm{id}\mathrm{\Delta })(\beta )\right)=(\alpha 1)(\mathrm{\Delta }\mathrm{id})(\beta ).$$
###### Proposition 13.24.
If $`L=He`$ with $`[ee]=\alpha _He`$, $`\alpha =r+s11s`$, then the only nonzero homomorphism $`LW(𝔡)`$ is given by $`er+1s`$.
###### Proof.
The statement of the proposition is equivalent to saying that all solutions $`\beta `$ of (13.7) with $`\beta H𝔡`$ are either trivial or of the form $`\beta =r1s`$. It is easy to check that the latter is indeed a solution (cf. Lemma 8.8).
Let us choose a basis $`\{_i\}`$ of $`𝔡`$, and write $`\beta =h^i_i`$ and $`r=_{i,j}r^{ij}_i_j`$ for some $`h^iH`$, $`r^{ij}𝐤`$. We will assume throughout the proof that $`\beta 0`$, and denote by $`d`$ the maximal degree of the $`h^i`$. Substituting the above expressions for $`\alpha `$ and $`\beta `$ in (13.7), we get:
$$\begin{array}{cc}\hfill \underset{i,j}{}h^j& h^i[_i,_j]+\underset{i,j}{}(h^jh^i_jh^i_jh^j)_i\hfill \\ & =\underset{i,j,k}{}r^{ij}_ih_{(1)}^k_jh_{(2)}^k_k+\underset{k}{}(sh_{(1)}^kh_{(2)}^kh_{(1)}^ksh_{(2)}^k)_k.\hfill \end{array}$$
If $`d>1`$, expressing all $`h^k`$ in the Poincaré–Birkhoff–Witt basis relative to the basis $`\{_i\}`$, we see that $`HH`$-coefficients of degree $`2d+1`$ in the second summation in the left-hand side cannot cancel with terms from other summations, which only contribute lower degree terms. Therefore
$$\underset{i,j}{}(h^jh^i_jh^i_jh^j)_i=0mod\mathrm{F}^{2d}(HH)𝔡.$$
This implies that $`_jh^jh^i_j=_jh^i_jh^jmod\mathrm{F}^{2d}(HH)`$ for every $`i`$, which gives a contradiction, since we can choose $`h^i`$ to have degree exactly $`d`$. So $`d1`$, and we can write $`\beta =_{i,j}\beta ^{ij}_i_j+1t`$ with $`\beta ^{ij}𝐤`$, $`t𝔡`$.
Substituting this into (13.7) and comparing degree four terms we get $`\beta ^{ij}\beta ^{kl}=r^{ij}\beta ^{kl}`$ for all $`i,j,k,l`$. Since $`\beta 0`$ we conclude that $`\beta ^{ij}=r^{ij}`$ for all $`i,j`$. We are only left with showing that $`t=s`$. Substitute $`\alpha =r+s_1s_2`$ and $`\beta =r+t_2`$ into (13.7), and then use (4.4) to obtain:
(13.8)
$$\begin{array}{cc}\hfill r_{12}(s_3+t_3)+[t_3,r_{13}r_{23}]+r_{23}t_2r_{13}t_1& s_1r_{13}+s_2r_{23}\hfill \\ & =t_3(t_1t_2+s_1s_2).\hfill \end{array}$$
Notice that $`r_{12}(s_3+t_3)`$ is the only term lying in $`𝔡𝔡𝔡`$ and everything else belongs to $`HH𝐤+H𝐤H+𝐤HH`$. Hence $`r_{12}(s_3+t_3)=0`$, which is only possible if $`r=0`$ or $`s+t=0`$. In the latter case we are done. In the former, the left-hand side of (13.8) becomes zero, and $`t0`$ since $`\beta 0`$. Thus $`t_1+s_1t_2s_2=0`$ and $`t+s=0`$. ∎
Proposition 13.24 shows that nonabelian Lie pseudoalgebras that are free of rank one over $`H`$ embed uniquely in $`W(𝔡)`$. We will show that the other simple pseudoalgebras of vector fields are spanned as Lie pseudoalgebras by subalgebras of rank one, and therefore also embed uniquely in $`W(𝔡)`$. Recall that any pseudoalgebra of vector fields is in fact simple (Theorem 13.18).
###### Theorem 13.25.
(i) For any subalgebra $`L`$ of $`W(𝔡)`$, there is a unique nonzero homomorphism $`LW(𝔡)`$.
(ii) There is at most one nonzero homomorphism between any two pseudoalgebras of vector fields.
###### Proof.
Part (ii) is an immediate consequence of (i) and Theorem 13.18.
By the above remarks, it remains to prove (i) in the cases when $`L`$ is a current pseudoalgebra over either $`W(𝔡^{})`$ or $`S(𝔡^{},\chi ^{})`$, where $`𝔡^{}`$ is a subalgebra of $`𝔡`$.
In the former case ($`L=\mathrm{Cur}W(𝔡^{}):=H_H^{}(H^{}𝔡^{})`$, $`H^{}=U(𝔡^{})`$), note that $`L`$ is spanned over $`H=U(𝔡)`$ by elements $`\stackrel{~}{a}=1_H^{}(1a)`$ for $`a𝔡^{}`$. Then $`[\stackrel{~}{a}\stackrel{~}{a}]=(a11a)_H\stackrel{~}{a}`$, and by Proposition 13.24 we know that the only nonzero homomorphism of the Lie $`H`$-pseudoalgebra $`H\stackrel{~}{a}`$ to $`W(𝔡)`$ maps $`\stackrel{~}{a}`$ to $`1a`$. Hence any embedding of $`L`$ in $`W(𝔡)`$ maps each $`\stackrel{~}{a}`$ to the corresponding element $`1a`$ of $`W(𝔡)`$.
Now let $`L`$ be a current pseudoalgebra over $`S(𝔡^{},\chi ^{})`$. We will give the proof in the case when $`L=S(𝔡,\chi )`$, the case of currents being completely analogous. We are going to make use of the following lemma.
###### Lemma 13.26.
If $`𝔡`$ is a finite-dimensional Lie algebra of dimension $`N>1`$, then there exist $`2`$-dimensional subalgebras $`𝔡_i`$ $`(i=1,\mathrm{},N1)`$ such that $`dim_{i=1}^r𝔡_i=r+1`$ for every $`r=1,\mathrm{},N1`$.
###### Proof.
If $`𝔡`$ has a semisimple element $`h`$, we complement it to a basis of $`\mathrm{ad}h`$ eigenvectors $`\{h,h_1,\mathrm{},h_{N1}\}`$. The subalgebras $`𝔡_i=𝐤h+𝐤h_i`$ then satisfy the statement of the lemma.
If $`𝔡`$ has no semisimple elements, then from Levi’s theorem we know that $`𝔡`$ must be solvable. In this case it has a $`1`$-dimensional ideal $`𝐤h`$. Complementing $`h`$ to a basis $`\{h,h_1,\mathrm{},h_{N1}\}`$, we conclude as before. ∎
Now consider a $`2`$-dimensional subalgebra of $`𝔡`$ with basis $`\{a,b\}`$. Then the element $`e_{ab}S(𝔡,\chi )`$ from Proposition 8.4 depends on the choice of basis only up to multiplication by a nonzero element of $`𝐤`$. Moreover, the $`H`$-span of this element is a (free) rank one subalgebra of $`S(𝔡,\chi )`$, as can be easily checked (cf. Remark 8.5).
Let $`S_i`$ be the rank one subalgebras of $`S(𝔡,\chi )`$ associated as above with the 2-dimensional subalgebras $`𝔡_i`$ of $`𝔡`$ constructed in Lemma 13.26. Then by comparing second tensor factors, we see that $`S_{r+1}_{i=1}^rS_i=0`$ for each $`r=1,\mathrm{},N2`$. Therefore the sum of all $`S_i`$ is a free $`H`$-submodule $`F`$ of $`S(𝔡,\chi )`$ of rank $`N1`$. Since the rank of $`S(𝔡,\chi )`$ is also $`N1`$, we see that $`S(𝔡,\chi )/F`$ is a torsion $`H`$-module.
Denote by $`S`$ the subalgebra of $`S(𝔡,\chi )`$ generated by $`F`$. Since $`S(𝔡,\chi )/S`$ is a torsion $`H`$-module, we conclude, by Corollary 10.17, that $`S`$ is an ideal of $`S(𝔡,\chi )`$. Hence $`S(𝔡,\chi )=S`$ by simplicity of $`S(𝔡,\chi )`$. Now by Proposition 13.24, each subalgebra $`S_i`$ embeds uniquely in $`W(𝔡)`$. Hence $`S=S(𝔡,\chi )`$ embeds uniquely in $`W(𝔡)`$.
This completes the proof of Theorem 13.25. ∎
Combining a number of previous results, we get an explicit description of all subalgebras of $`W(𝔡)`$, and of all isomorphisms between the simple Lie pseudoalgebras listed in Theorem 13.10.
###### Corollary 13.27.
A complete list of all subalgebras $`L`$ of $`W(𝔡)=H𝔡`$ is:
(a) $`L=H𝔡^{}\mathrm{Cur}_H^{}^HW(𝔡^{})`$, where $`𝔡^{}`$ is any subalgebra of $`𝔡`$ and $`H^{}=U(𝔡^{})`$;
(b) $`L=\{_ih_ia_iH𝔡^{}|_ih_i(a_i+\chi ^{}(a_i))=0\}\mathrm{Cur}_H^{}^HS(𝔡^{},\chi ^{})`$, where $`𝔡^{}`$ is any subalgebra of $`𝔡`$ and $`\chi ^{}(𝔡^{})^{}`$ is such that $`\chi ^{}([𝔡^{},𝔡^{}])=0`$;
(c) $`L=\{(h1)(r1s)|hH\}`$, where $`r𝔡𝔡`$ and $`s𝔡`$ satisfy (4.3), (4.4). In this case, $`L`$ is isomorphic to a current pseudoalgebra over $`H(𝔡^{},\chi ^{},\omega ^{})`$ or $`K(𝔡^{},\theta ^{})`$ (see Sections 8.6 and 8.7).
###### Corollary 13.28.
All nontrivial isomorphisms among the simple Lie $`H=U(𝔡)`$-pseudoalgebras listed in Theorem 13.10 are the following $`(H^{}=U(𝔡^{}))`$:
(i) $`\mathrm{Cur}𝔤^{}\mathrm{Cur}𝔤^{\prime \prime }`$ when $`𝔤^{}`$ and $`𝔤^{\prime \prime }`$ are isomorphic Lie algebras.
(ii) $`\mathrm{Cur}_H^{}^HH(𝔡^{},\chi ^{},\omega ^{})\mathrm{Cur}_H^{}^HH(𝔡^{},\chi ^{},\omega ^{\prime \prime })`$ when $`\omega ^{\prime \prime }=c\omega ^{}`$ for some nonzero $`c𝐤`$.
(iii) $`\mathrm{Cur}_H^{}^HK(𝔡^{},\theta ^{})\mathrm{Cur}_H^{}^HK(𝔡^{},\theta ^{\prime \prime })`$ when $`\theta ^{\prime \prime }=c\theta ^{}`$ for some nonzero $`c𝐤`$.
(iv) $`\mathrm{Cur}_H^{}^HW(𝔡^{})\mathrm{Cur}_H^{}^HK(𝔡^{},\theta ^{})`$ when $`dim𝔡^{}=1`$.
(v) $`\mathrm{Cur}_H^{}^HH(𝔡^{},\chi ^{},\omega )\mathrm{Cur}_H^{}^HS(𝔡^{},\chi ^{\prime \prime })`$ when $`dim𝔡^{}=2`$ and $`\chi ^{\prime \prime }=\chi ^{}+\mathrm{tr}\mathrm{ad}`$.
### 13.7. Finite simple and semisimple Lie $`(U(𝔡)\mathrm{}𝐤[\mathrm{\Gamma }])`$-pseudoalgebras
Let, as before, $`H=U(𝔡)`$ be the universal enveloping algebra of a finite-dimensional Lie algebra $`𝔡`$. Let $`\mathrm{\Gamma }`$ be a (not necessarily finite) group acting on $`𝔡`$ by automorphisms. The action of $`\mathrm{\Gamma }`$ on $`𝔡`$ can be extended to an action on $`H`$ which we denote by $`gf`$ for $`g\mathrm{\Gamma }`$, $`fH`$. Recall that the smash product $`\stackrel{~}{H}=H\mathrm{}𝐤[\mathrm{\Gamma }]`$ is a Hopf algebra, with the product determined by $`gf=gfg^1`$, and coproduct $`\mathrm{\Delta }(fg)=\mathrm{\Delta }(f)\mathrm{\Delta }(g)`$ ($`g\mathrm{\Gamma }`$, $`fH`$).
A left $`\stackrel{~}{H}`$-module $`L`$ is the same as an $`H`$-module together with an action of $`\mathrm{\Gamma }`$ on it which is compatible with that of $`H`$. An $`\stackrel{~}{H}`$-module $`L`$ will be called finite if it is finite as an $`H`$-module.
Let $`\stackrel{~}{L}`$ be a Lie $`\stackrel{~}{H}`$-pseudoalgebra with a pseudobracket denoted as $`[a\stackrel{~}{}b]`$. By Corollary 5.3, $`\stackrel{~}{L}`$ is also a Lie $`H`$-pseudoalgebra, which we denote as $`L`$ with a pseudobracket $`[ab]`$. $`L`$ is equipped with an action of $`\mathrm{\Gamma }`$, and $`[ab]`$ is $`\mathrm{\Gamma }`$-equivariant, see (5.5). As an $`\stackrel{~}{H}`$-module, $`L=\stackrel{~}{L}`$. The relationship between the two pseudobrackets is given by (5.7).
Then the following statements are easy to check.
###### Lemma 13.29.
(i) $`[a\stackrel{~}{}b]=0`$ iff $`[gab]=0`$ for all $`g\mathrm{\Gamma }`$.
(ii) $`IL=\stackrel{~}{L}`$ is an ideal of the Lie $`\stackrel{~}{H}`$-pseudoalgebra $`\stackrel{~}{L}`$ iff it is a $`\mathrm{\Gamma }`$-invariant ideal of the Lie $`H`$-pseudoalgebra $`L`$.
(iii) If $`I`$ is as in (ii), then its derived pseudoalgebra $`[I,I]`$ is the same with respect to both pseudobrackets $`[a\stackrel{~}{}b]`$ and $`[ab]`$.
(iv) $`\mathrm{Rad}\stackrel{~}{L}=\mathrm{Rad}L`$.
###### Proof.
(i) If $`[a\stackrel{~}{}b]=0`$ then all its coefficients in front of $`(gH𝐤)_{\stackrel{~}{H}}L`$ are zero for different $`g\mathrm{\Gamma }`$. Since $`[ab](HH)_HL`$, it follows that all $`[gab]=0`$.
(ii) and (iii) are clear by (5.7).
(iv) follows from (i)–(iii) and the fact that $`\mathrm{Rad}L`$ is $`\mathrm{\Gamma }`$-invariant. ($`\mathrm{Rad}L`$ is $`\mathrm{\Gamma }`$-invariant because $`[ab]`$ is $`\mathrm{\Gamma }`$-equivariant, see (5.5).) ∎
###### Proposition 13.30.
The Lie $`\stackrel{~}{H}`$-pseudoalgebra $`\stackrel{~}{L}`$ is solvable (respectively semisimple) if and only if the Lie $`H`$-pseudoalgebra $`L`$ is.
###### Proof.
Follows from Lemmas 13.29(iv) and 13.1(iii). ∎
###### Proposition 13.31.
The Lie $`\stackrel{~}{H}`$-pseudoalgebra $`\stackrel{~}{L}`$ is finite and simple if and only if the Lie $`H`$-pseudoalgebra $`L`$ is a finite direct sum of isomorphic finite simple Lie $`H`$-pseudoalgebras and $`\mathrm{\Gamma }`$ acts on them transitively.
###### Proof.
By Lemma 13.29, $`\stackrel{~}{L}`$ is simple iff $`L`$ is not abelian and has no nontrivial $`\mathrm{\Gamma }`$-invariant ideals. In particular, $`L`$ is semisimple. Using Theorem 13.15 and the fact that $`𝐤[\mathrm{\Gamma }]I`$ is an ideal of $`L`$ if $`I`$ is an ideal, we see that $`L`$ is a direct sum of isomorphic finite simple Lie $`H`$-pseudoalgebras. ∎
### 13.8. Examples of infinite simple subalgebras of $`\mathrm{gc}_n`$
In this subsection, $`H`$ is an arbitrary cocommutative Hopf algebra. Let us define a map $`\omega :HHHH`$ by the formula:
(13.9)
$$\omega (fa)=fa_{(1)}a_{(2)}=(f1)\mathrm{\Delta }(S(a)).$$
It is easy to check that $`\omega ^2=\mathrm{id}`$; this also follows from the identities $`\omega =^1(\mathrm{id}S)=(\mathrm{id}S)`$ where $``$ is the Fourier transform defined by (2.33).
###### Lemma 13.32.
The above $`\omega `$ is an anti-involution of $`\mathrm{Cend}_1=HH`$, i.e., it is an $`H`$-linear map satisfying $`\omega ^2=\mathrm{id}`$ and
(13.10)
$$\omega (a)\omega (b)=(\sigma _H\omega )(ba),a,b\mathrm{Cend}_1,$$
where, as before, $`\sigma :HHHH`$ is the permutation of the factors.
###### Proof.
It only remains to check (13.10), which is straightforward. ∎
When $`H=U(𝔡)`$, the annihilation algebra $`𝒜(\mathrm{Cend}_1)`$ is isomorphic to the associative algebra of all differential operators on $`X`$, and $`\omega `$ induces its standard anti-involution $``$ (formal adjoint).
Let $`\gamma :\mathrm{End}(𝐤^n)\mathrm{End}(𝐤^n)`$ be an anti-involution, i.e., $`\gamma ^2=\mathrm{id}`$ and $`\gamma (A)\gamma (B)=\gamma (BA)`$. Then we can define an anti-involution $`\omega `$ of $`\mathrm{Cend}_n=HH\mathrm{End}(𝐤^n)`$ by the formula (cf. (13.9)):
(13.11)
$$\omega (faA)=fa_{(1)}a_{(2)}\gamma (A).$$
Let $`\mathrm{gc}_n(\omega )`$ be the set of all $`a\mathrm{Cend}_n`$ such that $`\omega (a)=a`$. This is a subalgebra of the Lie pseudoalgebra $`\mathrm{gc}_n`$. Indeed, it is an $`H`$-submodule because $`\omega `$ is $`H`$-linear. If $`\omega (a)=a`$, $`\omega (b)=b`$, then:
$`(\mathrm{id}_H\omega )[ab]`$ $`=(\mathrm{id}_H\omega )\left(ab(\sigma _H\mathrm{id})(ba)\right)`$
$`=(\sigma _H\mathrm{id})\omega (b)\omega (a)\omega (a)\omega (b)=[ab].`$
Two important examples of Lie pseudoalgebras $`\mathrm{gc}_n(\omega )`$ are obtained when $`𝐤^n`$ is endowed with a symmetric or skew-symmetric nondegenerate bilinear form, and $`\gamma (A)`$ is the adjoint of $`A`$ with respect to this form. In these cases, we denote $`\mathrm{gc}_n(\omega )`$ by $`\mathrm{oc}_n`$ and $`\mathrm{spc}_n`$, respectively.
###### Proposition 13.33.
Let $`H=U(𝔡)`$, $`𝔡0`$. Then $`\mathrm{oc}_n`$ and $`\mathrm{spc}_n`$ are infinite subalgebras of $`\mathrm{gc}_n`$ that act irreducibly on $`H𝐤^n`$. We have: $`\mathrm{oc}_n\mathrm{Cur}𝔤𝔩_n=\mathrm{Cur}𝔬_n`$ and $`\mathrm{spc}_n\mathrm{Cur}𝔤𝔩_n=\mathrm{Cur}𝔰𝔭_n`$.
###### Proof.
The second statement is obvious by the definitions. Since $`𝔬_n`$ ($`n3`$) and $`𝔰𝔭_n`$ ($`n2`$) act irreducibly on $`𝐤^n`$, we only have to check that the action of $`\mathrm{oc}_n`$ on $`H𝐤^n`$ is irreducible for $`n=1,2`$. Using diagonal matrices, we see that it suffices to check that $`\mathrm{oc}_1`$ acts irreducibly on $`H`$.
Recall that this action is given by (see (10.12)):
$$\alpha h=(1h)\alpha _H1\text{for}\alpha \mathrm{gc}_1=HH,hH.$$
For $`a𝔡`$, let $`\alpha =1a\omega (1a)=2a+a1\mathrm{oc}_1`$. Then $`\alpha h=(1ha)_H1+(1h)_Ha`$. If $`MH`$ is an $`\mathrm{oc}_1`$-submodule, and $`hM`$, $`h0`$, then the previous formula implies $`1M`$. Therefore $`M=H`$. ∎
###### Remark 13.34.
It follows from Theorem 14.5 below that in the case $`H=U(𝔡)`$, $`𝔡0`$, the Lie pseudoalgebras $`\mathrm{gc}_n`$, $`\mathrm{oc}_n`$ and $`\mathrm{spc}_n`$ are semisimple. In fact, one can show that in this case they are simple. In the case $`H=𝐤[\mathrm{\Gamma }]`$ with a finite group $`\mathrm{\Gamma }`$, the Lie pseudoalgebra $`\mathrm{gc}_n`$ has a center that is a free $`H`$-module of rank $`1`$, the quotient by which is simple.
If $`I`$ is a left or right ideal of the associative pseudoalgebra $`\mathrm{Cend}_n`$ and $`L`$ is a subalgebra of the Lie pseudoalgebra $`\mathrm{gc}_n`$, then their intersection $`IL`$ is again a subalgebra of $`\mathrm{gc}_n`$. All ideals of $`\mathrm{Cend}_n`$ are described in the next proposition.
###### Proposition 13.35.
(i) Any left ideal of the associative pseudoalgebra $`\mathrm{Cend}_n`$ is a sum of ideals of the form $`HRE`$ where $`RH`$ is a right ideal and $`E\mathrm{End}(𝐤^n)`$ is a left ideal.
(ii) Any right ideal of $`\mathrm{Cend}_n`$ is of the form $`\omega (I)`$ for a unique left ideal $`I`$.
(iii) $`\mathrm{Cend}_n`$ has no two-sided ideals, i.e., it is a simple associative pseudoalgebra.
###### Proof.
Let $`I\mathrm{Cend}_n`$ be a left ideal, $`\alpha =1aA\mathrm{Cend}_n`$, and $`\beta =_ig_ib_iB_iI`$ with linearly independent $`g_i`$. Then:
$$\alpha \beta =\underset{i}{}(1g_ia_{(1)})_H(1b_ia_{(2)}AB_i)(HH)_HI.$$
Taking $`a=1`$, we see that all $`1b_iAB_iI`$. In particular, $`1b_iB_iI`$, and hence each element from $`I`$ is an $`H`$-linear combination of elements of the form $`\beta =1bB`$. For such $`\beta `$, we have $`\alpha \beta =(1a_{(1)})_H(1ba_{(2)}AB)`$. For $`a𝔡`$, $`A=\mathrm{Id}`$, we get that $`1baBI`$. This proves the first part of the proposition.
Part (ii) is obvious, and part (iii) follows easily from (i) and (ii). ∎
## 14. Representation Theory of Lie Pseudoalgebras
### 14.1. Conformal version of the Lie Lemma
Let $`L`$ be a Lie $`H`$-pseudoalgebra and $`V`$ be an $`L`$-module. In this subsection, $`H=U(𝔡)`$ will be the universal enveloping algebra of a finite-dimensional Lie algebra $`𝔡`$. In particular, $`H`$ is a Noetherian ring with no divisors of zero.
Let $`IL`$ be an ideal and $`\phi \mathrm{Hom}_H(I,H)`$ be such that
(14.1)
$$V_\phi :=\{vV|av=(\phi (a)1)_HvaI\}$$
is nonzero. We will call the elements of $`V_\phi `$ eigenvectors for $`I`$ with an eigenvalue $`\phi `$. Note that every $`vV_\phi `$ is an eigenvector for the action of $`X_HI𝒜(L)`$ on $`V`$. By abuse of notation, we will also write $`a_xv=\phi (a_x)v`$ for $`aI`$, $`xX`$, $`vV_\phi `$, where $`\phi (a_x)=S(x),\phi (a)`$, cf. (9.6).
Clearly, if $`\phi =0`$, then $`V_0`$ is an $`L`$-submodule of $`V`$.
###### Lemma 14.1.
If $`\phi 0`$, then $`HV_\phi `$ is a free $`H`$-module, isomorphic to $`HV_\phi `$ with $`H`$ acting on the first tensor factor.
###### Proof.
Assume that
(14.2)
$$_if_iv_i=0$$
for some $`f_iH`$, $`v_iV_\phi `$. Let (14.2) be a relation of this form with $`f_i\mathrm{F}^{n_i}H`$ so that $`_in_i`$ is minimal. We call $`_in_i`$ the degree of the relation (14.2). Assume that $`v_i`$’s are linearly independent, so that the degree of (14.2) is positive.
We can find $`aI`$, $`xX`$ such that $`\phi (a_x)0`$. Applying $`a_x`$ to (14.2) and using (9.16), we obtain
$$_if_{i}^{}{}_{(2)}{}^{}\phi (a_{f_{i}^{}{}_{(1)}{}^{}x})v_i=0.$$
Subtracting this from (14.2), we get a relation of lower degree than (14.2), because $`\mathrm{\Delta }(f)1f+_{j=1}^n\mathrm{F}^jH\mathrm{F}^{nj}H`$ for $`f\mathrm{F}^nH`$. ∎
The following result is an analogue of Lie’s Lemma.
###### Proposition 14.2.
If $`V`$ is finite as an $`H`$-module, then:
(14.3)
$$LV_\phi (H𝐤)_H(𝔡V_\phi +V_\phi ).$$
In other words, for every $`\beta 𝒜(L)`$, there exist $`_\beta 𝔡`$ and $`A_\beta \mathrm{End}V_\phi `$ such that
(14.4)
$$\beta v=(_\beta +A_\beta )v\text{for any}vV_\phi .$$
In particular, $`HV_\phi `$ is an $`L`$-submodule of $`V`$.
###### Proof.
Fix nonzero elements $`wV_\phi `$, $`\beta 𝒜(L)`$, and let $`w_n=\beta ^nw`$. Let $`W_n`$ be the linear span of $`w_0,\mathrm{},w_n`$; we set $`W_n=0`$ for $`n<0`$. For $`aI`$, $`xX`$, we have: $`a_xw=\phi (a_x)w`$, and by induction,
(14.5)
$$a_xw_n\phi (a_x)w_n+n\phi ([a_x,\beta ])w_{n1}+W_{n2}.$$
In particular, all $`HW_n`$ are $`I`$-modules.
Since $`V`$ is a Noetherian $`H`$-module, there exists $`N0`$ such that $`HW_{N1}HW_N=HW_{N+1}`$. In particular,
(14.6)
$$w_{N+1}(N+1)hw_N+HW_{N1}$$
for some $`hH`$.
Writing (14.5) for $`n=N+1`$ and using (14.6), we get
(14.7)
$$a_xw_{N+1}\phi (a_x)(N+1)hw_N+(N+1)\phi ([a_x,\beta ])w_N+HW_{N1}.$$
On the other hand, applying $`a_x`$ to both sides of (14.6) and using the $`H`$-sesqui-linearity gives
(14.8)
$$a_xw_{N+1}\phi (a_{h_{(1)}x})(N+1)h_{(2)}w_N+HW_{N1}.$$
Comparison of the last two equations gives
(14.9)
$$fw_NHW_{N1}\text{for}f=\phi (a_x)h+\phi ([a_x,\beta ])\phi (a_{h_{(1)}x})h_{(2)}.$$
If $`f0`$, then the module $`HW_N/HW_{N1}`$ is torsion, hence $`I`$ acts on it as zero by Corollary 10.17. This gives $`a_xw_NHW_{N1}`$ for all $`aI`$, $`xX`$. Then (14.5) implies $`\phi (a_x)w_NHW_{N1}`$. Since $`HW_{N1}HW_N`$, it follows that $`\phi =0`$, which contradicts to the assumption $`f0`$.
Therefore $`f=0`$. This is possible only when $`h\mathrm{F}^1H=𝔡+𝐤`$. Then for any $`vV_\phi `$, one has:
$$0=fv=ha_xv+[a_x,\beta ]vh_{(2)}a_{h_{(1)}x}v=[a_x,\beta h]v.$$
This implies that $`(\beta h)vV_\phi `$, proving (14.4). ∎
### 14.2. Conformal version of the Lie Theorem
###### Theorem 14.3.
Let $`H=U(𝔡)\mathrm{}𝐤[\mathrm{\Gamma }]`$ with $`dim𝔡<\mathrm{}`$. Let $`L`$ be a solvable Lie $`H`$-pseudoalgebra and $`V`$ be an $`L`$-module which is finite over $`U(𝔡)`$. Then there exists an eigenvector for the action of $`L`$ on $`V`$, i.e., $`vV\{0\}`$ and $`\phi \mathrm{Hom}_H(L,H)`$ such that $`av=(\phi (a)1)_Hv`$ for all $`aL`$.
###### Proof.
Using Corollary 5.3 and Proposition 13.30, we can assume that $`H=U(𝔡)`$. The proof will be by induction on the length of the derived series of $`L`$.
First consider the case when $`L`$ is abelian. By a Zorn’s Lemma argument, it is enough to find an eigenvector when $`L=Ha`$ is abelian generated by one element $`a`$. We may assume that $`\mathrm{ker}a=0`$; then by Lemma 10.15 all $`\mathrm{ker}_na`$ are finite dimensional. Let $`n`$ be such that $`\mathrm{ker}_na0`$. Then the statement follows from the usual Lie Theorem applied to the $`𝒜(L)`$-module $`\mathrm{ker}_na`$.
Now let $`L`$ be nonabelian, $`I=[L,L]0`$. By the inductive assumption, $`I`$ has a space of eigenvectors $`V_\phi 0`$. If $`\phi =0`$, then $`V_0`$ is an $`L`$-submodule of $`V`$ on which $`I`$ acts as zero. The abelian $`H`$-pseudoalgebra $`L/I`$ has an eigenvector in $`V_0`$, which is also an eigenvector for $`L`$.
Now assume that $`\phi 0`$. By Proposition 14.2, we have for $`\alpha ,\beta 𝒜(L)`$, $`vV_\phi `$:
$`\alpha v`$ $`=(_\alpha +A_\alpha )v,`$
$`\beta v`$ $`=(_\beta +A_\beta )v,`$
$`[\alpha ,\beta ]v`$ $`=\phi ([\alpha ,\beta ])v.`$
On the other hand, we can compute:
$`\alpha \beta v`$ $`=\alpha (_\beta +A_\beta )v=_\beta (\alpha v)(_\beta \alpha )v+\alpha (A_\beta v)`$
$`=_\beta (_\alpha +A_\alpha )v(_{_\beta \alpha }+A_{_\beta \alpha })v+(_\alpha +A_\alpha )A_\beta v`$
$`=_\beta _\alpha v_{_\beta \alpha }v+_\beta A_\alpha v+_\alpha A_\beta vA_{_\beta \alpha }v+A_\alpha A_\beta v.`$
It follows that
$$[_\alpha ,_\beta ]=_{_\alpha \beta }_{_\beta \alpha }.$$
Assume that $`_{a_x}0`$ for some $`aL`$, $`xX`$, and write $`_x=_{a_x}`$ for short. For $`\alpha =a_x`$, $`\beta =a_y`$, the above equation becomes:
$$[_x,_y]=_{_xy}_{_yx}$$
(recall that $`ha_x=a_{hx}`$ for $`hH`$). Note that $`_y=0`$ if $`y\mathrm{F}_nX`$ for sufficiently large $`n`$. Take the minimal such $`n`$, and let $`x\mathrm{F}_{n1}X`$ be such that $`_x0`$. By Lemma 6.10, there exists $`y\mathrm{F}_nX`$ such that $`x=_xy`$. Then $`_y=0`$ and $`_x=_{_xy}_{_yx}=[_x,_y]=0`$, which is a contradiction.
It follows that all $`_{a_x}=0`$, hence $`L`$ preserves $`V_\phi `$. By Lemma 14.1, $`dimV_\phi <0`$, and therefore $`L`$ has an eigenvector by the usual Lie Theorem for $`𝒜(L)`$. ∎
###### Corollary 14.4.
Let $`L`$ be a solvable Lie $`H`$-pseudoalgebra and $`V`$ be a finite $`L`$-module $`(`$i.e., finite over $`U(𝔡))`$. Then $`V`$ has a filtration by $`L`$-submodules $`0=V_0V_1\mathrm{}V_n=V`$ such that for any $`i`$ the $`L`$-module $`V_{i+1}/V_i`$ is generated over $`H`$ by eigenvectors of some given eigenvalue $`\phi _i\mathrm{Hom}_H(L,H)`$.
### 14.3. Conformal version of the Cartan–Jacobson Theorem
###### Theorem 14.5.
Let $`H=U(𝔡)`$ be the universal enveloping algebra of a finite-dimensional Lie algebra $`𝔡`$. Let $`L`$ be a Lie $`H`$-pseudoalgebra acting faithfully and irreducibly on the finite $`H`$-module $`V`$. Then one of the following two possibilities holds:
(i) $`L`$ is semisimple, either finite or infinite.
(ii) $`L`$ is finite, $`\mathrm{Rad}L`$ is abelian and of rank one as an $`H`$-module. In this case, there is a subspace $`\overline{V}`$ of $`V`$ such that $`VH\overline{V}`$ is a free $`H`$-module and $`L`$ can be identified with $`(A\mathrm{Cur}𝔤)(R\mathrm{id}_{\overline{V}})\mathrm{gc}V`$, where $`A`$ is a subalgebra of $`W(𝔡)`$, $`𝔤`$ is zero or a semisimple subalgebra of $`𝔰𝔩\overline{V}`$, and $`R`$ is a nonzero left ideal of $`H`$.
###### Proof.
Assume that $`L`$ is not semisimple, i.e., it has a nonzero abelian ideal $`I`$. Then, by Theorem 14.3, $`I`$ has an eigenvector in $`V`$. If $`\overline{V}=V_\phi `$ is the corresponding eigenspace in $`V`$, then, by Proposition 14.2, $`H\overline{V}`$ is an $`L`$-submodule of $`V`$. The irreducibility of $`V`$ implies that $`V=H\overline{V}`$. Now, by Lemma 14.1, $`VH\overline{V}`$ is a free $`H`$-module, since $`\phi 0`$ by the faithfulness of $`V`$.
Proposition 14.2 and the faithfulness of $`V`$ also show that $`L`$ embeds into $`W(𝔡)\mathrm{Cur}𝔤𝔩\overline{V}\mathrm{gc}V`$. In particular, $`L`$ is finite. Then $`\mathrm{Rad}L`$ exists, and we can assume that $`\overline{V}`$ is an eigenspace for $`\mathrm{Rad}L`$. For each element $`a\mathrm{Rad}L`$ and $`v\overline{V}`$ we have $`av=(\phi (a)1)_Hv`$, which means that $`\mathrm{Rad}L`$ is identified with $`R\mathrm{id}_{\overline{V}}\mathrm{Cur}𝔤𝔩\overline{V}`$ for $`R=\phi (\mathrm{Rad}L)`$. Note that $`R`$ is of rank one, because $`R0`$ and $`H`$ has no zero divisors.
Let $`L_1=L(W(𝔡)\mathrm{Cur}𝔰𝔩\overline{V})`$. Notice that $`L_1`$ is a subalgebra of $`L_2:=L/\mathrm{Rad}L`$ and $`L_1+\mathrm{Rad}L`$ is a semidirect sum, because $`W(𝔡)\mathrm{Cur}𝔤𝔩\overline{V}=(W(𝔡)\mathrm{Cur}𝔰𝔩\overline{V})(H\mathrm{id}_{\overline{V}})`$. Since $`\mathrm{Rad}L=R\mathrm{id}_{\overline{V}}`$, this also implies that $`RL`$ is contained in $`L_1+\mathrm{Rad}L`$. Then $`RL_2`$ embeds in $`L_1`$, i.e., $`R(L_2/L_1)=0`$. Hence $`L_2/L_1`$ is torsion, and by Corollary 10.17, $`L_1`$ is an ideal of $`L_2`$. But $`L_2`$ is semisimple, and by Proposition 13.20 a finite semisimple Lie pseudoalgebra does not have proper ideals of the same rank. Therefore, $`L_1=L_2`$ and $`L=L_1\mathrm{Rad}L`$ is a semidirect sum of pseudoalgebras.
Finally, to show that $`L_1`$ is of the form $`A\mathrm{Cur}𝔤`$, notice that $`L_1\mathrm{Cur}𝔰𝔩\overline{V}`$ is an ideal of $`L_1`$, since $`\mathrm{Cur}𝔰𝔩\overline{V}`$ is an ideal of $`W(𝔡)\mathrm{Cur}𝔰𝔩\overline{V}`$. This ideal is generated over $`H`$ by abelian elements, so by Propositions 13.16 and 13.20 if it is nonzero it is of the form $`\mathrm{Cur}𝔤`$ for some semisimple subalgebra $`𝔤`$ of $`𝔰𝔩\overline{V}`$. Hence $`\mathrm{Cur}𝔤L_1W(𝔡)\mathrm{Cur}𝔤`$. But any subalgebra of $`W(𝔡)\mathrm{Cur}𝔤`$ containing $`\mathrm{Cur}𝔤`$ is equal to $`A\mathrm{Cur}𝔤`$ for some subalgebra $`A`$ of $`W(𝔡)`$. This completes the proof. ∎
### 14.4. Conformal version of Engel’s Theorem
As an application of the results of Section 14.2, we can prove a conformal analogue of Engel’s Theorem.
###### Theorem 14.6.
Let $`H=U(𝔡)\mathrm{}𝐤[\mathrm{\Gamma }]`$ with $`dim𝔡<\mathrm{}`$, and let $`L`$ be a finite Lie $`H`$-pseudoalgebra $`(`$i.e., finite over $`U(𝔡))`$. Assume that the action of any element $`\alpha 𝒜(L)`$ on $`L`$ is nilpotent. Then $`L`$ is a nilpotent Lie pseudoalgebra.
###### Proof.
First of all, note that the property that any element $`\alpha `$ of $`𝒜(L)`$ acts nilpotently on $`L`$ remains valid when we replace $`L`$ by any quotient of $`L`$ by an ideal. In particular, $`L/\mathrm{Rad}L`$ will have that property. However, $`L/\mathrm{Rad}L`$ is semisimple, and from the classification of finite semisimple Lie pseudoalgebras we see that this is impossible, unless $`L/\mathrm{Rad}L=0`$.
Therefore $`L`$ is solvable. The nilpotence of all $`\alpha 𝒜(L)`$ imply that all eigenvalues for $`L`$ are zero. Now Corollary 14.4 implies that $`L`$ is a nilpotent Lie pseudoalgebra. ∎
### 14.5. Generalized weight decomposition for nilpotent Lie pseudoalgebras
Let $`L`$ be a (not necessarily finite) Lie $`H`$-pseudoalgebra, and $`V`$ be a finite $`L`$-module, where $`H=U(𝔡)`$ for a finite-dimensional Lie algebra $`𝔡`$.
Recall that for any $`\phi \mathrm{Hom}_H(L,H)`$, the eigenspace $`V_\phi `$ of $`V`$ is defined by:
(14.10)
$$V_\phi =\{vV|av=(\phi (a)1)_HvaL\}.$$
Let $`V_1^\phi =0`$ and set inductively
(14.11)
$$V_{i+1}^\phi =H\{vV|av(\phi (a)1)_Hv(HH)_HV_i^\phi aL\}.$$
Then $`V_0^\phi =HV_\phi `$ and $`V_{i+1}^\phi /V_i^\phi =H(V/V_i^\phi )_\phi `$. The $`V_i^\phi `$ form an increasing sequence of $`H`$-submodules of $`V`$ which stabilizes (because of noetherianity) to some $`H`$-submodule of $`V`$ denoted $`V^\phi `$. If $`V_{n1}^\phi V_n^\phi =V^\phi `$, then we set the depth of $`V^\phi `$ to be $`n`$. We call $`V^\phi `$ the generalized weight submodule of $`V`$ relative to the weight $`\phi `$.
When $`L`$ is nilpotent, it is obviously solvable, and, by Corollary 14.4, any finite $`L`$-module $`V`$ has a filtration by $`L`$-submodules so that the successive quotients are generalized weight modules.
The main result of this subsection is the following theorem.
###### Theorem 14.7.
Let $`L`$ be a nilpotent Lie $`H`$-pseudoalgebra and $`V`$ be a finite $`L`$-module. Then $`V`$ decomposes as a direct sum of generalized weight modules.
###### Proof.
In order to prove the statement, it is enough to show that all $`L`$-module extensions between generalized weight modules relative to distinct weights are trivial.
The strategy is to consider first the case when $`L=T`$ is the Lie pseudoalgebra generated by one element $`T\mathrm{gc}V`$. Then in the general case, we show that the generalized weight spaces $`V^\phi `$ relative to some element $`TL`$ are $`L`$-invariant.
###### Lemma 14.8.
Let $`V`$ be a finite $`H`$-module, $`T\mathrm{gc}V`$, and $`L=T`$ be a nilpotent Lie pseudoalgebra. If $`V`$ contains a $`T`$-generalized weight module $`V^\phi `$ and $`V/V^\phi =W=W^\psi `$ with $`\psi \phi `$, then $`VV^\phi W`$ as $`L`$-modules.
###### Proof.
Since $`W_{i+1}^\psi /W_i^\psi =H(W/W_i^\psi )_\psi `$ for any $`i`$, it suffices to prove the statement when $`W=W_0^\psi =HW_\psi `$.
Let us first consider the case when $`W=H\overline{v}`$ is a cyclic $`H`$-module. In order to prove that the extension is trivial, we need to find a lifting $`vV`$ of $`\overline{v}`$ such that $`Tv=(\psi 1)_Hv`$ and to show that $`Hv+V^\phi `$ is a direct sum of $`H`$-modules (here and below, we write just $`\psi `$ instead of $`\psi (T)`$). We will prove this by induction on the depth of $`V^\phi `$, the basis of induction being trivial.
Let thus the statement be true for all $`T`$-generalized weight modules of depth $`n`$ and consider a module $`V^\phi `$ of depth $`n+1`$. Fix an arbitrary lifting $`vV`$ of $`\overline{v}`$; then:
(14.12)
$$Tv=(\psi 1)_Hvmod(HH)_HV^\phi .$$
Set
$$T_1=T,T_{m+1}=[T_mT]H^{(m+1)}_HL\text{for}m1.$$
Then we claim that for $`m>1`$, $`T_{m+1}vH^{(m+2)}_HV_n^\phi `$ implies $`T_mvH^{(m+1)}_HV_n^\phi `$. We are going to show this first in the case when $`\phi 0`$, the proof for $`\phi =0`$ only requiring minor changes.
So, let $`\phi 0`$. Then $`V^\phi /V_n^\phi =V_{n+1}^\phi /V_n^\phi =H(V^\phi /V_n^\phi )_\phi `$ is a free $`H`$-module, because it is generated by its $`\phi `$-eigenspace and we can apply Lemma 14.1. We pick some $`H`$-basis $`\{w^j\}`$ for $`V^\phi `$ modulo $`V_n^\phi `$. If $`\{h^i\}`$ is some $`𝐤`$-basis of $`H`$ compatible with its filtration, we write
(14.13)
$$T_mv=_{i,j}(\alpha _j^ih^i)_Hw^jmodH^{(m+1)}_HV_n^\phi ,$$
where $`\alpha _j^iH^m`$.
Notice that for $`m>1`$, $`T_m`$ belongs to $`H^m_H[L,L]`$ where $`[L,L]`$ is the derived algebra of $`L`$, hence all weights are zero on it. This means that
(14.14)
$$T_mV^\phi H^{(m+1)}_HV_n^\phi \text{for}m>1.$$
We have:
$$T_{m+1}v=[T_mT]v=T_m(Tv)((\sigma \mathrm{id})_H\mathrm{id})T(T_mv).$$
We compute the right-hand side, using (3.16), (3.19) and (14.12)–(14.14), and obtain:
$$T_{m+1}v=_{i,j}(\alpha _j^i\psi h_{(1)}^ih_{(2)}^i\alpha _j^i\phi h^i)_Hw^jmodH^{(m+2)}_HV_n^\phi .$$
Now the assumption $`T_{m+1}vH^{(m+2)}_HV_n^\phi `$ implies that coefficients of all $`w^j`$ must be zero. Let us choose the highest degree $`d`$ for which there is some $`h^i`$ of degree $`d`$ such that $`\alpha _j^i0`$ for some $`j`$. Then we get $`\alpha _j^i(\psi \phi )=0`$ for all $`j`$ and all $`h^i`$ of degree $`d`$, hence $`\alpha _j^i=0`$, giving a contradiction. This proves that all $`\alpha _j^i=0`$, and therefore $`T_mvH^{(m+1)}_HV_n^\phi `$.
Now, because of nilpotence of $`L`$, $`T_N=0`$ for $`N0`$, and obviously $`0vH^{(N+1)}_HV_n^\phi `$. Thus we can pull the statement back to $`m=2`$ to obtain that $`[TT]`$ maps any lifting $`v`$ of $`\overline{v}`$ inside $`H^3_HV_n^\phi `$.
Now we can choose the lifting $`v`$ of $`\overline{v}`$ so that $`Tv(\psi 1)_HvH^2_HV_n^\phi `$. Indeed, performing the same computation as above, using instead of (14.12)
$$Tv=(\psi 1)_Hv+\underset{i,j}{}(\alpha _j^ih^i)_Hw^jmodH^2_HV_n^\phi $$
for some $`\alpha _j^iH`$, we get $`\alpha _j^i(\phi \psi )(\phi \psi )\alpha _j^i=0`$. This shows that $`\alpha _j^i=c_j^i(\phi \psi )`$ for some choice of $`c_j^i𝐤`$. Now choose $`v`$ to be the lifting of $`\overline{v}`$ minimizing the top degree $`d`$ of $`h^i`$ such that some $`\alpha _j^i`$ is nonzero. Then if we replace $`v`$ by $`v^{}=v+c_j^ih^iw^j`$, all coefficients $`\alpha _j^i`$ in degree $`d`$ vanish, against minimality of $`v`$. This contradiction shows that the lifting $`v`$ can be chosen in such a way that $`\alpha _j^i=0`$ for all $`i,j`$, and $`Tv=(\psi 1)_Hv`$ modulo $`H^2_HV_n^\phi `$.
This shows that $`Hv+V_n^\phi `$ is indeed a submodule of $`V`$, and it satisfies the hypotheses of our claim. Moreover, $`V_n^\phi `$ is of depth $`n`$ and we can apply the inductive assumption to show that $`Hv+V_n^\phi `$ decomposes as a direct sum of $`L`$-submodules. This means that we can find a lifting $`\stackrel{~}{v}`$ of $`v+V_n^\phi `$ for which $`T\stackrel{~}{v}=(\psi 1)_H\stackrel{~}{v}`$ holds exactly.
We have found a lifting $`\stackrel{~}{v}`$ of $`\overline{v}`$ proving that $`V=H\stackrel{~}{v}+V^\phi `$. We are left with showing that this is a direct sum of $`H`$-modules. This is clear if $`\psi 0`$ since in this case $`H\stackrel{~}{v}`$ is free, hence projective. If instead $`\psi =0`$, assume the sum not to be free. This means that some multiple $`h\stackrel{~}{v}`$ of $`\stackrel{~}{v}`$ lies in $`V^\phi `$. Since $`\stackrel{~}{v}`$ is killed by $`T`$, so is $`h\stackrel{~}{v}`$, showing $`h\stackrel{~}{v}=0`$ as no other vector in a generalized weight module of nonzero weight $`\phi `$ is killed by $`T`$. This concludes the proof in case $`\phi 0`$.
If $`\phi =0`$, then we choose a $`𝐤`$-basis of $`V^\phi `$ modulo $`V_n^\phi `$, and use in (14.13) coefficients of the form $`\alpha _j1`$. The rest of the proof is the same.
Finally, consider the general case of a non-cyclic $`H`$-module $`W`$. We distinguish two cases. If $`\psi 0`$, then $`W=HW_\psi `$ is free by Lemma 14.1, and it decomposes as a direct sum of cyclic modules to which we can apply the above argument independently. If $`\psi =0`$, then we choose generators $`\overline{v}^i`$ of $`W`$ over $`H`$, lift them to elements $`v^i`$ of $`V`$ in such a way that each of them is mapped by $`T`$ to zero, and then argue that if $`h_i\overline{v}^i=0`$ then $`h_iv^i`$ is an element of $`V^\phi `$ killed by $`T`$, hence is zero. Therefore the extension of $`H`$-modules splits, and so does that of $`L`$-modules, by the above computation. ∎
Now let $`L`$ be any nilpotent Lie $`H`$-pseudoalgebra, $`V`$ be a finite $`L`$-module, and $`TL`$, $`T0`$.
###### Lemma 14.9.
Every $`T`$-generalized weight submodule of $`V`$ is stabilized by the action of $`L`$.
###### Proof.
We set
(14.15)
$$L_{(1)}=0,L_{(i+1)}=\{aL|[Ta](HH)_HL_{(i)}\}$$
and
(14.16)
$$V_{(1)}=0,V_{(i+1)}=\{vV|Tv(\phi (T)1)_Hv(HH)_HV_{(i)}\}.$$
Then the $`L_{(i)}`$ are $`H`$-submodules of $`L`$ whose union is $`L`$ (since $`L`$ is nilpotent), and the $`V_{(i)}`$ are vector subspaces of $`V`$ whose $`H`$-span is the $`T`$-generalized weight space $`V^\phi `$ (because $`V_i^\phi =HV_{(i)}`$ for all $`i`$).
It is easy to show by induction on $`n=i+j`$ that
(14.17)
$$L_{(i)}V_{(j)}(HH)_HV_{(i+j)}.$$
Indeed, the basis of induction (say $`n=1`$) is trivial, and the inductive step follows from (14.15), (14.16) and the identity $`[Ta]v=T(av)\left((\sigma \mathrm{id})_H\mathrm{id}\right)a(Tv)`$. Equation (14.17) implies that $`LV^\phi (HH)_HV^\phi `$, as desired. ∎
We are now able to complete the proof of Theorem 14.7. Let $`V=V_i`$ be finest among all decompositions into direct sum of $`L`$-submodules of $`V`$ such that all of the $`H`$-torsion of $`V`$ is contained in one of the $`V_i`$. Note that such a finest decomposition always exists, because any decomposition defines a partition of $`\mathrm{rank}V`$ into non-negative integers, and finer decompositions define finer partitions.
We claim that each $`V_i`$ is a generalized weight module for $`L`$. Otherwise, there must be some element $`TL`$ for which some of the $`V_i`$ is not a $`T`$-generalized weight module. But if so, then $`V_i`$ decomposes into a direct sum of its $`T`$-generalized weight submodules, and all torsion elements lie in the $`T`$-eigenspace of eigenvalue $`0`$. Since all $`T`$-generalized weight submodules are $`L`$-invariant, we obtain a contradiction. Therefore $`V`$ is a direct sum of its generalized weight submodules. ∎
### 14.6. Representations of a Lie pseudoalgebra and of its annihilation algebra
Let $`H=U(𝔡)`$ be the universal enveloping algebra of a finite-dimensional Lie algebra $`𝔡`$, and $`L`$ be a finite Lie $`H`$-pseudoalgebra.
Recall that the annihilation algebra $`=𝒜(L)`$ of $`L`$ possesses a filtration by subspaces $`=_1_0\mathrm{}`$ satisfying (7.14):
$$[_i,_j]_{i+js}\text{for all }i,j\text{ and some fixed }s,$$
that make $``$ a linearly compact Lie algebra (Proposition 7.12). Moreover, $``$ is an $`H`$-differential algebra, i.e., $`𝔡`$ acts on it by derivations. The semidirect sum $`^e:=𝔡`$ is called the extended annihilation algebra. Letting $`_n^e=_n`$ for all $`n`$ makes $`^e`$ a topological Lie algebra as well.
An $`^e`$-module (or $``$-module) $`V`$ is called conformal if any $`vV`$ is killed by some $`_n`$; in other words, if $`V`$ is a topological $`^e`$-module when endowed with the discrete topology. Now Proposition 9.4 can be reformulated as follows.
###### Proposition 14.10.
Any module $`V`$ over the Lie pseudoalgebra $`L`$ has a natural structure of a conformal $`^e`$-module, and vice versa. Moreover, $`V`$ is irreducible as an $`L`$-module iff it is irreducible as an $`^e`$-module.
Together with the next two lemmas, this proposition is an important tool in the study of representation theory of Lie pseudoalgebras.
###### Lemma 14.11.
Let $`L`$ be a finite Lie pseudoalgebra and $`V`$ be a finite $`L`$-module. For $`n1`$, let
$$\mathrm{ker}_nV=\{vV|_nv=0\},$$
so that, for example, $`\mathrm{ker}_1V=\mathrm{ker}V`$ and $`V=_n\mathrm{ker}_nV`$. Then all vector spaces $`\mathrm{ker}_nV/\mathrm{ker}V`$ are finite dimensional.
###### Proof.
The proof is an application of Lemma 10.15, using the following fact. Let $`A`$ be a vector space and $`A_iB_i`$ ($`i=1,\mathrm{},k`$) be subspaces of $`A`$ such that all $`A_i/B_i`$ are finite dimensional, then $`A_i/B_i`$ is finite dimensional. It is enough to show this for $`k=2`$, in which case it follows from the isomorphism
$$\frac{(A_1A_2)/(B_1B_2)}{(A_1B_2)/(B_1B_2)}\frac{A_1A_2}{A_1B_2}.$$
Note that $`[_s,_n]_n`$ for any $`n`$, and in particular $`_s`$ is a Lie algebra.
###### Lemma 14.12.
Let $`L`$ be a finite Lie pseudoalgebra and $`V`$ be a finite $`L`$-module such that $`\mathrm{ker}V=0`$. Then $`V`$ is locally finite as an $`_s`$-module, i.e., any vector $`vV`$ is contained in a finite-dimensional subspace invariant under $`_s`$.
###### Proof.
Any $`vV`$ is contained in some $`\mathrm{ker}_nV`$, which is finite dimensional by Lemma 14.11, and $`_s`$-invariant because $`[_s,_n]_n`$. ∎
Let $`V`$ be a finite irreducible $`L`$-module. Then $`\mathrm{ker}V=0`$. Take some $`n`$ such that $`\mathrm{ker}_nV0`$. This space is finite dimensional and $`_s`$-invariant; let $`U`$ be an irreducible $`_s`$-submodule of $`\mathrm{ker}_nV`$. The $`^e`$-submodule of $`V`$ generated by $`U`$ is a factor of the induced module $`\mathrm{Ind}__s^^eU`$. Thefore, $`V`$ is a factor module of $`\mathrm{Ind}__s^^eU`$.
In many cases $`𝔡`$ acts on $``$ by inner derivations so that we have an injective homomorphism $`𝔡`$. In this case, $`^e`$ is isomorphic to the direct sum of $`𝔡`$ and $``$, and we have $`\mathrm{Ind}__s^^eUH\mathrm{Ind}__s^{}U`$.
The above results, combined with the results of Rudakov \[Ru1, Ru2\] and \[Ko\], will allow us to classify all finite irreducible representations of all finite semisimple Lie pseudoalgebras (work in progress).
## 15. Cohomology of Lie Pseudoalgebras
### 15.1. The complexes $`C^{}(L,M)`$ and $`\stackrel{~}{C}^{}(L,M)`$
Recall that in Section 3 we defined cohomology of a Lie algebra in any pseudotensor category (Definition 3.6). Now we will spell out this definition for the case of Lie $`H`$-pseudoalgebras, i.e., for the pseudotensor category $`^{}(H)`$ (see (3.4)). As before, $`H`$ is a cocommutative Hopf algebra. Let $`L`$ be a Lie $`H`$-pseudoalgebra and $`M`$ be an $`L`$-module.
By definition, $`C^n(L,M)`$, $`n1`$, consists of all
(15.1)
$$\gamma \mathrm{Lin}(\{\underset{n}{\underset{}{L,\mathrm{},L}}\},M):=\mathrm{Hom}_{H^n}(L^n,H^n_HM)$$
that are skew-symmetric (see Figure 4). Explicitly, $`\gamma `$ has the following defining properties (cf. (3.23), (3.24)):
(15.2)
$$\gamma (h_1a_1\mathrm{}h_na_n)=((h_1\mathrm{}h_n)_H1)\gamma (a_1\mathrm{}a_n)$$
for $`h_iH`$, $`a_iL`$.
(15.3)
$$\begin{array}{cc}\hfill \gamma (a_1\mathrm{}& a_{i+1}a_i\mathrm{}a_n)\hfill \\ & =(\sigma _{i,i+1}_H\mathrm{id})\gamma (a_1\mathrm{}a_ia_{i+1}\mathrm{}a_n),\hfill \end{array}$$
where $`\sigma _{i,i+1}:H^nH^n`$ is the transposition of the $`i`$th and $`(i+1)`$st factors.
For $`n=0`$, we put $`C^0(L,M)=𝐤_HMM/H_+M`$, where $`H_+=\{hH|\epsilon (h)=0\}`$ is the augmentation ideal. The differential $`d:C^0(L,M)=𝐤_HMC^1(L,M)=\mathrm{Hom}_H(L,M)`$ is given by:
(15.4)
$$\begin{array}{cc}\hfill \left(d(1_Hm)\right)(a)=& _i(\mathrm{id}\epsilon )(h_i)m_iM\hfill \\ & \text{if}am=_ih_i_Hm_iH^2_HM\hfill \end{array}$$
for $`aL`$, $`mM`$.
For $`n1`$, the differential $`d:C^n(L,M)C^{n+1}(L,M)`$ is given by Figure 5. Explicitly:
(15.5)
$$\begin{array}{cc}\hfill (d\gamma )& (a_1\mathrm{}a_{n+1})\hfill \\ & =\underset{1in+1}{}(1)^{i+1}(\sigma _{1i}_H\mathrm{id})a_i\gamma (a_1\mathrm{}\widehat{a}_i\mathrm{}a_{n+1})\hfill \\ & +\underset{1i<jn+1}{}(1)^{i+j}(\sigma _{1i,\mathrm{\hspace{0.17em}2}j}_H\mathrm{id})\hfill \\ & \times \gamma ([a_ia_j]a_1\mathrm{}\widehat{a}_i\mathrm{}\widehat{a}_j\mathrm{}a_{n+1}),\hfill \end{array}$$
where $`\sigma _{1i}`$ is the permutation $`h_ih_1\mathrm{}h_{i1}h_{i+1}\mathrm{}h_{n+1}h_1\mathrm{}h_{n+1}`$, and $`\sigma _{1i,\mathrm{\hspace{0.17em}2}j}`$ is the permutation $`h_ih_jh_1\mathrm{}h_{i1}h_{i+1}\mathrm{}h_{j1}h_{j+1}\mathrm{}h_{n+1}h_1\mathrm{}h_{n+1}`$.
In (15.5) we also use the following conventions. If $`ab=_if_i_Hc_iH^2_HM`$ for $`aL`$, $`bM`$, then for any $`fH^n`$ we set:
$$a(f_Hb)=_i(1f)(\mathrm{id}\mathrm{\Delta }^{(n1)})(f_i)_Hc_iH^{(n+1)}_HM,$$
where $`\mathrm{\Delta }^{(n1)}=(\mathrm{id}\mathrm{}\mathrm{id}\mathrm{\Delta })\mathrm{}(\mathrm{id}\mathrm{\Delta })\mathrm{\Delta }:HH^n`$ is the iterated comultiplication ($`\mathrm{\Delta }^{(0)}:=\mathrm{id}`$). Similarly, if $`\gamma (a_1\mathrm{}a_n)=_ig_i_Hv_iH^n_HM`$, then for $`gH^2`$ we set:
$$\begin{array}{c}\gamma ((g_Ha_1)a_2\mathrm{}a_n)\hfill \\ \hfill =_i(g1^{(n1)})(\mathrm{\Delta }\mathrm{id}^{(n1)})(g_i)_Hv_iH^{(n+1)}_HM.\end{array}$$
These conventions reflect the compositions of polylinear maps in $`^{}(H)`$, see (3.8). Note that (15.5) holds also for $`n=0`$ if we define $`\mathrm{\Delta }^{(1)}:=\epsilon `$.
The fact that $`d^2=0`$ is most easily checked using Figure 5 and the same argument as in the usual Lie algebra case. The cohomology of the resulting complex $`C^{}(L,M)`$ is called the cohomology of $`L`$ with coefficients in $`M`$ and is denoted by $`\mathrm{H}^{}(L,M)`$.
One can also modify the above definition by replacing everywhere $`_H`$ with $``$. Let $`\stackrel{~}{C}^n(L,M)`$ consist of all skew-symmetric $`\gamma \mathrm{Hom}_{H^n}(L^n,H^nM)`$, cf. (15.2), (15.3). Then we can define a differential $`d:\stackrel{~}{C}^n(L,M)\stackrel{~}{C}^{n+1}(L,M)`$ by (15.5) with $`_H`$ replaced everywhere by $``$; then again $`d^2=0`$. (In fact, one can define a pseudotensor category $`\stackrel{~}{}^{}(H)`$ by replacing $`_H`$ with $``$ everywhere in the definition of $`^{}(H)`$.) The corresponding cohomology $`\stackrel{~}{\mathrm{H}}^{}(L,M)`$ will be called the basic cohomology of $`L`$ with coefficients in $`M`$. In contrast, $`\mathrm{H}^{}(L,M)`$ is sometimes called the reduced cohomology (cf. \[BKV\]).
### 15.2. Extensions and deformations
We will show that the cohomology theory of Lie pseudoalgebras defined in Section 15.1 describes extensions and deformations, just as any cohomology theory. This result is a straightforward generalization of Theorem 3.1 from \[BKV\].
###### Theorem 15.1.
(i) The isomorphism classes of $`H`$-split extensions
$$0MEN0$$
of finite modules over a Lie $`H`$-pseudoalgebra $`L`$ are in one-to-one correspondence with elements of $`\mathrm{H}^1(L,\mathrm{Chom}(N,M))`$.
(ii) Let $`C`$ be an $`L`$-module, considered as a Lie $`H`$-pseudoalgebra with respect to the zero pseudobracket. Then the equivalence classes of $`H`$-split “abelian” extensions
$$0C\widehat{L}L0$$
of the Lie $`H`$-pseudoalgebra $`L`$ correspond bijectively to $`\mathrm{H}^2(L,C)`$.
(iii) The equivalence classes of first-order deformations of a Lie $`H`$-pseudoalgebra $`L`$ $`(`$leaving the $`H`$-action intact$`)`$ correspond bijectively to $`\mathrm{H}^2(L,L)`$.
###### Proof.
(i) Let
$$0M\stackrel{i}{}E\stackrel{p}{}N0$$
be an extension of $`L`$-modules, which is split over $`H`$. Choose a splitting $`E=MN=\{m+n|mM,nN\}`$ as $`H`$-modules. The fact that $`i`$ and $`p`$ are homomorphisms of $`L`$-modules implies ($`aL`$, $`mM`$, $`nN`$):
(15.6) $`a_Em`$ $`=a_Mm,`$
(15.7) $`a_Ena_Nn`$ $`=:\gamma (a)(n)H^2_HM.`$
It is clear that $`\gamma (a)\mathrm{Chom}(N,M)`$ and $`\gamma :L\mathrm{Chom}(N,M)`$ is $`H`$-linear; in other words, $`\gamma C^1(L,\mathrm{Chom}(N,M))=\mathrm{Hom}_H(L,\mathrm{Chom}(N,M))`$.
For $`a,bL`$, $`nN`$, we have (cf. (3.26)):
$`[ab]_En`$ $`=a_E(b_En)((\sigma \mathrm{id})_H\mathrm{id})(b_E(a_En)),`$
$`[ab]_Nn`$ $`=a_N(b_Nn)((\sigma \mathrm{id})_H\mathrm{id})(b_N(a_Nn)).`$
Subtracting these two equations and using (15.6), (15.7), we get:
$`\gamma ([ab])(n)`$ $`=a_M\gamma (b)(n)((\sigma \mathrm{id})_H\mathrm{id})\gamma (b)(a_Nn)`$
$`((\sigma \mathrm{id})_H\mathrm{id})b_M\gamma (a)(n)+\gamma (a)(b_Nn)`$
$`=\left((a\gamma )(b)\right)(n)((\sigma \mathrm{id})_H\mathrm{id})\left((b\gamma )(a)\right)(n)`$
(recall that the action of $`L`$ on $`\mathrm{Chom}(N,M)`$ was defined in Remark 10.7). The last equation means that $`d\gamma =0`$.
If we choose another splitting of $`H`$-modules $`E=M^{}N=\{m+^{}n|mM,nN\}`$, then it will differ by an element $`\phi `$ of $`\mathrm{Hom}_H(N,M)`$: $`m+n=(m+\phi (n))+^{}n`$. Then the corresponding
$$\gamma (a)(n)=a_M\phi (n)(\mathrm{id}_{HH}_H\phi )(a_Nn)+\gamma ^{}(a)(n).$$
Since $`\mathrm{Hom}_H(N,M)𝐤_H\mathrm{Chom}(N,M)=C^0(L,\mathrm{Chom}(N,M))`$ (see Remark 10.3), we get $`\gamma (a)=a\phi +\gamma ^{}(a)`$, i.e., $`\gamma =d\phi +\gamma ^{}`$.
Conversely, given an element of $`\mathrm{H}^1(L,\mathrm{Chom}(N,M))`$, we can choose a representative $`\gamma C^1(L,\mathrm{Chom}(N,M))`$ and define an action $`_E`$ of $`L`$ on $`E=MN`$ by (15.6), (15.7), which will depend only on the cohomology class of $`\gamma `$. This proves (i).
The proof of (ii) is similar. Write $`\widehat{L}=LC=\{a+c|aL,cC\}`$ as $`H`$-modules. Denoting the pseudobracket of $`\widehat{L}`$ by $`[a\widehat{}b]`$, we have for $`a,bL`$, $`c,c_1C`$:
$`[a\widehat{}c]`$ $`=ac,`$
$`[c\widehat{}c_1]`$ $`=0,`$
$`[a\widehat{}b][ab]`$ $`=:\gamma (ab)H^2_HC.`$
It is clear that $`\gamma C^2(L,C)`$, and the Jacobi identity for $`\widehat{L}`$ implies $`d\gamma =0`$.
(iii) A first-order deformation of $`L`$ is the structure of a Lie $`H`$-pseudoalgebra on $`\widehat{L}=L[ϵ]/(ϵ^2)=LLϵ`$, where $`H`$ acts trivially on $`ϵ`$, such that the map $`\widehat{L}L`$ given by putting $`ϵ=0`$ is a homomorphism of Lie pseudoalgebras. This means that
$$0Lϵ\widehat{L}L0$$
is an abelian extension of Lie pseudoalgebras, so (iii) follows from (ii). ∎
### 15.3. Relation to Gelfand–Fuchs cohomology
Let again $`L`$ be a Lie $`H`$-pseudoalgebra and $`=𝒜(L):=X_HL`$ be its annihilation Lie algebra. Recall that (by Proposition 9.4) any $`L`$-module $`M`$ has a natural structure of an $``$-module, given by $`(x_Ha)m=a_xm`$ $`(aL,xX,mM)`$, where $`a_xm`$ is the $`x`$-product defined by (cf. (9.6)):
$$a_xm=_iS(x),g_iv_i\text{if}am=_i(g_i1)_Hv_iH^2_HM.$$
Similarly, for $`\gamma \stackrel{~}{C}^n(L,M)`$ and $`x_1,\mathrm{},x_nX`$, we define
$$\gamma _{x_1,\mathrm{},x_n}(a_1\mathrm{}a_n)=_iS(x_1),g_{i,1}\mathrm{}S(x_n),g_{i,n}v_i$$
if
$$\gamma (a_1\mathrm{}a_n)=_i(g_{i,1}\mathrm{}g_{i,n})v_iH^nM.$$
The $`H`$-polylinearity (15.2) of $`\gamma `$ implies that the map $`𝒜\gamma :^nM`$, given by
$$(𝒜\gamma )\left((x_1_Ha_1)\mathrm{}(x_n_Ha_n)\right):=\gamma _{x_1,\mathrm{},x_n}(a_1\mathrm{}a_n),$$
is well defined. Moreover, $`𝒜\gamma `$ is skew-symmetric (i.e., it is map from $`^n`$ to $`M`$) because of skew-symmetry (15.3) of $`\gamma `$.
Therefore, we can consider $`𝒜\gamma `$ as an $`n`$-cochain for the Lie algebra $``$ with coefficients in $`M`$. It is not difficult to check that the map $`𝒜:\stackrel{~}{C}^n(L,M)C^n(,M)`$ commutes with the differentials (this also follows from the results of Section 7.2). The following result is proved in the same way as Proposition 9.4.
###### Proposition 15.2.
The above map $`𝒜:\stackrel{~}{C}^{}(L,M)C^{}(,M)`$ is an isomorphism from the complex $`\stackrel{~}{C}^{}(L,M)`$ to the subcomplex $`C_{\mathrm{GF}}^{}(,M)`$ of $`C^{}(,M)`$ consisting of local cochains, i.e., cochains $`𝒜\gamma `$ satisfying
(15.8)
$$(𝒜\gamma )\left((x_1_Ha_1)\mathrm{}(x_n_Ha_n)\right)=0,\text{for}x_1\mathrm{F}_kX,k0,$$
for any fixed $`x_2,\mathrm{},x_n`$ and $`a_1,\mathrm{},a_n`$.
Note that the locality condition (15.8) means that $`𝒜\gamma `$ is continuous when $`M`$ is endowed with the discrete topology and $``$ with the topology defined in Section 7.4. Therefore we have:
###### Corollary 15.3.
The basic cohomology $`\stackrel{~}{\mathrm{H}}^{}(L,M)`$ of a Lie pseudoalgebra $`L`$ is isomorphic to the Gelfand–Fuchs cohomology $`\mathrm{H}_{\mathrm{GF}}^{}(,M)`$ of its annihilation Lie algebra $``$.
Recall that $`H`$ acts on $`=X_HL`$ via its left action on $`X`$: $`h(x_Ha)=hx_Ha`$ ($`hH`$, $`xX`$, $`aL`$). Using the comultiplication $`\mathrm{\Delta }^{(n1)}(h)=h_{(1)}\mathrm{}h_{(n)}`$, we also get an action of $`H`$ on $`^n`$. It follows from (2.18), (2.25) that for $`hH`$, $`\alpha ^n`$, $`\gamma \stackrel{~}{C}^n(L,M)`$, one has:
$$(𝒜\gamma )(h\alpha )=(𝒜(\gamma h))(\alpha ),$$
where $`\gamma h\stackrel{~}{C}^n(L,M)`$ is defined by:
$`(\gamma h)(a_1\mathrm{}a_n)`$ $`=_ig_i\mathrm{\Delta }^{(n1)}(h)v_i`$
if
$`\gamma (a_1\mathrm{}a_n)`$ $`=_ig_iv_iH^nM.`$
Considering $`C^n(L,M)`$ instead of $`\stackrel{~}{C}^n(L,M)`$ amounts to replacing $``$ with $`_H`$, i.e., to factoring by the relations
$$(\gamma h)(a_1\mathrm{}a_n)(1^nh)\gamma (a_1\mathrm{}a_n),hH.$$
In terms of $`𝒜\gamma `$, this corresponds to factoring by
$$h\left((𝒜\gamma )(\alpha )\right)(𝒜\gamma )(h\alpha )=\left(h(𝒜\gamma )(𝒜\gamma )h\right)(\alpha ).$$
This implies the next result.
###### Proposition 15.4.
The isomorphism $`𝒜:\stackrel{~}{C}^{}(L,M)\stackrel{}{}C_{\mathrm{GF}}^{}(,M)`$ induces an isomorphism from $`C^{}(L,M)`$ to the quotient complex of $`C_{\mathrm{GF}}^{}(,M)`$ with respect to the subcomplex $`\{hcch|cC_{\mathrm{GF}}^{}(,M),hH\}`$.
When $`H=U(𝔡)`$, we can define an action of $`H`$ on $`C_{\mathrm{GF}}^{}C_{\mathrm{GF}}^{}(,M)`$ by $`hc:=hc+cS(h)`$. This action commutes with the differential $`d`$, and $`𝒜`$ induces an isomorphism from $`C^{}(L,M)`$ to the quotient complex $`C_{\mathrm{GF}}^{}/HC_{\mathrm{GF}}^{}`$. The Lie algebra $`𝔡`$ acts on $`C_{\mathrm{GF}}^{}`$, and clearly $`C_{\mathrm{GF}}^{}/HC_{\mathrm{GF}}^{}=C_{\mathrm{GF}}^{}/𝔡C_{\mathrm{GF}}^{}`$. We have an exact sequence of complexes
$$0𝔡C_{\mathrm{GF}}^{}C_{\mathrm{GF}}^{}C_{\mathrm{GF}}^{}/𝔡C_{\mathrm{GF}}^{}0,$$
which gives a long exact sequence for cohomology
(15.9)
$$\begin{array}{cc}\hfill \mathrm{}\mathrm{H}^i(𝔡C_{\mathrm{GF}}^{})\mathrm{H}^i(C_{\mathrm{GF}}^{})& \mathrm{H}^i(C_{\mathrm{GF}}^{}/𝔡C_{\mathrm{GF}}^{})\hfill \\ \hfill \mathrm{H}^{i+1}(𝔡C_{\mathrm{GF}}^{})& \mathrm{H}^{i+1}(C_{\mathrm{GF}}^{})\mathrm{}.\hfill \end{array}$$
###### Remark 15.5 (\[BKV\]).
If $`dim𝔡=1`$, then $`𝔡`$ acts freely on $`C_{\mathrm{GF}}^i`$ for $`i>0`$, and we have $`\mathrm{H}^i(𝔡C_{\mathrm{GF}}^{})\mathrm{H}^i(C_{\mathrm{GF}}^{})`$ for $`i>0`$. When $`M`$ is a free $`H`$-module, this is also true for $`i=0`$.
###### Proposition 15.6.
Assume that $`𝔡`$ acts on $``$ by inner derivations and that the action of $`𝔡`$ on $`M`$ coincides with that of its image in $``$. Then for any $`i0`$, we have isomorphisms
(15.10) $`\mathrm{H}^i(L,M)`$ $`\mathrm{H}_{\mathrm{GF}}^i(,M)\mathrm{H}^{i+1}(𝔡C_{\mathrm{GF}}^{}).`$
If, in addition, $`dim𝔡=1`$, then we have for $`i>0`$
(15.11) $`\mathrm{H}^i(L,M)`$ $`\mathrm{H}_{\mathrm{GF}}^i(,M)\mathrm{H}_{\mathrm{GF}}^{i+1}(,M).`$
This also holds for $`i=0`$, provided that $`M`$ is a free $`H`$-module.
###### Proof.
Since the adjoint action of $``$ on $`\mathrm{H}_{\mathrm{GF}}^{}(,M)`$ is trivial, we obtain that $`\mathrm{H}^i(𝔡C_{\mathrm{GF}}^{})`$ maps to zero in the exact sequence (15.9). Therefore we have exact sequences
$$0\mathrm{H}^i(C_{\mathrm{GF}}^{})\mathrm{H}^i(C_{\mathrm{GF}}^{}/𝔡C_{\mathrm{GF}}^{})\mathrm{H}^{i+1}(𝔡C_{\mathrm{GF}}^{})0,$$
which lead to isomorphisms (15.10). Formula (15.11) follows from Remark 15.5. ∎
Note that in general we have:
$$dim\mathrm{H}^i(L,M)dim\mathrm{H}_{\mathrm{GF}}^i(,M)+dim\mathrm{H}^{i+1}(𝔡C_{\mathrm{GF}}^{}).$$
The above results provide a tool for computing the cohomology of Lie pseudoalgebras, by making use of the known results on Gelfand–Fuchs cohomology of Lie algebras of vector fields \[Fu\].
### 15.4. Central extensions of finite simple Lie pseudoalgebras
In this section we determine by a direct computation all nontrivial central extensions of a finite simple Lie pseudoalgebra $`L`$ with trivial coefficients (see Theorem 15.13 below).
Such a central extension of $`L`$ is isomorphic as an $`H`$-module to $`\widehat{L}=L𝐤1`$, where the action of $`H`$ on $`1`$ is given by $`h1=\epsilon (h)1`$. The pseudobracket is then
(15.12)
$$[a\widehat{}b]=[ab]+\gamma (a,b)_H1,a,bL,$$
where $`\gamma (a,b)HH`$.
Notice that a tensor product $`(h^1\mathrm{}h^n)_H1H^n_H𝐤`$ can always be re-expressed as $`(h^1h_{(1)}^n\mathrm{}h^{n1}h_{((n1))}^n1)_H1`$, and this coefficient is unique in $`H^{(n1)}1`$ (see Lemma 2.5).
Therefore, the above bracket is uniquely determined by the unique $`\beta (a,b)H`$ such that
(15.13)
$$\gamma (a,b)_H1=(\beta (a,b)1)_H1;$$
we will call this map $`\beta :LLH`$ the cocycle representing the central extension. Then $`H`$-bilinearity and skew-symmetry of the pseudobracket give the following properties of this cocycle:
(15.14)
$$\beta (ha,b)=h\beta (a,b),\beta (a,hb)=\beta (a,b)S(h),\beta (a,b)=S\left(\beta (b,a)\right),$$
for all $`a,bL`$, $`hH`$.
We consider two central extensions equivalent if they are isomorphic Lie pseudoalgebras. An isomorphism of $`\widehat{L}`$ is given by an embedding $`L\widehat{L}`$ projecting to the identity on $`L`$. All such embeddings are uniquely determined by $`H`$-linear maps $`\varphi :L𝐤`$. Then the cocycles representing the two equivalent central extensions differ by $`\tau _\varphi (a,b)`$ such that
(15.15)
$$(\tau _\varphi (a,b)1)_H1=(\mathrm{id}_{HH}_H\varphi )([ab]).$$
This is called a trivial cocycle.
If $`L=He`$ is an $`H=U(𝔡)`$-module which is free on the generator $`e`$, such that
$$[ee]=\alpha _He,\alpha =r+s11s,r𝔡𝔡,s𝔡,$$
then a cocycle $`\beta (a,b)`$ is completely determined by its value $`\beta =\beta (e,e)H`$. Trivial cocycles are of the form
$$\tau =\tau (e,e)=\varphi (e)(2sx),\varphi (e)𝐤,$$
where
(15.16)
$$x=\frac{1}{2}\underset{i,j}{}r^{ij}[_i,_j]\text{if}r=\underset{i,j}{}r^{ij}_i_j.$$
###### Lemma 15.7.
Let $`L=He`$ be a Lie pseudoalgebra as above. Then $`\mathrm{H}^2(L,𝐤)B/𝐤(2sx)`$, where $`B`$ is the space of elements $`\beta H`$ satisfying the following two conditions:
(15.17) $`\beta `$ $`=S(\beta ),`$
(15.18) $`\alpha \mathrm{\Delta }(\beta )`$ $`=(\beta 1+1\beta )\alpha +\beta (3sx)(3sx)\beta .`$
Moreover, when $`r0`$, then $`\beta 𝔡`$, and (15.18) becomes equivalent to the following system of equations:
(15.19) $`[s,\beta ]`$ $`=0,`$
(15.20) $`[r,\mathrm{\Delta }(\beta )]`$ $`=\beta (3sx)(3sx)\beta .`$
###### Proof.
Let $`\widehat{L}=He+𝐤1`$ be a central extension of $`L`$ with a pseudobracket
$$[e\widehat{}e]=\alpha _He+(\beta 1)_H1,$$
where $`h1=\epsilon (h)1`$ for $`hH`$.
The skew-symmetry of $`[e\widehat{}e]`$ is equivalent to (15.17). The Jacobi identity is equivalent to Jacobi identity for $`[ee]`$ together with the following cocycle condition for $`\gamma =\beta 1`$ (cf. Proposition 4.1):
(15.21)
$$\begin{array}{cc}\hfill (\alpha 1)(\mathrm{\Delta }\mathrm{id})(\gamma )_H1& =(1\alpha )(\mathrm{id}\mathrm{\Delta })(\gamma )_H1\hfill \\ & (\sigma \mathrm{id})\left((1\alpha )(\mathrm{id}\mathrm{\Delta })(\gamma )\right)_H1.\hfill \end{array}$$
With the usual notation $`r_{12}=r1`$, $`s_1=s11`$, etc., we have:
$`(\alpha 1)(\mathrm{\Delta }\mathrm{id})(\gamma )_H1`$ $`=(\alpha \mathrm{\Delta }(\beta )1)_H1,`$
$`(1\alpha )(\mathrm{id}\mathrm{\Delta })(\gamma )_H1`$ $`=(r_{23}+s_2s_3)\beta _1_H1=\beta _1(r_{23}+s_2s_3)_H1`$
$`=\beta _1(`$ $`r_{21}x_2+s_1+2s_2)_H1=\beta _1(\alpha _{12}+3s_2x_2)_H1.`$
From here it is easy to see that (15.21) is equivalent to (15.18).
Let now $`r`$ be nonzero. Rewrite (15.18) in the form
$$\alpha \left(\mathrm{\Delta }(\beta )\beta 11\beta \right)=[\beta 1+1\beta ,\alpha ]+\beta (3sx)(3sx)\beta .$$
If $`\beta 𝔡+𝐤`$, then the degree of the left-hand side equals $`\mathrm{deg}\beta +2`$ while that of the right-hand side is at most $`\mathrm{deg}\beta +1`$, giving a contradiction. So $`\beta 𝔡+𝐤`$, and (15.17) shows that $`\beta 𝔡`$. ∎
###### Proposition 15.8.
Let $`𝔡^{}𝔡`$ be finite-dimensional Lie algebras, $`H=U(𝔡)`$, $`H^{}=U(𝔡^{})`$, and let $`L=\mathrm{Cur}_H^{}^HW(𝔡^{})`$.
(i) If $`dim𝔡^{}=1`$, then $`\mathrm{H}^2(L,𝐤)`$ is $`1`$-dimensional.
(ii) If $`𝔡`$ is abelian and $`dim𝔡^{}>1`$, then $`\mathrm{H}^2(L,𝐤)=0`$.
###### Proof.
(i) The Lie pseudoalgebra $`L=He`$ is free of rank one, with $`e=1s`$, $`s𝔡^{}\{0\}`$, hence we can use Lemma 15.7. In this case $`\alpha =s11s`$, and equation (15.18) becomes
$$(s11s)(\mathrm{\Delta }(\beta )\beta 11\beta )=3(\beta ss\beta )$$
for $`\beta H`$. We choose a basis $`\{_i\}`$ of $`𝔡`$ such that $`_1=s`$, and express $`\beta `$ in a Poincaré–Birkhoff–Witt basis as $`\beta =_I\beta _I^{(I)}`$, $`\beta _I𝐤`$ (see Example 2.3). Then the above equation becomes:
$$\underset{I,J0}{}\beta _{I+J}(_1^{(I)}^{(J)}^{(I)}_1^{(J)})=\underset{I}{}3\beta _I(^{(I)}_1_1^{(I)}).$$
Comparing terms of the form $`h_j`$ ($`j1`$) we find that $`\beta _I`$ is zero unless $`I=(i,0,\mathrm{},0)`$ for some $`i`$. Hence $`\beta =_i\beta _is^i`$, $`\beta _i𝐤`$. Substituting and comparing coefficients, we obtain that $`\beta =\beta _1s+\beta _3s^3`$. This obviously satisfies (15.17). The trivial cocycles are multiples of $`2s`$, hence $`s^3`$ is the unique central extension up to scalar multiples. This is the well-known Virasoro central extension.
(ii) Choose a basis of $`𝔡^{}`$ and let $`\beta `$ be a cocycle representing a central extension of $`LH𝔡^{}`$. Then for each basis element $`a`$, $`\beta `$ restricts to a cocycle of $`HaL`$, which is a current Lie pseudoalgebra over $`W(𝐤a)`$. By part (i) we can then add to $`\beta `$ a trivial cocycle as to make $`\beta (1a,1a)=c_aa^3`$, $`c_a𝐤`$, for every such basis element $`a𝔡^{}`$. Denoting $`\beta =\beta (1a,1b)`$, the Jacobi identity for elements $`1a,1a,1b`$ then gives:
(15.22)
$$(a11a)\mathrm{\Delta }(\beta )=c_a(a^3bba^3)+(\beta 11\beta )\mathrm{\Delta }(a).$$
Let $`a,b`$ be distinct elements in the above basis, which we extend to a basis $`\{_i\}`$ of $`𝔡`$ with $`_1=a`$, $`_2=b`$. We substitute the Poincaré–Birkhoff–Witt basis expression $`\beta =_I\beta _I^{(I)}`$ in (15.22), to get:
$`c_a(_1^3_2_2_1^3)`$ $`={\displaystyle \underset{I,J}{}}\beta _{I+J}(_1^{(I)}^{(J)}^{(I)}_1^{(J)})`$
$`{\displaystyle \underset{J}{}}\beta _J(_1^{(J)}1+^{(J)}_1_1^{(J)}1_1^{(J)}).`$
Comparing coefficients of the form $`h_j`$ for $`j1`$, we find that $`\beta _I`$ can be nonzero only when $`I=(2,1,0,\mathrm{},0)`$, in which case $`\beta _I=2c_a`$, and when $`I=(i,0,\mathrm{},0)`$ for some $`i`$. This means that $`\beta =\beta (1a,1b)=f(a)+c_aa^2b`$ for some polynomial $`f`$.
We can repeat the same argument after switching the roles of $`a`$ and $`b`$, to get: $`\beta (1b,1a)=g(b)+c_bb^2a`$. Then the skew-symmetry $`\beta (1a,1b)=S\left(\beta (1b,1a)\right)`$ implies: $`f(a)+c_aa^2b=g(b)+c_bab^2`$. This is possible only when $`f=0`$, $`c_a=0`$. Therefore $`\beta `$ is identically zero. ∎
###### Proposition 15.9.
Let $`𝔡^{}𝔡`$ be abelian finite-dimensional Lie algebras, $`H=U(𝔡),H^{}=U(𝔡^{})`$, and let $`L=\mathrm{Cur}_H^{}^HH(𝔡^{},\chi ,\omega )`$. Then $`\mathrm{H}^2(L,𝐤)`$ is isomorphic to $`𝔡`$ if $`\chi =0`$, and is trivial otherwise.
###### Proof.
$`L`$ is free of rank one and $`r0`$, hence (15.18) becomes $`3(\beta ss\beta )=0`$, $`\beta 𝔡`$. This is satisfied only by multiples of $`s`$ if $`s0`$ and by all elements of $`𝔡`$ otherwise. Since $`𝔡`$ is abelian, then $`x=0`$ and trivial cocycles are multiples of $`s`$. ∎
###### Proposition 15.10.
Let $`𝔡^{}`$ be the Heisenberg Lie algebra of dimension $`N=2n+13`$, and $`𝔡=𝔡^{}𝔡_0`$ be the direct sum of $`𝔡^{}`$ and an abelian Lie algebra $`𝔡_0`$. Let $`H=U(𝔡),H^{}=U(𝔡^{})`$, and $`L=\mathrm{Cur}_H^{}^HK(𝔡^{},\theta )`$. Then $`\mathrm{H}^2(L,𝐤)=0`$.
###### Proof.
$`L`$ is free of rank one, and
$$\alpha =\underset{i=1}{\overset{n}{}}(a_ib_ib_ia_i)c1+1c,$$
where $`\{a_i,b_i,c\}`$ is a basis of $`𝔡^{}`$ with the only nonzero commutation relations $`[a_i,b_i]=c`$, $`1in`$ (see Example 8.19).
It is immediate to check that $`[r,d1+1d]=cddc`$ for all $`d𝔡^{}`$. Moreover, the element $`x`$ from (15.16) equals $`nc`$. Then, if $`\beta =\beta ^{}+\beta _0`$ with $`\beta ^{}𝔡^{}`$, $`\beta _0𝔡_0`$, equation (15.20) becomes:
$$\beta ^{}cc\beta ^{}=(n+3)(\beta cc\beta ).$$
All solutions $`\beta `$ of this equation are multiples of $`c`$. Trivial cocycles are multiples of $`2sx=(n+2)c`$, hence all cocycles are trivial. ∎
###### Proposition 15.11.
Let $`𝔡^{}𝔡`$ be abelian finite-dimensional Lie algebras such that $`dim𝔡^{}>2`$, let $`H=U(𝔡),H^{}=U(𝔡^{})`$, and let $`L=\mathrm{Cur}_H^{}^HS(𝔡^{},0)`$. Then $`\mathrm{H}^2(L,𝐤)=0`$.
###### Proof.
By Proposition 8.4, $`L`$ is spanned over $`H`$ by elements
$$e_{ab}=abba,a,b𝔡^{},$$
satisfying the relations $`e_{ab}=e_{ba}`$ and
(15.23)
$$ae_{bc}+be_{ca}+ce_{ab}=0.$$
The pseudobrackets are (see (8.24)):
$`[e_{ab}e_{cd}]`$ $`=(ad)_He_{bc}+(bc)_He_{ad}(ac)_He_{bd}(bd)_He_{ac},`$
and in particular
$`[e_{ab}e_{ac}]`$ $`=(ac)_He_{ab}+(ba)_He_{ac}(aa)_He_{bc},`$
$`[e_{ab}e_{ab}]`$ $`=(baab)_He_{ab}.`$
Trivial cocycles $`\tau _\varphi `$ are determined by the identity (see (15.15)):
$$\begin{array}{c}(\tau _\varphi (e_{ab},e_{cd})1)_H1\hfill \\ \hfill =(ad)_H\varphi _{bc}+(bc)_H\varphi _{ad}(ac)_H\varphi _{bd}(bd)_H\varphi _{ac},\end{array}$$
where $`\varphi _{ab}=\varphi (e_{ab})=\varphi _{ba}𝐤`$, which gives:
(15.24)
$$\tau _\varphi (e_{ab},e_{cd})=ad\varphi _{bc}bc\varphi _{ad}+ac\varphi _{bd}+bd\varphi _{ac}.$$
Let $`\beta `$ be a cocycle for $`L`$ representing a central extension. Write $`\beta _{ab,cd}=\beta (e_{ab},e_{cd})`$ for short. Equations (15.14), (15.23) give the identities:
(15.25)
$$\beta _{ab,cd}=\beta _{ba,cd}=\beta _{ab,dc}=S(\beta _{cd,ab}),$$
(15.26)
$$a\beta _{bc,cd}+b\beta _{ca,cd}+c\beta _{ab,cd}=0.$$
Using this, Jacobi identity for the elements $`e_{ab},e_{ab},e_{ac}`$ gives the following equation for $`\beta `$:
(15.27)
$$\begin{array}{cc}\hfill (baab)(\mathrm{\Delta }(\beta _{ab,ac})\beta _{ab,ac}1& 1\beta _{ab,ac})\hfill \\ \hfill =ab\beta _{ab,ac}\beta _{ab,ac}ab& +\beta _{ab,bc}a^2a^2\beta _{ab,bc}\hfill \\ & +\beta _{ab,ab}acac\beta _{ab,ab}.\hfill \end{array}$$
This is a homogeneous equation and can be solved degree by degree. If $`\beta `$ is homogeneous of degree one, then the left-hand side is zero, and we immediately see that $`\beta _{ab,ac}=\beta _{ab,bc}=\beta _{ab,ab}=0`$. Then, by (15.26), $`\beta _{ab,cd}=0`$.
If $`\beta `$ is homogeneous of degree other than one, then $`\beta _{ab,ab}=0`$, since $`\beta `$ restricts to a cocycle of the free rank one Lie pseudoalgebra $`He_{ab}`$, which has been shown in Proposition 15.9 to take values in $`𝔡`$. Then equations (15.25), (15.26) give $`a\beta _{ab,bc}=b\beta _{ab,ac}`$. Hence if $`a`$ and $`b`$ are linearly independent, we can find some $`p=p_{bc}^aH`$ such that
(15.28)
$$\beta _{ab,ac}=ap_{bc}^a,\beta _{ab,bc}=bp_{bc}^a.$$
We substitute this into (15.27) and after simplification obtain $`\mathrm{\Delta }(p)=p1+1p`$. Therefore, $`p𝔡`$, hence the only nonzero solutions to (15.27) occur in degree two.
Now using (15.26) and (15.28), we get:
(15.29)
$$\beta _{ab,cd}=ap_{bd}^cbp_{ad}^c.$$
The skew-symmetry $`\beta _{ab,cd}=\beta _{cd,ab}`$ gives the equations $`p_{bc}^a=p_{cb}^a`$ and
(15.30)
$$ap_{bd}^cbp_{ad}^c=cp_{bd}^adp_{bc}^a.$$
From this we obtain that $`p_{bd}^a`$ lies in the linear span of $`a,b,d`$. Comparing the coefficients in front of $`ac`$ in (15.30), we see that the coefficient of $`a`$ in $`p_{bd}^a`$ is equal to the coefficient of $`c`$ in $`p_{bd}^c`$. Call this coefficient $`\varphi _{bd}`$; then $`\varphi _{bd}=\varphi _{db}`$. Then comparison of other coefficients in (15.30) shows that
(15.31)
$$p_{bc}^a=a\varphi _{bc}+b\varphi _{ca}+c\varphi _{ab}$$
for all $`a,b,c𝔡^{}`$. Substitute this in (15.29) to obtain that $`\beta =\tau _\varphi `$ is trivial (cf. (15.24)). ∎
###### Proposition 15.12.
Let $`H=U(𝔡)`$, and let $`𝔤`$ be a simple finite-dimensional Lie algebra. If $`L=\mathrm{Cur}𝔤`$, then $`\mathrm{H}^2(L,𝐤)𝔡`$.
###### Proof.
Let $`\beta `$ be a cocycle for $`L`$. We will write $`\beta (a,b)=\beta (1a,1b)`$ for $`a,b𝔤`$. Then Jacobi identity leads to the equation
(15.32)
$$\beta (a,[b,c])11\beta (b,[a,c])=\mathrm{\Delta }\left(\beta ([a,b],c)\right).$$
This immediately implies: $`\beta ([a,b],c)𝔡+𝐤`$. Since $`[𝔤,𝔤]=𝔤`$ this shows that $`\beta (a,b)𝔡+𝐤`$ for all $`a,b𝔤`$.
We can now solve the homogeneous equation (15.32) degree by degree. Solutions of degree zero are cocycles of the Lie algebra $`𝔤`$, hence they are all trivial. Solutions of degree one satisfy $`\beta (a,[b,c])=\beta ([a,b],c)`$, and skew-symmetry implies $`\beta (a,b)=S\left(\beta (b,a)\right)=\beta (b,a)`$. Therefore every such $`\beta `$ is an invariant symmetric bilinear form on $`𝔤`$ with values in $`𝔡`$. Any such bilinear form can be written as $`\beta (a,b)=(a|b)d`$ where $`(|)`$ is the Killing form on $`𝔤`$ and $`d`$ is some element of $`𝔡`$. Such cocycles $`\beta `$ are nontrivial, hence inequivalent central extensions are in one-to-one correspondence with elements of $`𝔡`$. ∎
###### Theorem 15.13.
Let $`H=U(𝔡)`$ and $`L`$ be a simple Lie $`H`$-pseudoalgebra. Then $`L`$ may have a nontrivial central extension (given by (15.12), (15.13)) only if:
(i) $`L=\mathrm{Cur}𝔤`$, in which case $`\mathrm{H}^2(L,𝐤)𝔡`$ and cocycles $`\beta `$ are given by $`\beta _d(1a,1b)=(a|b)d`$ for $`a,b𝔤`$, where $`d𝔡`$ and $`(|)`$ is the Killing form.
(ii) $`L=\mathrm{Cur}_H^{}^HW(𝔡^{})`$ with $`𝐤s=𝔡^{}𝔡`$, $`dim𝔡^{}=1`$, in which case $`\mathrm{H}^2(L,𝐤)`$ is of dimension one, generated by the Virasoro cocycle $`\beta (1s,1s)=s^3`$.
(iii) $`L=\mathrm{Cur}_H^{}^HH(𝔡^{},\chi ^{},\omega ^{})`$ with $`𝔡^{}𝔡`$, in which case $`\mathrm{H}^2(L,𝐤)`$ is isomorphic to the quotient of the space of all solutions $`\beta 𝔡`$ to equations (15.19), (15.20) by the subspace $`𝐤(2sx)`$, where $`r𝔡^{}𝔡^{}`$ is dual to $`\omega ^{}`$, $`s𝔡^{}`$ is such that $`\chi ^{}=\iota _s\omega ^{}`$, and $`x`$ is given by (15.16).
###### Proof.
(i) and (ii) follow from Propositions 15.12 and 15.8(i), and (iii) from a direct application of Lemma 15.7.
For any other simple pseudoalgebra $`L`$, the strategy is to construct a continuous family of pseudoalgebras $`L_t`$, indexed by $`t𝐤`$ endowed with the Zariski topology, that are all isomorphic to $`L`$ when $`t0`$, and whose fiber at $`t=0`$ is one of the pseudoalgebras already considered in Propositions 15.8(ii), 15.10 and 15.11. Then, since $`\mathrm{H}^2(L_t,𝐤)=0`$ for $`t=0`$, it will follow that $`\mathrm{H}^2(L_t,𝐤)=0`$ whenever $`t`$ lies in a neighborhood of $`0`$, hence for all $`t𝐤`$.
In the case of a current pseudoalgebra over a $`W`$ or $`S`$ type Lie pseudoalgebra, choose a basis $`\{_i\}`$ of $`𝔡`$ that contains a basis of $`𝔡^{}`$, and construct the family $`𝔡_t^{}𝔡_t`$ of Lie algebras indexed by $`t𝐤`$ generated by elements $`\{_i^t\}`$ with Lie bracket $`[_i^t,_j^t]=t[_i,_j]^t`$. Then for $`t0`$ we have an isomorphism $`𝔡_t_i^tt_i𝔡`$, whereas $`𝔡_0`$ is an abelian Lie algebra. Then $`\{\mathrm{Cur}_{H_t^{}}^{H_t}W(𝔡_t^{})\}_{t𝐤}`$, where $`H_t=U(𝔡_t)`$, $`H_t^{}=U(𝔡_t^{})`$, is a family of pseudoalgebras all isomorphic to $`\mathrm{Cur}_H^{}^HW(𝔡^{})`$ for $`t0`$. The fiber of this family at $`t=0`$ has been shown in Proposition 15.8(ii) to have no nontrivial central extensions. In the same way, if we set $`\chi _t(_i^t)=t\chi (_i)`$, then $`\{\mathrm{Cur}_{H_t^{}}^{H_t}S(𝔡_t^{},\chi _t)\}_{t𝐤}`$ is a family of pseudoalgebras all isomorphic to $`\mathrm{Cur}_H^{}^HS(𝔡^{},\chi )`$ for $`t0`$, and the fiber at $`t=0`$ is $`\mathrm{Cur}_{H_0^{}}^{H_0}S(𝔡_0^{},0)`$ where $`𝔡_0^{}𝔡_0`$ are isomorphic to $`𝔡^{}𝔡`$ as vector spaces but have trivial bracket.
If $`L`$ is a current pseudoalgebra over $`K(𝔡^{},\theta )`$, for finite-dimensional Lie algebras $`𝔡^{}𝔡`$, choose a basis $`\{a_i,b_i,s\}`$ of $`𝔡^{}`$ as in Lemma 8.9, and complete it with $`\{d_1,\mathrm{},d_r\}`$ to a basis of $`𝔡`$. Then a continuous family $`\{𝔡_t\}`$ of Lie algebras can be constructed for $`t0`$ as $`𝔡_t𝔡`$ spanned by $`a_i^t=ta_i`$, $`b_i^t=tb_i`$, $`s^t=t^2s`$, $`d_i^t=t^2d_i`$, and by setting $`a_i^0,b_i^0,c^0=s^0`$ to span a Heisenberg algebra, and all brackets involving $`d_i^0`$ to be zero. Define $`\theta _t(𝔡_t^{})^{}`$ by $`\theta _t(a_i^t)=\theta _t(b_i^t)=0`$, $`\theta _t(s^t)=1`$. Then $`\mathrm{Cur}_{H_0^{}}^{H_0}K(𝔡_0^{},\theta _0)`$ is the limit of the Lie pseudoalgebras $`\{\mathrm{Cur}_{H_t^{}}^{H_t}K(𝔡_t^{},\theta _t)\}_{t0}`$, which are all isomorphic to $`\mathrm{Cur}_H^{}^HK(𝔡^{},\theta )`$, and the former Lie pseudoalgebra is of the type treated in Proposition 15.10. ∎
## 16. Application to the Classification of Poisson Brackets in Calculus of Variations
In calculus of variations the phase space consists of $`C^{\mathrm{}}`$ vector functions $`u=(u_1(x),\mathrm{},u_r(x))`$ where $`u_i(x)`$ are, for example, functions with compact support on a closed $`N`$-dimensional manifold. We will consider linear local Poisson brackets:
(16.1)
$$\{u_i(x),u_j(y)\}=\underset{\alpha }{}B_{\alpha ij}(y)_y^\alpha \delta (xy)$$
where the sum runs over a finite set of multi-indices $`\alpha =(\alpha _1,\mathrm{},\alpha _N)_+^N`$, the $`B_{\alpha ij}`$ are linear combinations of the $`u_k`$ and of their derivatives $`u_k^{(\gamma )}:=_x^\gamma u_k`$, where $`_x^\gamma :=\left(\frac{}{x_1}\right)^{\gamma _1}\mathrm{}\left(\frac{}{x_N}\right)^{\gamma _N}`$, and $`\delta (xy)`$ is the delta-function (defined by $`f(x)\delta (xy)dx=f(y)`$). By Leibniz rule and bilinearity, the Poisson bracket (16.1) extends to arbitrary polynomials $`P`$ and $`Q`$ in the $`u_i`$ and their derivatives. Explicitly:
(16.2)
$$\{P(x),Q(y)\}=\underset{\alpha ,\beta ,i,j}{}\frac{P(x)}{u_i^{(\alpha )}}\frac{Q(y)}{u_j^{(\beta )}}_x^\alpha _y^\beta \{u_i(x),u_j(y)\}.$$
This bracket, apart from bilinearity and Leibniz rule, should satisfy skew-commutativity and the Jacobi identity.
The basic quantities in calculus of variations are local functionals (Hamiltonians) $`I_P=P(x)dx`$. Using bilinearity and integration by parts ($`\frac{P}{x_i}Qdx=P\frac{Q}{x_i}dx`$), we get from (16.2) the following well-known formula:
(16.3)
$$\{I_P,I_Q\}=\underset{i,j}{}\frac{\delta P(x)}{\delta u_i}\frac{\delta Q(y)}{\delta u_j}\{u_i(x),u_j(y)\}dxdy,$$
where
(16.4)
$$\frac{\delta P(x)}{\delta u_i}=\underset{\alpha }{}(_x)^\alpha \frac{P(x)}{u_i^{(\alpha )}}$$
is the variational derivative.
More generally, one usually considers a class of Poisson brackets of the form:
(16.5)
$$\{u_i(x),u_j(y)\}=B_{ij}(y,_y^\alpha )\delta (xy),$$
where $`B_{ij}`$ are differential operators in $`_y^\alpha `$ whose coefficients are polynomials in $`u_k^{(\gamma )}(y)`$. Then the $`r\times r`$ matrix $`B=(B_{ij})`$ is called a Hamiltonian operator, and, integrating by parts, formula (16.3) can be rewritten in its most familiar form:
(16.6)
$$\{I_P,I_Q\}=\left(B\frac{\delta P(x)}{\delta u}\right)\frac{\delta Q(x)}{\delta u}dx.$$
Given a Hamiltonian $`h=P(x)dx`$, we have the corresponding Hamiltonian system of evolutionary partial differential equations:
(16.7)
$$\dot{u}_i=\{h,u_i\}\underset{j}{}B_{ij}\frac{\delta P}{\delta u_j},$$
so that if another Hamiltonian $`h_1`$ is in involution with $`h`$, i.e., $`\{h,h_1\}=0`$, then $`h_1`$ is an integral of motion of (16.7), i.e., $`\dot{h}_1=0`$.
It is shown in \[DN1\] and \[M\] that for $`r2`$, any Poisson bracket of hydrodynamic type (i.e., linear in the derivatives) under certain nondegeneracy conditions can be transformed into a linear Poisson bracket of hydrodynamic type by a change of the field variables. The latter Poisson brackets have been studied rather extensively (see \[Do\], \[DN2\] and references there, \[GD\], \[M\], \[Z\]).
Let $`H=[\frac{}{x_1},\mathrm{},\frac{}{x_N}]`$ be the universal enveloping algebra of the $`N`$-dimensional abelian Lie algebra $`𝔡`$. Let $`F=_{i=1}^rHu_i`$ be the free $`H`$-module of rank $`r`$ on generators $`u_i`$. Consider a linear Poisson bracket (16.1). For any three subspaces $`A,B,C`$ of $`F`$, we will use the notation $`\{A,B\}C`$ if for any $`aA`$, $`bB`$ all coefficients in front of $`_y^\alpha \delta (xy)`$ in $`\{a(x),b(y)\}`$ belong to $`C`$. We call a linear Poisson algebra any $`H`$-submodule $`L`$ of $`F`$ which is closed under the Poisson bracket, i.e., such that $`\{L,L\}L`$. By an isomorphism of two such algebras we mean a $``$–linear isomorphism preserving Poisson brackets.
If $`L`$ is a linear Poisson algebra, we define the $`\lambda `$-bracket ($`\lambda =(\lambda _1,\mathrm{},\lambda _N)`$) on $`L`$ as the Fourier transform of the linear Poisson bracket (16.1):
(16.8)
$$[u_i{}_{\lambda }{}^{}u_{j}^{}]=_\alpha \lambda ^\alpha B_{\alpha ij}.$$
Then we get a Lie conformal algebra in $`N`$ (commuting) indeterminates (defined in the same way as for the $`N=1`$ case in the introduction). Thus, the classification of linear Poisson algebras follows from the classification of Lie $`U(𝔡)`$-conformal algebras, where $`𝔡`$ is the $`N`$-dimensional abelian Lie algebra.
Recall that the structure of a Lie conformal algebra is equivalent to the structure of a Lie pseudoalgebra (see Section 9). The relationship between the linear Poisson bracket (16.1) and the pseudobracket can be described explicitly as follows:
(16.9) $`[u_iu_j]`$ $`=_kP_{ij}^k(1,1)_Hu_k,`$
if
(16.10) $`\{u_i(x),u_j(y)\}`$ $`=_kP_{ij}^k(_x,_y)\left(u_k(y)\delta (xy)\right)`$
for some polynomials $`P_{ij}^k`$. Note that equation (16.1) can always be written in the form (16.10). Indeed, if
(16.11) $`\{u_i(x),u_j(y)\}`$ $`=_kQ_{ij}^k(_y,_t)\left(u_k(t)\delta (xy)\right)|_{t=y}`$
for some polynomials $`Q_{ij}^k`$, then we have (16.10) with $`P_{ij}^k(z,w)=Q_{ij}^k(z,z+w)`$. In this case, the $`\lambda `$-bracket (16.8) is given by:
(16.12) $`\{u_i(x),u_j(y)\}`$ $`=_kQ_{ij}^k(\lambda ,)u_k.`$
###### Remark 16.1.
The constant terms of $`B_{\alpha ij}`$ in (16.1) give a central extension of the linear Poisson algebra corresponding to the $`B_{\alpha ij}`$ with constant terms removed. In terms of the associated Lie pseudoalgebras this corresponds to a central extension by $``$ with a trivial action of $`H`$. By Theorem 15.1, these central extensions are parameterized by $`\mathrm{H}^2(L,)`$.
###### Examples 16.2 (cf. \[M\], \[DN2\], \[K4\], \[BKV\]).
1. General Poisson algebra $`W_{r,N}`$, where $`1rN`$ $`(1i,jr)`$:
$$\{u_i(x),u_j(y)\}=\frac{u_j(y)}{y_i}\delta (xy)+u_j(y)\frac{}{y_i}\delta (xy)+u_i(y)\frac{}{y_j}\delta (xy).$$
2. Special Poisson algebra $`S_{r,N,\chi },`$ where $`2rN`$ and $`\chi =(\chi _1,\mathrm{},\chi _r)^r`$, is the following subalgebra of $`W_{r,N}`$:
$$\left\{\underset{i=1}{\overset{r}{}}P_i(_x)u_i(x)\right|\underset{i=1}{\overset{r}{}}\left(\frac{}{x_i}+\chi _i\right)P_i(_x)=0\}.$$
It is generated over $`H=[\frac{}{x_1},\mathrm{},\frac{}{x_N}]`$ by elements
$$u_{ij}(x)=\left(\frac{}{x_i}+\chi _i\right)u_j(x)\left(\frac{}{x_j}+\chi _j\right)u_i(x).$$
3. Hamiltonian Poisson algebra $`H_{2s,N}`$, $`2r=2sN`$:
$$\{u(x),u(y)\}=\underset{i=1}{\overset{s}{}}\left(\frac{u(y)}{y_i}\frac{}{y_{i+s}}\delta (xy)\frac{u(y)}{y_{i+s}}\frac{}{y_i}\delta (xy)\right).$$
We have an inclusion $`H_{2s,N}W_{2s,N}`$ by letting
$$u(x)=\underset{i=1}{\overset{s}{}}\left(\frac{u_{i+s}(x)}{x_i}\frac{u_i(x)}{x_{i+s}}\right).$$
4. Current Poisson algebra $`\mathrm{Cur}_N𝔤`$ associated to a simple $`r`$-dimensional Lie algebra $`𝔤`$ with structure constants $`c_{ij}^k`$ $`(1i,j,kr)`$:
$$\{v_i(x),v_j(y)\}=\underset{k=1}{\overset{r}{}}c_{ij}^kv_k(y)\delta (xy).$$
5. Semidirect sum of $`W_{r,N}`$ or one of its subalgebras $`S_{r,N,\chi },H_{2s,N}`$ with $`\mathrm{Cur}_N𝔤`$ defined by ($`1ir,v(x)\mathrm{Cur}_N𝔤`$):
$$\{u_i(x),v(y)\}=\frac{v(y)}{y_i}\delta (xy)+v(y)\frac{}{y_i}\delta (xy).$$
A subspace $`I`$ of a Poisson algebra $`L`$ is called an ideal if it is invariant under taking Poisson brackets with elements of $`L`$, i.e., if $`\{L,I\}I`$. A Poisson algebra $`L`$ is called simple (respectively semisimple) if the Poisson bracket is not identically zero and $`L`$ contains no nonzero $`H`$-invariant ideals $`I`$ such that $`IL`$ (respectively $`\{I,I\}0`$). Note that the Poisson algebras that we consider here are finite, i.e., finitely generated as $`H`$-modules.
Then Theorems 13.10 and 13.15 and Corollary 13.27 imply:
###### Theorem 16.3.
(i) Any simple linear Poisson algebra is isomorphic to one of the Poisson algebras $`W_{r,N},S_{r,N,\chi },H_{2s,N},\mathrm{Cur}_N𝔤`$.
(ii) Any semisimple linear Poisson algebra is a direct sum of simple ones and of the semidirect sums described in Example 16.2(5).
###### Remark 16.4.
It follows from Remark 16.1 and the results of Section 15.4 that all nontrivial central extensions of simple Poisson algebras are described by the following $`2`$-cocycles.
For $`\alpha ^r`$ let
$$\psi _\alpha (x,y)=\underset{i=1}{\overset{r}{}}\alpha _i\frac{}{y_i}\delta (xy).$$
Then all nontrivial $`2`$-cocycles for $`H_{r,N}`$ are:
$$\gamma _\alpha (u(x),u(y))=\psi _\alpha (x,y).$$
All nontrivial $`2`$-cocycles for $`\mathrm{Cur}_N𝔤`$ are:
$$\gamma _\alpha (v_i(x),v_j(y))=b_{ij}\psi _\alpha (x,y),$$
where $`b_{ij}=(v_i|v_j)`$ is the invariant scalar product.
The Poisson algebra $`W_{r,N}`$ has a nontrivial central extension iff $`r=1`$, and in the latter case it is given by the well-known Virasoro cocycle:
$$\gamma (u(x),u(y))=\left(\frac{}{y}\right)^3\delta (xy).$$
The Poisson algebra $`S_{r,N,\chi }`$ has no nontrivial central extensions if $`r>2`$ or $`\chi 0`$, and $`S_{2,N,0}H_{2,N}`$ has nontrivial central extensions described above.
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# 1 Introduction
## 1 Introduction
A definite pattern for the theoretical description of particle physics has emerged based on the framework of the Standard Model (SM) of high energy physics \[4-22\], which is built up from observation for prediction and correlation of the new data. Although the SM has proven to be in spectacular agreement with experimental measurements and quite successful in a predicting a wide range of phenomena, however, it is not exception to the rule that as the phenomenological approach it suffers from own difficulties. There were still many open key questions arisen inevitably that we have no understanding why the SM is as it is? Why is the gauge symmetry? Why is this the particle spectrum? The mechanism of the electroweak symmetry breaking is a complete mystery. The most problematic ingredient of such a breaking is the Higgs boson (in simplest version), which has not yet been discovered experimentally. The untested aspects of SM are the mass spectrum of the particles, the mixing patterns and the CP violation. The latter is introduced through complex Yukawa couplings of fermions to Higgs bosons, resulting in complex parameters in the CKM matrix. The SM contains a large number of a priori free parameters, while a consistent complete theory would not have so many free parameters. It is important then to develop the other schemes that attempt to reduce the number of free parameters. All the approaches proposed towards the unified gauge field theory, e.g. \[11-42\], have own advantages and difficulties. But, it is for one thing, the basic concepts and right symmetries, as a rule, were inserted in all of them in ad hoc manner. However, their physical essence and origin still remain unknown and, therefore, the resulting theory is beset by various difficulties. The key of solution rested entirely in the searches for the new mathematical structures enabling an insight, in microscopic sense, to the conceptual foundation of particle physics.
To address to some of the above mentioned nagging questions of the SM recently the MSM is developed in , wherein the proliferation of lepton and quark flavours prompts us to consider the fields as composites. Certainly, it may seem foolhardy to set up such a picture in the spacetime continuum. The difficulties arisen here are well-known. The first problem is closely related to the expected mass differences of particles, which in this case would be too large ($`1TeV`$). Another problem concerns the transformations of particles. Our idea is to remove these difficulties by employing the multiworld (MW) geometry ensued from the operator manifold (OM) formalism . The revised, expanded and completed version of is suggested in the present article, which is organized as follows: The OM formalism is given in the part I, wherein to facilitate writing the main new features of the new version are firstly represented in the restructuring of the old one. We forbear to write out in the main text the pieces including only the technical issues, which have now replaced into corresponding appendices as they are a somewhat lengthy and so standard. Both the quantum field and differential geometric aspects of the OM formalism are studied. The OM enables to develop an approach to the unification of the geometry and the field theory. The former is equivalent to the configuration space wave mechanics incorporated with geometric properties leading to the quantization of geometry strongly different from earlier suggested schemes (subsec.2.1, App.A). Within the framework of algebraic approach we have reached to rigorous definition of the OM. While, we consider one parameter group of operator diffeomorphisms, the operator differential forms and their integration, also operator exterior differentiation. We generalize this formalism via the concept of operator multimanifold (OMM) (sec.5), which yields the MW geometry (subsec.5.2) involving the spacetime continuum and internal worlds of the given number. In an enlarged framework of the OMM we define and clarify the conceptual basis of subquarks and their characteristics stemming from the various symmetries of the internal worlds (subsec.5.3). The OMM formalism has the following features:
$``$ It provides a natural unification of the geometry- yielding the 1) special and 2) general relativity principles, and the fermion fields serving as the basis for the constituent subquarks (subsec.5.3).
$``$ It has cleared up the physical conditions in which the geometry and particles come into being (subsec.2.2, 5.2).
$``$ The subquarks emerge in the geometry only in certain permissible combinations utilizing the idea of the 3) subcolour (subquark) confinement principle, and have undergone the transformations yielding the internal symmetries and 4) gauge principle.
These results are used in the part II to develop the updated version of the MSM . The main new features of the new version are as follows: 1) In dealing with a new round of the experiments, in this paper, we supplement this previous discussion, first by including a detailed description of the remarkable specific mechanism of electroweak symmetry breaking arisen in the MSM (sec.8); 2) The two solid phenomenological testable implications of the MSM are given in the sec.9; 3) Finally, we derive a physically more realistic mass spectrum of the leptons and quarks (sec.12) instead of the former one inferred within the simplified scheme .
All the fields including the leptons and quarks along with the spacetime component have also the nontrivial composite internal MW-structure (sec.6). While, the various subquarks then are defined on the corresponding internal worlds. The possible elementary particles are thought to be composite dynamical systems in analogy to quantum mechanical stationary states of compound atom, but, now a dynamical treatment built up on the MW geometry is quite different and more amenable to qualitative understanding. The microscopic structure of leptons, quarks and other particles will be governed by the only possible conjunctions of constituent subquarks implying concrete symmetries. Although within considered schemes the subquarks are defined on the internal worlds, however the resulting spacetime components of particles, which we are going to deal with to describe the leptons and quarks defined on the spacetime continuum, are affected by them (subsec.6.4) in such a way that they carry exactly all the quantum numbers of the various constituent subquarks of the given composition. The MSM enables an insight to the key concepts of particle physics (subsec.6.1-6.6), and to conclude that the leptons are particles with integer electric and leptonic charges and free of confinement, while the quarks carry fractional electric and baryonic charges and imply the confinement. The theoretical significance of the MSM resides in the microscopic interpretation of all physical parameters. We derive the Gell-Mann-Nishijima relation and the flavour group. The MW structure of primary field (subsec.6.6, 6.7) is described by the gauge invariant Lagrangian involving nonlinear fermion interactions of the internal field components somewhat similar to the theory by Heisenberg and his co-workers , but still it will be defined on the MW geometry. This Lagrangian is the whole story since all the complexity of the leptons, quarks and their interactions arises from it. The number of free parameters in this Lagrangian is reduced to primary coupling constant of the nonlinear interaction and gauge coupling. Based on it, we consider the microscopic theory of the unified electroweak interactions (subsec.6.8-sub.12). It follows that contemporary phenomenological SM is an approximation to the suggested microscopic approach. We exploit the background of the local symmetry $`SU^{loc}(2)U^{loc}(1)`$ (subsec.6.8), the weak hypercharge and P (mirror symmetry)-violation. The Weinberg mixing angle determining the symmetry reduction coefficient is shown to have a value fixed at $`30^o`$ (subsec.6.9). We develop the microscopic approach to the isospinor Higgs boson with self-interaction and Yukawa couplings (subsec.6.10, sec.7). It involves Higgs boson as the collective excitations of bound quasi-particle pair. Tracing a resemblance with the Cooper pairs \[53-55\], within the framework of local gauge invariance of the theory incorporated with the phenomenon of P-violation in weak interactions we suggest a mechanism providing the Bose-condensation of relativistic fermion pairs. In contrast to the SM, it predicts both the electroweak symmetry breaking in the $`W`$-world by the vacuum expectation value (VEV) of spin zero Higgs bosons and its transmission from the $`W`$world to the spacetime continuum (sec.8). This is the most remarkable feature of our approach especially in the view of existing great belief of the conventional theories for a discovery of the Higgs boson with other new particles at next round of experiments at LEP2, at the Tevatron, at LHC and other colliders, which will explore the TeV energy range (e.g. see ). The LEP2 data (is currently running at 189 GeV) provide a lower limit $`m_H>89.3GeV`$ on its mass in simplest version. Furthermore, there is a tight upper limit $`(m_{h^0}<150GeV)`$ on the mass of the lightest Higgs boson $`h^0`$ among the 5 physical Higgs bosons predicted by the models of Minimal Supersymmetric Exstention of the SM (MSSM). The current direct search limits from LEP2 give $`m_{h^0}>75GeV.`$ Therefore, the future searches for this boson (if the mass is below 150 GeV or so) would be a crucial point in testing the efforts made in the conventional models building as well in the present MSM based on a quite different approaches. Actually, reflecting upon the results far obtained in the sec.9, in strong contrast to conventional theories, the MSM rejects drastically any expectation of discovery of the Higgs boson, but in the same time it expects to include a rich spectrum of new particles at higher energies. Thus, if the MSM proves viable it becomes an crucial issue to hold in experiments the following two solid tests:
$``$ The Higgs bosons never could emerge in spacetime continuum since they have arisen only on the internal $`W`$-world, i.e., thus, the unobserved effects produced by such bosons cannot be discovered in experiments nor at any energy range.
$``$ For each of the three Standard Model families of quarks and leptons, there are corresponding heavy family partners with the same quantum numbers lying far above the electroweak scale.
Regarding to the last phenomenological implication of the MSM, it is remarkable that the similar in many respects prediction is made in somewhat different context by S.L.Adler within a phenomenological scheme of a compositeness of the quarks and leptons. It based on the generic group theoretical framework of rishon type models exploring the preon constituents. But, therein the most important specification of this scale is absent. Although one admits that such a scale could be much higher than electroweak scale, however, it is also necessary special argumentations in support of validity of this prediction in the case if this scale has turned out to be low enough, namely, if these heavy partners lie not too far above the electroweak scale. Even thus, as it is notified in , one must not worry for the existence of 6 heavy flavors, which is then marginally compatible with the current LEP data . A complete analysis of this question, naturally, is possible now in suggested microscopic approach. The MSM enables oneself to study in detail the phenomenology associated with such extra heavy families and to estimate the value of energy threshold of their creation. While, the low energy scale could not be realized since it lies far below the energy threshold of the next pole necessary for appearing of the heavy partners. The estimate gives the common mass-shift coefficients $`(1+k)`$, where $`k`$ reads for the next few low energy poles with respect to the lowest one: $`k_0=0,k_1>\sqrt{2},k_2=\sqrt{8/3}`$ and $`k_3=2`$. The first one obviously does not produce the extra families, but the energy thresholds corresponded to the next non-trivial poles can be respectively written: $`E_1>(419.6\pm 12.0)GeV,E_2=(457.6\pm 13.2)GeV`$ and $`E_3=(521.4\pm 15.0)GeV.`$ Thus, which of these schemes above, if any, is realized either exactly or at least approximately in nature remains to be seen in the years to come.
Finally we attempt to predict the mixing angles in the six-quark KM model (sec.10), the appearance of the CP-violation phase (sec.11) and derive the mass spectrum of the leptons and quarks (sec.12). The physical outlook on suggested approach and concluding remarks are given in the sec.13, which have involved in order once again to resume a whole physical picture and to provide a sufficient background for its understanding without undue hardship. The appendices will complete the mathematical framework.
Our approach still should be considered as a preliminary one, wherein we have contended ourselves with a rather modest task and do not profess to have any clear-cut answers to all the problems of particle physics, the complete picture of which is largely beyond the scope of the present paper. The only argument that prompts us to consider present approach seriously is the remarkable feature that the most important properties of particle phenomenology can be derived naturally within its framework. Therefore we hope that it will be an attractive basis for the future theory. Although many key problems are elucidated within outlined approach, nevertheless some issues still remain to be solved. To go any further in exploring the significance of obtained results it is entirely feasible, for instance, to promote the MSM into the supersymmetric framework in order to solve its technical aspects of the vacuum zero point energy and hierarchy problems, and attempt to develop the realistic viable microscopic theory (VMSM), which will be carried out in subsequent paper. Given therein the VMSM will make a new predictions about observing of the supersymmetric partners drastically different from those of conventional MSSM-models.
Part I
Introduction to the Operator Manifold Formalism
The mathematical framework of the OM formalism reveals primordial deeper structures underlying the fundamental concepts of the particle physics. Here we explore the query how did the geometry and fields, as they are, come into being. In the first our major purpose is to prove the idea that the geometry and fields, with the internal symmetries and all interactions, as well the four major principles of relativity (special and general), quantum, gauge and colour confinement are derivative, and they come into being simultaneously. The substance out of which the geometry and fields are made is the “primordial structures” involved into reciprocal “linkage” establishing processes (sec.4)
## 2 Preliminaries
This section contains some of the necessary preliminaries on generic of the OM formalism , which one to know in order to understand our approach. We adopt then all the ideas and conventions of and to be brief we often suppress the indices without notice.
We start by tracing at elementary level the relevant steps of motivation of the OM formalism:
$``$ First step is an extension of the Minkowski space $`\underset{x}{𝑀}{}_{4}{}^{}M_8=\underset{x}{𝑀}{}_{4}{}^{}\underset{u}{𝑀}_4`$ in order to introduce the particle mass operator defined on the internal world $`\underset{u}{𝑀}_4`$. For example, in the case of Dirac’s particle one proceeds at once:
$$(\underset{x}{\underset{}{\gamma p}}m)\underset{x}{\Psi }=0\gamma p\psi =0,$$
provided by $`\psi =\underset{x}{\Psi }\underset{u}{\Psi },\gamma p=\underset{x}{\underset{}{\gamma p}}\underset{u}{\underset{}{\gamma p}},m\underset{u}{\Psi }\underset{u}{\underset{}{\gamma p}}\underset{u}{\Psi }`$ and
$$dx^2=invdx_8^2=dx^2du^2=0,x_8M_8.$$
The same holds for the other fields of arbitrary spin.
$``$ Next, a two-steps passage $`M_4M_6\stackrel{45^0}{}G_6`$ will be performed for each sample of the $`M_4`$.
a) A passage $`M_4M_6`$ restores the complete equivalence between the three spacial and three time components:
$$\begin{array}{c}e_4=(\stackrel{}{e},e_0)\stackrel{}{e}_6=(\stackrel{}{e},\stackrel{}{e}_0)M_6,x_4=(\stackrel{}{x},x_0)x_6=(\stackrel{}{x},\stackrel{}{x}_0)M_6.\hfill \end{array}$$
b) A rotation $`M_6\stackrel{45^0}{}G_6`$ of the basis vectors on the angle $`45^0`$ provides an adequate algebra for quantization of the geometry (subsec.2.1):
$$\begin{array}{c}\stackrel{}{e}_6\stackrel{45^0}{}e_{(\lambda \alpha )},\lambda =\pm ,\alpha =1,2,3,\hfill \\ \\ e_{\pm \alpha }=\frac{1}{\sqrt{2}}(e_{0\alpha }\pm e_\alpha )=O_\pm \sigma _\alpha ,<O_\lambda ,O_\tau >=1\delta _{\lambda \tau },<\sigma _\alpha ,\sigma _\beta >=\delta _{\alpha \beta }.\hfill \end{array}$$
Accordingly one gets $`M_8G_{12}`$. Thus, within a simplified scheme (one $`u`$\- channel) of the following it is convenient to deal in terms of smooth differentiable manifold
$$G=\underset{\eta }{𝐺}\underset{u}{𝐺},$$
$`DimG=12,Dim\underset{i}{𝐺}=6(i=\eta ,u)`$.
$``$ Finally, in suggested approach we will be dealing in terms of first degree of the line element, which entails an additional phase multiplier $`\mathrm{\Phi }(\zeta `$) for the vector defined on $`G`$:
$$d\zeta ^2d\stackrel{}{\zeta }e^{iS},\stackrel{}{\zeta }\stackrel{}{\mathrm{\Phi }}(\zeta )=\stackrel{}{\zeta }\mathrm{\Phi }(\zeta ),\mathrm{\Phi }(\zeta )e^{iS},$$
where $`\stackrel{}{\zeta }=\stackrel{}{e}\zeta ,\stackrel{}{e}=(\stackrel{}{\underset{\eta }{𝑒}},\stackrel{}{\underset{u}{𝑒}}),S(\zeta )`$ is the invariant action defined on $`G`$.
### 2.1 Quantization of Geometry
The $`\{e_{(\lambda ,\mu ,\alpha )}=O_{\lambda ,\mu }\sigma _\alpha \}G`$ $`(\lambda ,\mu =1,2;\alpha =1,2,3)`$ are linear independent $`12`$ unit vectors at the point $`p`$ of the 12 dimensional smooth differentiable manifold $`G`$, provided by the linear unit bipseudovectors $`O_{\lambda ,\mu }`$ and the ordinary unit vectors $`\sigma _\alpha `$ implying
$$<O_{\lambda ,\mu },O_{\tau ,\nu }>={}_{}{}^{}\delta _{\lambda ,\tau }^{}{}_{}{}^{}\delta _{\mu ,\nu }^{}<\sigma _\alpha ,\sigma _\beta >=\delta _{\alpha \beta },{}_{}{}^{}\delta =1\delta ,$$
where $`\delta `$ is Kronecker symbol, $`\{O_{\lambda ,\mu }=O_\lambda O_\mu \}`$ is the basis for tangent vectors of $`2\times 2`$ dimensional linear pseudospace $`{}_{}{}^{}𝐑_{}^{4}={}_{}{}^{}𝐑_{}^{2}{}_{}{}^{}𝐑_{}^{2}`$, the $`\sigma _\alpha `$ refer to three dimensional ordinary space $`𝐑^3`$. Henceforth, we always let the first two subscripts in the parentheses to denote the pseudovector components, while the third refers to the ordinary vector components. The metric on $`G`$ is $`\widehat{𝐠}:𝐓_p𝐓_pC^{\mathrm{}}(G)`$ a section of conjugate vector bundle $`S^2𝐓`$. Any vector $`𝐀_p𝐓_p`$ reads $`𝐀=eA`$, provided with components $`A`$ in the basis $`\{e\}`$. In holonomic coordinate basis $`\left(/\zeta \right)_p`$ one gets $`A={\displaystyle \frac{d\zeta }{dt}}|_p`$ and $`\widehat{g}=gd\zeta d\zeta `$. The manifold $`G`$ decomposes as follows:
$$G={}_{}{}^{}𝐑_{}^{2}{}_{}{}^{}𝐑_{}^{2}𝐑^3=\underset{\eta }{𝐺}\underset{u}{𝐺}=\underset{\lambda ,\mu =1}{\overset{2}{}}𝐑_{\lambda \mu }^3=\underset{x}{𝐑}^3\underset{x_0}{𝐑}^3\underset{u}{𝐑}^3\underset{u_0}{𝐑}^3$$
with corresponding basis vectors $`\underset{i}{𝑒}{}_{(\lambda \alpha )}{}^{}=\underset{i}{𝑂}{}_{\lambda }{}^{}\sigma _\alpha \underset{i}{𝐺}`$ $`(\lambda =\pm ,i=\eta ,u)`$ of tangent sections, where
$$\underset{i}{𝑂}{}_{+}{}^{}=\frac{1}{\sqrt{2}}(O_{1,1}+\epsilon _iO_{2,1}),\underset{i}{𝑂}{}_{}{}^{}=\frac{1}{\sqrt{2}}(O_{1,2}+\epsilon _iO_{2,2}),\epsilon _\eta =1,\epsilon _u=1.$$
Hence $`<\underset{i}{𝑂}{}_{\lambda }{}^{},\underset{j}{𝑂}{}_{\tau }{}^{}>=\epsilon _i\delta _{ij}{}_{}{}^{}\delta _{\lambda \tau }^{}`$. The $`\underset{\eta }{𝐺}`$ decomposes into three dimensional ordinary and time flat spaces
$$\underset{\eta }{𝐺}=\underset{x}{𝐑}^3\underset{x_0}{𝐑}^3$$
with the signatures $`sgn(\underset{x}{𝐑}^3)=(+++)`$ and $`sgn(\underset{x_0}{𝐑}^3)=()`$. The same holds for the $`\underset{u}{𝐺}`$ with opposite signatures $`sgn(\underset{u}{𝐑}^3)=()`$ and $`sgn(\underset{u_0}{𝐑}^3)=(+++)`$.
The positive metric forms are defined on manifolds $`\underset{i}{𝐺}:`$ $`\eta ^2\underset{\eta }{𝐺},u^2\underset{u}{𝐺}.`$ The passage to the four-dimensional Minkowski space is a further step as follows: since all directions in $`\underset{x_0}{𝐑}^3`$ are equivalent, then by notion time one implies the projection of time-coordinate on fixed arbitrary universal direction in $`\underset{x_0}{𝐑}^3`$. This clearly respects the physical ground. By such a reduction $`\underset{x_0}{𝐑}^3\underset{x_0}{𝐑}^1`$ the passage
$$\underset{\eta }{𝐺}M_4=\underset{x}{𝐑}^3\underset{x_0}{𝐑}^1$$
may be performed whenever it will be necessary.
In the other case of the six dimensional curved manifold $`\stackrel{~}{G}`$ (subsec.4.2), the passage to four dimensional Riemannian geometry $`R_4`$ is straightforward by making use of reduction of three time components $`e_{0\alpha }={\displaystyle \frac{1}{\sqrt{2}}}(e_{(+\alpha )}+e_{(\alpha )})`$ of basis sixvector $`e_{(\lambda \alpha )}`$ to the single one $`e_0`$ in the given universal direction, which merely has fixed the time coordinate. Actually, since Lagrangian of the fields defined on $`\stackrel{~}{G}`$ is a function of scalars such as $`A_{(\lambda \alpha )}B^{(\lambda \alpha )}=A_{0\alpha }B^{0\alpha }+A_\alpha B^\alpha `$, thus, taking into account that $`A_{0\alpha }B^{0\alpha }=A_{0\alpha }<e^{0\alpha },e^{0\beta }>B_{0\beta }=A_0<e^0,e^0>B_0=A_0B^0`$, one readily may perform the required passage. Hence
$$d\zeta ^2=d\eta ^2du^2=0,d\eta ^2|_{64}ds^2=g_{\mu \nu }dx^\mu dx^\nu =du^2=inv.$$
For more discussion see App.B or .
Unifying the geometry and particles into one framework the OM formalism is analogous to the method of secondary quantization with appropriate expansion over the geometric objects. We proceed at once to the secondary quantization of geometry by substituting the basis elements for the creation and annihilation operators acting in the configuration space of occupation numbers. Instead of pseudo vectors $`O_\lambda `$ we introduce the operators supplied by additional index ($`r`$) referring to the quantum numbers of corresponding state
$$\widehat{O}_1^r=O_1^r\alpha _1,\widehat{O}_2^r=O_2^r\alpha _2,\widehat{O}_r^\lambda ={}_{}{}^{}\delta _{}^{\lambda \mu }\widehat{O}_\mu ^r=(\widehat{O}_\lambda ^r)^+,\{\widehat{O}_\lambda ^r,\widehat{O}_\tau ^r^{}\}=\delta _{rr^{}}{}_{}{}^{}\delta _{\lambda \tau }^{}I_2.$$
(2.1.1)
The matrices $`\alpha _\lambda `$ satisfy the condition $`\{\alpha _\lambda ,\alpha _\tau \}={}_{}{}^{}\delta _{\lambda \tau }^{}I_2,`$ where $`\alpha ^\lambda ={}_{}{}^{}\delta _{}^{\lambda \mu }\alpha _\mu =(\alpha _\lambda )^+,`$ and $`I_2=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$. For example $`\alpha _1=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),\alpha _2=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).`$ The creation $`\widehat{O}_1^r`$ and annihilation $`\widehat{O}_2^r`$ operators are acting as follows:
$$\widehat{O}_1^r0>=O_1^r1>,\widehat{O}_2^r1>=O_2^r0>,$$
where $`0>0,0,\mathrm{}>`$ and $`1>0,\mathrm{},1,\mathrm{}>`$ are respectively the nonoccupied vacuum state and the one occupied state. Thus, $`\widehat{O}_1^r1>=0,\widehat{O}_2^r0>=`$0. A matrix realization of such states, for instance, can be: $`0>\chi _1=\left(\begin{array}{c}0\\ 1\end{array}\right),1>\chi _2=\left(\begin{array}{c}1\\ 0\end{array}\right).`$ Hence $`\chi _00>={\displaystyle \underset{r=1}{\overset{N}{}}}(\chi _1)_r`$ and $`\chi _r^{}1>=(\chi _2)_r^{}{\displaystyle \underset{rr^{}}{}}(\chi _1)_r`$.
Also, instead of ordinary basis vectors we introduce the operators $`\widehat{\sigma }_\alpha ^r\delta _{\alpha \beta \gamma }\sigma _\beta ^r\stackrel{~}{\sigma }_\gamma `$, where $`\stackrel{~}{\sigma }_\gamma `$ are Pauli’s matrices such that
$$<\sigma _\alpha ^r,\sigma _\beta ^r^{}>=\delta _{rr^{}}\delta _{\alpha \beta },\widehat{\sigma }_r^\alpha =\delta ^{\alpha \beta }\widehat{\sigma }_\beta ^r=(\widehat{\sigma }_\alpha ^r)^+=\widehat{\sigma }_\alpha ^r,\{\widehat{\sigma }_\alpha ^r,\widehat{\sigma }_\beta ^r^{}\}=2\delta _{rr^{}}\delta _{\alpha \beta }I_2.$$
(2.1.2)
The vacuum state $`0>\phi _{1(\alpha )}`$ and the one occupied state $`1_{(\alpha )}>\phi _{2(\alpha )}`$ read:
$`\phi _{1(\alpha )}\chi _1,\phi _{2(1)}=\left(\begin{array}{c}1\\ 0\end{array}\right),\phi _{2(2)}=\left(\begin{array}{c}i\\ 0\end{array}\right),\phi _{2(3)}=\left(\begin{array}{c}0\\ 1\end{array}\right),`$ thus,
$$\widehat{\sigma }_\alpha ^r\phi _{1(\alpha )}=\sigma _\alpha ^r\phi _{2(\alpha )}=(\sigma _\alpha ^r\stackrel{~}{\sigma }_\alpha )\phi _{1(\alpha )},\widehat{\sigma }_\alpha ^r\phi _{2(\alpha )}=\sigma _\alpha ^r\phi _{1(\alpha )}=(\sigma _\alpha ^r\stackrel{~}{\sigma }_\alpha )\phi _{2(\alpha )}.$$
Whence, the single eigenvalue $`(\sigma _\alpha ^r\stackrel{~}{\sigma }_\alpha )`$ associates with different $`\phi _{\lambda (\alpha )},`$ namely it is degenerate with degeneracy degree equal 2. Thus, among quantum numbers $`r`$ there is also the quantum number of the half integer spin $`\stackrel{}{\sigma }`$ $`(\sigma _3={\displaystyle \frac{1}{2}}s,s=\pm 1).`$ This consequently gives rise to the spins of particles. The one occupied state reads $`\phi _{r^{}(\alpha )}=(\phi _{2(\alpha )})_r^{}{\displaystyle \underset{rr^{}}{}}(\chi _1)_r.`$ Next, we introduce the operators
$$\widehat{\gamma }_{(\lambda ,\mu ,\alpha )}^r\widehat{O}_\lambda ^{r_1}\widehat{O}_\mu ^{r_2}\widehat{\sigma }_\alpha ^{r_3}$$
and the state vectors
$$\chi _{\lambda ,\mu ,\tau (\alpha )}\lambda ,\mu ,\tau (\alpha )>=\chi _\lambda \chi _\mu \phi _{\tau (\alpha )},$$
where $`\lambda ,\mu ,\tau ,\nu =1,2;\alpha ,\beta =1,2,3`$ and $`r(r_1,r_2,r_3)`$. Omitting two valuedness of state vectors we apply $`\lambda ,\tau ,\delta (\beta )>\lambda ,\tau >`$, and remember that always the summation must be extended over the double degeneracy of the spin states $`(s=\pm 1)`$. The explicit matrix elements of basis vectors read
$$<\lambda ,\mu \widehat{\gamma }{}_{(\tau ,\nu ,\alpha )}{}^{r}\tau ,\nu >={}_{}{}^{}\delta _{\lambda \tau }^{}{}_{}{}^{}\delta _{\mu \nu }^{}e_{(\tau ,\nu ,\alpha )}^r,<\tau ,\nu \widehat{\gamma }{}_{r}{}^{(\tau ,\nu ,\alpha )}\lambda ,\mu >={}_{}{}^{}\delta _{\lambda \tau }^{}{}_{}{}^{}\delta _{\mu \nu }^{}e_r^{(\tau ,\nu ,\alpha )}.$$
(2.1.3)
for given $`\lambda ,\mu .`$ The operators of occupation numbers
$$\underset{1}{\widehat{N}}{}_{\alpha \beta }{}^{rr^{}}=\widehat{\gamma }{}_{(1,1,\alpha )}{}^{r}\widehat{\gamma }{}_{(2,2,\beta )}{}^{r^{}},\underset{2}{\widehat{N}}{}_{\alpha \beta }{}^{rr^{}}=\widehat{\gamma }{}_{(2,1,\alpha )}{}^{r}\widehat{\gamma }_{(1,2,\beta )}^r^{}$$
(2.1.4)
have the expectation values implying Pauli’s exclusion principle
$$\begin{array}{cc}<2,2\underset{1}{\widehat{N}}{}_{\alpha \beta }{}^{rr^{}}2,2>=<1,2\underset{2}{\widehat{N}}{}_{\alpha \beta }{}^{rr^{}}1,2>=\delta _{rr^{}}\delta _{\alpha \beta },\hfill & \\ & \\ <1,1\underset{1}{\widehat{N}}{}_{\alpha \beta }{}^{rr^{}}1,1>=<2,1\underset{2}{\widehat{N}}{}_{\alpha \beta }{}^{rr^{}}2,1>=0.\hfill & \end{array}$$
(2.1.5)
The operators $`\{\widehat{\gamma }{}_{}{}^{r}\}`$ are the basis for tangent operator vectors $`\widehat{\mathrm{\Phi }}(\zeta )=\widehat{\gamma }{}_{}{}^{r}\mathrm{\Phi }_{r}^{}(\zeta )`$ of the 12 dimensional flat OM: $`\widehat{G}`$, where we introduce the vector function belonging to the ordinary class of functions of $`C^{\mathrm{}}`$ smoothness defined on the manifold $`G`$: $`\mathrm{\Phi }_r^{(\lambda ,\mu ,\alpha )}(\zeta )=\zeta ^{(\lambda ,\mu ,\alpha )}\mathrm{\Phi }_r^{\lambda ,\mu }(\zeta ),\zeta G`$. But, the operators $`\{\widehat{\gamma }{}_{r}{}^{}\}`$ is a dual basis for operator covectors $`\overline{\widehat{\mathrm{\Phi }}}(\zeta )=\widehat{\gamma }{}_{r}{}^{}\mathrm{\Phi }_{}^{r}(\zeta )`$, where $`\mathrm{\Phi }^r=\overline{\mathrm{\Phi }}_r`$ (charge conjugated). Hence
$$<\lambda ,\mu \widehat{\mathrm{\Phi }}(\zeta )\overline{\widehat{\mathrm{\Phi }}}(\zeta )\lambda ,\mu >={}_{}{}^{}\delta _{\lambda \tau }^{}{}_{}{}^{}\delta _{\mu \nu }^{}\mathrm{\Phi }_r^{(\tau ,\nu ,\alpha )}(\zeta )\mathrm{\Phi }_{(\tau ,\nu ,\alpha )}^r(\zeta ),$$
(2.1.6)
for given $`\lambda ,\mu .`$ Considering the state vectors $`\chi _\pm >`$ eq.(A.1.16) we get the matrix elements
$$\begin{array}{c}<\chi _+\widehat{\mathrm{\Phi }}(\zeta )\overline{\widehat{\mathrm{\Phi }}}(\zeta )\chi _+>\mathrm{\Phi }_+^2(\zeta )=\mathrm{\Phi }_r^{(\lambda ,1,\alpha )}(\zeta )\mathrm{\Phi }_{(\lambda ,1,\alpha )}^r(\zeta ),\hfill \\ \\ <\chi _{}\widehat{\mathrm{\Phi }}(\zeta )\overline{\widehat{\mathrm{\Phi }}}(\zeta )\chi _{}>\mathrm{\Phi }_{}^2(\zeta )=\mathrm{\Phi }_r^{(\lambda ,2,\alpha )}(\zeta )\mathrm{\Phi }_{(\lambda ,2,\alpha )}^r(\zeta ).\hfill \end{array}$$
(2.1.7)
The basis $`\{\widehat{\gamma }{}_{}{}^{r}\}`$ decomposes into $`\{\underset{i}{\widehat{\gamma }}{}_{}{}^{r}\}(i=\eta ,u),`$ where
$$\underset{i}{\widehat{\gamma }}{}_{(+\alpha )}{}^{r}=\frac{1}{\sqrt{2}}(\gamma _{(1,1\alpha )}^r+\epsilon _i\gamma _{(2,1\alpha )}^r),\underset{i}{\widehat{\gamma }}{}_{(\alpha )}{}^{r}=\frac{1}{\sqrt{2}}(\gamma _{(1,2\alpha )}^r+\epsilon _i\gamma _{(2,2\alpha )}^r).$$
The expansion of operator vectors $`\underset{i}{\widehat{\mathrm{\Psi }}}\underset{i}{\widehat{G}}`$ and operator covectors $`\overline{\underset{i}{\widehat{\mathrm{\Psi }}}}`$ are written $`\underset{i}{\widehat{\mathrm{\Psi }}}=\underset{i}{\widehat{\gamma }}{}_{}{}^{r}\underset{i}{\Psi }{}_{r}{}^{},\overline{\underset{i}{\widehat{\mathrm{\Psi }}}}=\underset{i}{\widehat{\gamma }}{}_{r}{}^{}\underset{i}{\Psi }{}_{}{}^{r},`$ where the following vector functions of $`C^{\mathrm{}}`$ smoothness are defined on the manifolds $`\underset{i}{𝐺}:`$
$$\underset{\eta }{\Psi }{}_{r}{}^{(\pm \alpha )}(\eta ,p_\eta )=\eta ^{(\pm \alpha )}\underset{\eta }{\Psi }{}_{r}{}^{\pm }(\eta ,p_\eta ),\underset{u}{\Psi }{}_{r}{}^{(\pm \alpha )}(u,p_u)=u^{(\pm \alpha )}\underset{u}{\Psi }{}_{r}{}^{\pm }(u,p_u).$$
(2.1.8)
Namely, the probability of finding the vector function in the state $`r`$ with given sixvector of coordinate ($`\eta `$ or $`u`$) and momentum ($`p_\eta `$ or $`p_u`$) is determined by the square of its state wave function $`\underset{\eta }{\Psi }{}_{r}{}^{\pm }(\eta ,p_\eta ),`$ or $`\underset{u}{\Psi }{}_{r}{}^{\pm }(u,p_u).`$ Due to the spin states, the $`\underset{i}{\Psi }_r^\pm `$ can be regarded as the Fermi field of the positive and negative frequencies $`\underset{i}{\Psi }{}_{r}{}^{\pm }\underset{i}{\Psi }{}_{\pm p}{}^{r}.`$
### 2.2 Realization of the Flat Manifold $`G`$
The bispinor $`\mathrm{\Psi }(\zeta )`$ defined on the manifold $`G=\underset{\eta }{𝐺}\underset{u}{𝐺}`$ can be written $`\mathrm{\Psi }(\zeta )=\underset{\eta }{\Psi }(\eta )\underset{u}{\Psi }(u)`$, where $`\underset{i}{\Psi }`$ is the bispinor defined on the manifold $`\underset{i}{𝐺}.`$ The free state of $`i`$-type fermion with definite values of momentum $`p_i`$ and spin projection $`s`$ is described by plane waves (App.A). The relations of orthogonality and completeness hold for the spinors. Considering also the solutions of negative frequencies, we make use of localized wave packets constructed by means of superposition of plane wave solutions furnished by creation and annihilation operators in agreement with Pauli’s principle
$$\underset{i}{\widehat{\mathrm{\Psi }}}=\underset{\pm s}{}\frac{d^3p_i}{(2\pi )^{3/2}}\left(\underset{i}{\widehat{\gamma }}_{(+\alpha )}\underset{i}{\Psi }^{(+\alpha )}+\underset{i}{\widehat{\gamma }}_{(\alpha )}\underset{i}{\Psi }^{(\alpha )}\right),$$
etc, where the summation is extended over all dummy indices. In such a manner we can treat as well the wave packets of operator vector fields $`\widehat{\mathrm{\Phi }}(\zeta )`$. While the matrix element of the anticommutator of expansion coefficients reads
$$<\chi _{}\{\underset{i}{\widehat{\gamma }}{}_{}{}^{(+\alpha )}(p_i,s),\underset{j}{\widehat{\gamma }}{}_{(+\beta )}{}^{}(p_j^{},s^{})\}\chi _{}>=\epsilon _i\delta _{ij}\delta _{ss^{}}\delta _{\alpha \beta }\delta ^{(3)}(\stackrel{}{p}_i\stackrel{}{p^{}}_i).$$
(2.2.1)
In the aftermath, we get the most important relation
$$\begin{array}{c}\underset{\lambda =\pm }{}<\chi _\lambda \widehat{\mathrm{\Phi }}(\zeta )\overline{\widehat{\mathrm{\Phi }}}(\zeta )\chi _\lambda >=\underset{\lambda =\pm }{}<\chi _\lambda \overline{\widehat{\mathrm{\Phi }}}(\zeta )\widehat{\mathrm{\Phi }}(\zeta )\chi _\lambda >=\hfill \\ \\ =i\zeta ^2\underset{\zeta }{𝐺}(0)=i\left(\eta ^2\underset{\eta }{𝐺}(0)u^2\underset{u}{𝐺}(0)\right),\hfill \end{array}$$
(2.2.2)
where $`\underset{i}{𝐺}(0)\underset{ii^{}}{lim}\underset{i}{𝐺}(ii^{}),(i=\zeta ,\eta ,u,)`$, etc., the Green’s function $`\underset{i}{𝐺}(ii^{})=(i\underset{i}{\widehat{}}+m)\underset{i}{\Delta }(ii^{})`$ is provided by the usual invariant singular functions $`\underset{i}{\Delta }(ii^{}).`$
Realization of the flat manifold $`G`$ ensued from the constraint imposed upon the matrix element eq.(2.2.2) that, as the geometric object, it is required to be finite
$$\underset{\lambda =\pm }{}<\chi _\lambda \widehat{\mathrm{\Phi }}(\zeta )\overline{\widehat{\mathrm{\Phi }}}(\zeta )\chi _\lambda ><\mathrm{},$$
(2.2.3)
which gives rise to
$$\zeta ^2\underset{\zeta }{𝐺}{}_{F}{}^{}(0)<\mathrm{},$$
(2.2.4)
and
$$\begin{array}{c}\underset{\zeta }{𝐺}{}_{F}{}^{}(0)=\underset{\eta }{𝐺}{}_{F}{}^{}(0)=\underset{u}{𝐺}{}_{F}{}^{}(0)=\hfill \\ =\underset{uu^{}}{lim}[i\underset{\stackrel{}{p}_u}{}\underset{u}{\Psi }{}_{p_u}{}^{}(u)\underset{u}{\overline{\mathrm{\Psi }}}{}_{p_u}{}^{}(u^{})\theta (u_0u_0^{})+i\underset{\stackrel{}{p}_u}{}\underset{u}{\overline{\mathrm{\Psi }}}{}_{p_u}{}^{}(u^{})\underset{u}{\Psi }{}_{p_u}{}^{}(u)\theta (u_0^{}u_0)],\hfill \end{array}$$
(2.2.5)
where the $`\underset{\zeta }{𝐺}{}_{F}{}^{},\underset{\eta }{𝐺}_F`$ and $`\underset{u}{𝐺}_F`$ are causal Green’s functions characterized by the boundary condition that only positive frequency occur for $`\eta _0>0(u_0>0)`$, only negative for $`\eta _0<0(u_0<0)`$. Here $`\eta _0=\stackrel{}{\eta }_0`$, $`\eta _{0\alpha }={\displaystyle \frac{1}{\sqrt{2}}}(\eta _{(+\alpha )}+\eta _{(\alpha )})`$ and the same holds for $`u_0`$. Satisfying the condition eq.(2.2.4) the length of each vector $`\zeta =e\zeta G`$ (see eq.(2.2.2)) compulsory must be equaled zero
$$\zeta ^2=\eta ^2u^2=0.$$
(2.2.6)
Thus, the requirement eq.(2.2.3) provided by eq.(2.2.5) yields the realization of the flat manifold $`G`$, which subsequently leads to Minkowski flat space $`M_4`$ (subsec.2.1) where, according to eq.(2.2.6), the relativity principle holds
$$d\eta ^2|_{64}ds^2=du^2=inv.$$
## 3 Mathematical Background
$``$ Field Aspect
The quantum field theory of the OM is equivalent to configuration space wave mechanics employing the antisymmetric state functions incorporated with geometric properties of corresponding objects (see App.A or . Therein, by applying the algebraic approach we reach to rigorous definition of the OM: $`\widehat{G}`$, construct the explicit forms of wave state functions and calculate the matrix elements of field operators. While, the $`\widehat{G}`$ reads
$$\widehat{G}=\underset{n=0}{\overset{\mathrm{}}{}}\widehat{G}^{(n)}=\underset{n=0}{\overset{\mathrm{}}{}}\left(\widehat{𝒰}^{(n)}\overline{}^{(n)}\right),$$
(3.1)
where $`\widehat{𝒰}_{(r_1,\mathrm{},r_n)}^{(n)}=\widehat{𝒰}_{r_1}^{(1)},\mathrm{}\widehat{𝒰}_{r_n}^{(1)}`$ is the open neighbourhood of the n-points $`\widehat{\zeta }_{r_i}`$ of the OM, $`\overline{}_{(r_1,\mathrm{},r_n)}^{(n)}=_{r_1}^{(1)}\mathrm{}_{r_n}^{(1)}`$ is the Hilbert space for description of n particle system. Meanwhile, one has to modify the basis operators (the creation $`\widehat{\gamma }_r`$ and annihilation $`\widehat{\gamma }^r`$ operators) in order to provide an anticommutation in arbitrary states. For example, acting on free state $`0>_{r_i}`$ the creation operator $`\widehat{\gamma }_{r_i}`$ now yields the one occupied state $`1>_{r_i}`$ with the phase $`{}_{}{}^{}+_{}^{}`$ or $`{}_{}{}^{}_{}^{}`$ depending of parity of the number of quanta in the states $`r<r_i`$. Modified operators satisfy the same anticommutation relations of the basis operators (eq.(A.1.12)). Defining the secondary quantized form of one particle observable $`A`$ on the $``$ we consider a set of identical samples $`\widehat{}_i`$ of one particle space $`^{(1)}`$ and operators $`A_i`$ acting on them. The vacuum state given in eq.(A.1.16) satisfies the normalization condition. The state vectors eq.(A.1.17) are the eigenfunctions of modified operators. They form a whole set of orthogonal vectors. Considering an arbitrary superposition of state vectors we get a whole set of explicit forms of the matrix elements of operator vector and covector fields (App.A).
$``$ Differential Geometric Aspect
For illustrative purposes here we consider a few examples from the differential geometric aspect of the OM by referring to the Appendix A for more details,.
The operators $`\{\widehat{\gamma }{}_{}{}^{r}\}`$ are the basis for all operator vectors of tangent section $`\widehat{𝐓}_{\mathrm{\Phi }_p}`$ of principle bundle with the base $`\widehat{G}`$ at the point $`𝚽_p=𝚽(\zeta (t))|_{t=0}\widehat{G}`$. The smooth field of tangent operator vector $`\widehat{𝐀}(𝚽(\zeta ))`$ is a class of equivalence of the curves $`𝐟(𝚽(\zeta ))`$, $`𝐟(𝚽(\zeta (0)))=𝚽_p`$. While, the operator differential $`\widehat{d}A_p^t`$ of the flux $`A_p^t:\widehat{G}\widehat{G}`$ at the point $`𝚽_p`$ with the velocity fields $`\widehat{𝐀}(𝚽(\zeta ))`$ is defined by one parameter group of operator diffeomorphisms given for the curve $`𝚽(\zeta (t)):R^1\widehat{G}`$. Provided one has $`𝚽(\zeta (0))=𝚽_p`$ and $`\widehat{\dot{𝚽}}(\zeta (0))=\widehat{𝐀}_p`$
$$\widehat{d}A_p^t(𝐀)=\frac{\widehat{d}}{dt}|_{t=0}A^t(𝚽(\zeta (t)))=\widehat{𝐀}(𝚽(\zeta ))=\widehat{\gamma }{}_{}{}^{r}A_{p}^{},$$
(3.2)
where the $`\{A_p\}`$ are the components of $`\widehat{𝐀}`$ in the basis $`\{\widehat{\gamma }{}_{}{}^{r}\}`$. According to eq.(3.2), in holonomic coordinate basis $`\widehat{\gamma }{}_{}{}^{r}(\widehat{}/\mathrm{\Phi }_r(\zeta (t)))_p`$ one gets $`A_p={\displaystyle \frac{\mathrm{\Phi }_r}{\zeta _r}}{\displaystyle \frac{d\zeta _r}{dt}}|_p.`$ Hence, for any function $`f:𝐑^n𝐑^n`$ of the ordinary class of functions of $`C^{\mathrm{}}`$ smoothness on $`\widehat{G}`$ one may define an operator differential
$$<\widehat{d}f,\widehat{𝐀}>=\widehat{(Af)},$$
by means of smooth reflection
$$\widehat{d}f:\widehat{𝐓}\left(\widehat{G}\right)\widehat{R}\left(\widehat{𝐓}\left(\widehat{G}\right)=\underset{\mathrm{\Phi }_p}{}\widehat{𝐓}_{\mathrm{\Phi }_p}\right),$$
where (see eq.(A.1.25))
$$<\chi \widehat{d}f,\widehat{𝐀}\chi ^0>=\underset{\lambda ,\mu =1}{\overset{2}{}}\underset{r^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}\widehat{c}^{}(r^{\lambda \mu })<df,𝐀>_{r^{\lambda \mu }}=\underset{\lambda ,\mu =1}{\overset{2}{}}\underset{r^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}\widehat{c}^{}(r^{\lambda \mu })(𝐀f)_{r^{\lambda \mu }}.$$
(3.3)
In coordinate basis $`<d\mathrm{\Phi }^{\widehat{ı}},\widehat{}/\mathrm{\Phi }^j>={\displaystyle \frac{\mathrm{\Phi }^{\widehat{ı}}}{\mathrm{\Phi }^j}}=\delta _j^{\widehat{ı}},`$ provided by $`d\mathrm{\Phi }^{\widehat{ı}}\widehat{d}\mathrm{\Phi }^i`$ and
$`<\chi \widehat{\delta }{}_{j}{}^{ı}\chi ^0>={\displaystyle \underset{\lambda ,\mu =1}{\overset{2}{}}}{\displaystyle \underset{r^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}}\widehat{c}^{}(r^{\lambda \mu })\delta _j^i,`$ where the $`i`$ and $`j`$ stand for a set of $`(\lambda _i,\mu _i,\alpha _i)`$.
The operator tensor $`\widehat{𝐓}`$ of $`\widehat{(n,0)}`$-type at the point $`𝚽_p`$ is a linear function of the space $`\widehat{𝐓}_0^n=\underset{n}{\underset{}{\widehat{𝐓}_{\mathrm{\Phi }_p}\mathrm{}\widehat{𝐓}_{\mathrm{\Phi }_p}}},`$ where the $``$ denotes the tensor product. It enables a correspondence between the element $`(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n)`$ of $`\widehat{𝐓}_0^n`$ and the number $`T(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n)`$ furnished by linearity. Constructing matrix elements of operator tensors of $`\widehat{G}`$ one produces the Cartan’s exterior forms (A.1.28). Whence, the matrix elements of symmetric operator tensors equal zero. The differential operator $`n`$ form $`\widehat{\omega }^n|_{\mathrm{\Phi }_p}`$ at the point $`𝚽_p\widehat{G}`$ can be defined as the exterior operator $`n`$ form on tangent operator space $`\widehat{𝐓}_{\mathrm{\Phi }_p}`$ of tangent operator vectors $`\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n`$. That is, if the $`\widehat{𝐓}_{\mathrm{\Phi }_p}^{}\left(\widehat{G}\right)`$ means the exterior algebra on $`\widehat{𝐓}_{\mathrm{\Phi }_p}^{}\left(\widehat{G}\right)`$, then operator $`n`$ form $`\widehat{\omega }^n|_{\mathrm{\Phi }_p}`$ is an element of $`n`$-th degree out of $`\widehat{𝐓}_{\mathrm{\Phi }_p}^{}`$ depending of the point $`𝚽_p\widehat{G}`$. Hence $`\widehat{\omega }^n={\displaystyle \underset{\mathrm{\Phi }_p}{}}\widehat{\omega }^n|_{\mathrm{\Phi }_p}`$. Any differential operator $`n`$ form of dual operator space $`\underset{n}{\underset{}{\widehat{𝐓}_{\mathrm{\Phi }_p}^{}\mathrm{}\widehat{𝐓}_{\mathrm{\Phi }_p}^{}}}`$ may be written $`\widehat{\omega }^n={\displaystyle \underset{i_1<\mathrm{}<i_n}{}}\alpha _{i_1\mathrm{}i_n}(\mathrm{\Phi })d\mathrm{\Phi }^{\widehat{ı}_1}\mathrm{}d\mathrm{\Phi }^{\widehat{ı}_n},`$ provided by the smooth differentiable functions $`\alpha _{i_1\mathrm{}i_n}(\mathrm{\Phi })C^{\mathrm{}}`$ and basis $`d\mathrm{\Phi }^{\widehat{ı}_1}\mathrm{}d\mathrm{\Phi }^{\widehat{ı}_n}={\displaystyle \underset{\sigma S_n}{}}sgn(\sigma )\gamma ^{\sigma (\widehat{ı}_1}\mathrm{}\gamma ^{\widehat{ı}_n)}.`$ The matrix elements of some of the geometric objects of the $`\widehat{G}`$ are given in the App.A.
## 4 Beyond the Geometry and Fields
To facilitate the physical picture and provide sufficient background it seems worth to bring few formal matters in concise form which one will have to know in order to understand the general structure of our approach without undue hardship. Here we only outline briefly the relevant steps. In the mean time we refer to for more detailed justification of some of the procedures and complete exposition. Before proceeding further, it is profitable to define the pulsating gauge functions and fields denoted by wiggles as follows:
$``$ The function $`\stackrel{~}{W}(x)`$ defined on the space $`M`$ ($`xM`$ ) and being an invariant with respect to the coordinate transformations is called the pulsating gauge function if it undergoes local gauge transformations
$$\stackrel{~}{W}^{}(x)=U(x)\stackrel{~}{W}(x).$$
(4.1)
Here $`U(x)`$ is the element of some simple Lie group $`G`$ the generators of which imply the algebra $`[F^a,F^b]=iC^{abc}F^c`$, where $`C^{abc}`$ are wholly antisymmetric structure constants.
$``$ A smooth function $`\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{W}(x))`$ belonged to some representation of the group $`G`$, where the generators are presented by the matrices $`T^a`$, is called the pulsating field if under the transformation eq.(4.1) it transforms
$$\stackrel{~}{\mathrm{\Phi }}^{}\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{W}^{}(x))=U(x)\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{W}(x)).$$
(4.2)
Let $`L_0(\mathrm{\Phi },\mathrm{\Phi })`$ is the invariant Lagrangian of free field $`\mathrm{\Phi }`$ defined on $`M`$. Then, a simple gauge invariant Lagrangian of the pulsating field $`\stackrel{~}{\mathrm{\Phi }}`$ can be written
$$L=\stackrel{~}{W}^+\stackrel{~}{W}L_0(\mathrm{\Phi },\mathrm{\Phi }),$$
(4.3)
which reduces to
$$LL(\stackrel{~}{\mathrm{\Phi }},\stackrel{~}{D\mathrm{\Phi }})=L_0(\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{W}(x)),D\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{W}(x))).$$
(4.4)
Here we have noticed that due to eq.(4.1) and eq.(4.2) $`\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{W})=\stackrel{~}{W}\mathrm{\Phi }`$, and introduced the covariant derivative $`\stackrel{~}{D\mathrm{\Phi }}D\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{W})=\stackrel{~}{W}\mathrm{\Phi }.`$ Whence
$$D=igT^aW^a,T^aW^a=\frac{i}{g}\mathrm{ln}\stackrel{~}{W},D\stackrel{~}{W}=\left(D\stackrel{~}{W}\right)^+=0,$$
(4.5)
where $`W^a`$ is the gauge field, g is the coupling constant. Hence, all the conventional matter fields interacting by gauge fields are the pulsating fields.
### 4.1 The Regular Primordial Structures
In we have chosen a simple setting and considered the primordial structures designed to possess certain physical properties satisfying the stated general rules. These structures are the substance out of which the geometry and particles are made. We distinguish the “$`\eta `$\- and $`u`$-types primordial structures” involved in the linkage establishing processes occurred between the structures of different types. Let us recall that the $`\eta `$-type structure may accept the linkage only from $`u`$-type structure, which is described by the link function $`\underset{\eta }{\Psi }(s)`$ belonged to the ordinary class of functions of $`C^{\mathrm{}}`$ smoothness, where $`s\eta =\underset{\eta }{𝑒}{}_{(\lambda \alpha )}{}^{}\eta _{}^{(\lambda \alpha )},(\lambda =\pm ;\alpha =1,2,3,`$ see subsec.2.1), $`\eta `$ is the link coordinate. Respectively the $`u`$-type structure may accept the linkage only from $`\eta `$-type structure described by the link function $`\underset{u}{\Psi }(s)`$ (u-channel), where $`su=\underset{u}{𝑒}u`$. We assume that $`s`$ is the pulsating gauge function associated with the Abelian local gauge group $`U(1)`$ and $`\mathrm{\Psi }(s)`$ is the pulsating field (the wiggles are left implicit). Thus, under local gauge transformations
$$s^{}=e^{i\alpha }s,\alpha 0,$$
the link function $`\mathrm{\Psi }(s)`$ transforms
$$\mathrm{\Psi }(s^{})=e^{i\alpha }\mathrm{\Psi }(s),$$
and the Lagrangian eq.(4.3) is invariant under gauge transformations. It includes the covariant derivative $`D(s)=+igb(s)`$ and gauge field $`b(s)={\displaystyle \frac{i}{g}}\mathrm{ln}s`$ undergone gauge transformations $`b(s^{})=b(s)+{\displaystyle \frac{1}{g}}\alpha .`$ Then, $`\underset{i}{\Psi }(s)=s\underset{i}{\Psi }=\underset{i}{𝑒}{}_{(\lambda \alpha )}{}^{}\underset{i}{\Psi }_{}^{(\lambda \alpha )}`$ ($`i=\eta ,u`$), where the eq.(2.1.8) holds
$$\underset{\eta }{\Psi }^{(\pm \alpha )}(\eta ,p_\eta )=\eta ^{(\pm \alpha )}\underset{\eta }{\Psi }^\pm (\eta ,p_\eta ),\underset{u}{\Psi }^{(\pm \alpha )}(u,p_u)=u^{(\pm \alpha )}\underset{u}{\Psi }^\pm (u,p_u),$$
(4.1.1)
a bispinor $`\underset{i}{\Psi }^\pm `$ is the invariant state wave function of positive or negative frequencies, $`p_i`$ is the corresponding link momentum. Thus, a primordial structure can be considered as a fermion found in external gauge field $`b(s).`$
The simplest system made of two structures of different types becomes stable only due to the stable linkage
$$\left|\underset{\eta }{𝑝}\right|=(\underset{\eta }{𝑝}^{(\lambda \alpha )},\underset{\eta }{𝑝}{}_{(\lambda \alpha )}{}^{})^{1/2}=\left|\underset{u}{𝑝}\right|=(\underset{u}{𝑝}^{(\lambda \alpha )},\underset{u}{𝑝}{}_{(\lambda \alpha )}{}^{})^{1/2}.$$
(4.1.2)
Otherwise they are unstable. There is not any restriction on the number of primordial structures of both types getting into the link establishing processes simultaneously. In the stable system the link stability condition must be held for each linkage separately.
The persistent processes of creation and annihilation of the primordial structures occur in different states $`s,s^{},s^{\prime \prime },\mathrm{}`$ The ”creation” of structure in the given state $`(s)`$ is due to its transition to this state from other states $`(s^{},s^{\prime \prime },\mathrm{})`$, while the ”annihilation” means a vice versa. Satisfying eq.(4.1.2) the primordial structures from the arbitrary states may establish a stable linkage. Among the states $`(s,s^{},s^{\prime \prime },\mathrm{})`$ there is a lowest one ($`s_0`$), in which all structures are regular. That is, they are in free (pure) state and described by the plane wave functions $`\underset{\eta }{\Psi }^\pm (\eta _f,p_\eta )`$ or $`\underset{u}{\Psi }^\pm (u_f,p_u)`$ defined respectively on flat manifolds $`\underset{\eta }{𝐺}`$ and $`\underset{u}{𝐺}`$. The index (f) specifies the points of corresponding flat manifolds $`\eta _f\underset{\eta }{𝐺}`$, $`u_f\underset{u}{𝐺}`$. For example, in accordance with subsec.2.2, the equation of regular structure $`\mathrm{\Psi }(s_+)(s=s_++s_{})`$ reads
$$\left[i\gamma _f(+igb(s_+))m\right]\mathrm{\Psi }(s_+)=0,$$
the matrices $`\gamma _f`$ are given in eq.(B.1.4). Whence the equation of plane wave function $`\mathrm{\Psi }_p^+`$ of positive frequencies stems
$$(i\gamma _fm)\mathrm{\Psi }_p^+=0.$$
The processes of creation and annihilation of regular structures in lowest state are described by the OM formalism given in the previous sections.
### 4.2 The Distorted Primordial Structures
In all the higher states the primordial structures are distorted ones (interaction states) and described by distorted link functions defined on distorted manifolds $`\stackrel{~}{\underset{\eta }{𝐺}}`$ and $`\stackrel{~}{\underset{u}{𝐺}}`$. A distortion $`G\stackrel{~}{G}`$ with hidden Abelian local group $`G=U^{loc}(1)=SO^{loc}(2)`$ and one dimensional trivial algebra $`\widehat{g}=R^1`$ has studied in . It involves a drastic revision of a role of local internal symmetries in the concept of curved geometry. Under the reflection of fields and their dynamics from Minkowski space to Riemannian a standard gauge principle of local internal symmetries is generalized. The gravitation gauge group is proposed, which is generated by hidden local internal symmetry. This suggests an opportunity for the unification of all interactions on an equal footing. Our scheme is implemented as follows: Considering the principle bundle $`p:EG`$ the basis $`e^f`$ is transformed $`e=De^f,`$ under massless gauge distortion field $`a_f`$ associated with $`U^{loc}(1)`$. The matrix $`D`$ is in the form $`D=CR`$, where the distortion transformations $`O_{(\lambda \alpha )}=C_{(\lambda \alpha )}^\tau O_\tau `$ and $`\sigma _{(\lambda \alpha )}=R_{(\lambda \alpha )}^\beta \sigma _\beta `$ are defined. Here $`C_{(\lambda \alpha )}^\tau =\delta _\lambda ^\tau +\kappa a_{(\lambda \alpha )}{}_{}{}^{}\delta _{\lambda }^{\tau }`$, but $`R`$ is a matrix of the group $`SO(3)`$ of ordinary rotations of the planes involving two arbitrary basis vectors of the spaces $`R_\pm ^3`$ around the orthogonal third axes. The rotation angles are determined from the constraint imposed upon distortion transformations that a sum of distorted parts of corresponding basis vectors $`O_\lambda `$ and $`\sigma _\beta `$ should be zero for given $`\lambda `$ (App.B). Whence $`\mathrm{tan}\theta _{(\lambda \alpha )}=\kappa a_{(\lambda \alpha )}`$, where $`\theta _{(\lambda \alpha )}`$ is the particular rotation around the axis $`\sigma _\alpha `$. Next we construct the diffeomorphism $`G\stackrel{~}{G}`$ and introduce the invariant action of the fields. The passage from six dimensional curved manifold $`\stackrel{~}{G}`$ to four dimensional Riemannian geometry $`R^4`$ is straightforward (subsec.2.1). Given a diffeomorphism $`u(u_f):\underset{u}{𝐺}\stackrel{~}{\underset{u}{𝐺}}`$ we consider the reflection of the Fermi fields and their dynamics from the flat manifold $`\underset{u}{𝐺}`$ to distorted one $`\stackrel{~}{\underset{u}{𝐺}}`$, and vice versa (App.B). In the aftermath, the relation between the wave functions of distorted and regular structures reads
$$\underset{u}{\Psi }^\lambda (\theta _{+k})=f_{(+)}(\theta _{+k})\underset{u}{\Psi }^\lambda ,\underset{u}{\Psi }{}_{\lambda }{}^{}(\theta _k)=\underset{u}{\Psi }{}_{\lambda }{}^{}f_{()}^{}(\theta _k).$$
(4.2.1)
The $`\underset{u}{\Psi }^\lambda (\underset{u}{\Psi }{}_{\lambda }{}^{})`$ is the plane wave function of regular ordinary structure (antistructure) and
$$f_{(+)}(\theta _{+k})=e^{\chi _R(\theta _{+k})i\chi _J(\theta _{+k})},f_{()}(\theta _k)=f_{(+)}^{}(\theta _{+k})|_{\theta _{+k}=\theta _k},$$
(4.2.2)
where the $`\chi _R`$ and $`\chi _J`$ are given in Appendix C.
Next, we supplement the previous assumptions made in sec. 4 by a new one that now the $`\eta `$-type (fundamental) regular structure can not directly form a stable system with the regular $`u`$-type (ordinary) structures. Instead of it the $`\eta `$-type regular structure forms a stable system with the infinite number of distorted ordinary structures, where the link stability condition held for each linkage separately. Such structures take part in realization of the flat manifold $`G`$ (subsec.2.2). The laws regarding to this change apply in use of functions of distorted ordinary structures
$$\underset{u}{\Psi }^{(\lambda \alpha )}(\theta _+)=u^{(\lambda \alpha )}\underset{u}{\Psi }^\lambda (\theta _+),\underset{u}{\Psi }{}_{(\lambda \alpha )}{}^{}(\theta _{})=u_{(\lambda \alpha )}\underset{u}{\Psi }{}_{\lambda }{}^{}(\theta _{}),$$
(4.2.3)
where $`u\stackrel{~}{\underset{u}{𝐺}}`$. For our immediate purposes we employ the wave packets constructed by superposition of these functions furnished by generalized operators of creation and annihilation as the expansion coefficients
$$\begin{array}{c}\underset{u}{\widehat{\mathrm{\Psi }}}(\theta _+)=\underset{\pm s}{}\frac{d^3p_u}{(2\pi )^{3/2}}(\underset{u}{\widehat{\gamma }}{}_{(+\alpha )}{}^{k}\underset{u}{\Psi }{}_{}{}^{(+\alpha )}(\theta _{+k})+\underset{u}{\widehat{\gamma }}{}_{(\alpha )}{}^{k}\underset{u}{\Psi }{}_{}{}^{(\alpha )}(\theta _{+k})),\hfill \\ \\ \overline{\underset{u}{\widehat{\mathrm{\Psi }}}}(\theta _{})=\underset{\pm s}{}\frac{d^3p_u}{(2\pi )^{3/2}}(\underset{u}{\widehat{\gamma }}{}_{k}{}^{(+\alpha )}\underset{u}{\Psi }{}_{(+\alpha )}{}^{}(\theta _k)+\underset{u}{\widehat{\gamma }}{}_{k}{}^{(\alpha )}\underset{u}{\Psi }{}_{(\alpha )}{}^{}(\theta _k)),\hfill \end{array}$$
(4.2.4)
where as usual the summation is extended over all dummy indices. The matrix element of anticommutator of generalized expansion coefficients reads
$$<\chi _{}\{\underset{u}{\widehat{\gamma }}{}_{k}{}^{(+\alpha )}(p,s),\underset{u}{\widehat{\gamma }}{}_{(+\beta )}{}^{k^{}}(p^{},s^{})\}\chi _{}>=\delta _{ss^{}}\delta _{\alpha \beta }\delta _{kk^{}}\delta ^3(\stackrel{}{p}\stackrel{}{p^{}}).$$
(4.2.5)
The wave packets eq.(4.2.4) yield the causal Green’s function $`\underset{u}{𝐺}_F^\theta (\theta _+\theta _{})`$ of distorted ordinary structure. Geometry realization requirement (eq.(2.2.5)) now should be satisfied for each ordinary structure in terms of
$$\underset{u}{𝐺}{}_{F}{}^{\theta }(0)=\underset{\theta _+\theta _{}}{lim}\underset{u}{𝐺}{}_{F}{}^{\theta }(\theta _+\theta _{})=\underset{\eta }{𝐺}{}_{F}{}^{}(0)=\underset{\eta _f^{}\eta _f}{lim}\underset{\eta }{𝐺}{}_{F}{}^{}(\eta _f^{}\eta _f).$$
(4.2.6)
They are valid if following relations hold for each distorted ordinary structure:
$$\underset{k}{}\underset{u}{\Psi }(\theta _{+k})\overline{\underset{u}{\Psi }}(\theta _k)=\underset{k}{}\underset{u}{\Psi }^{}(\theta _{+k}^{})\overline{\underset{u}{\Psi }}^{}(\theta _k^{})=\mathrm{}=inv.$$
(4.2.7)
Namely, the distorted ordinary structures have met in the permissible combinations to realize the geometry in the stable system. Below, in simplified schematic way we exploit the background of the known colour confinement and gauge principles. This scheme still should be considered as a preliminary one, which will be further elaborated in the sec.5.
### 4.3 “Quarks” and “Colour” Confinement
At the very first to avoid irrelevant complications, here, for illustrative purposes, we will attempt to introduce temporarily skeletonized “quark” and “antiquark” fields emerged in confined phase in the simplified geometry with the one-u channel given in the previous subsections. The complete picture of such a dynamics is beyond the scope of this subsection, but some relevant discussions on this subject will also be presented in the subsec.5.3. We may think of the function $`\underset{u}{\Psi }{}_{}{}^{\lambda }(\theta _{+k})`$ at fixed $`(k)`$ as being the $`u`$-component of bispinor field of “quark” $`q_k`$, and of $`\overline{\underset{u}{\Psi }}{}_{\lambda }{}^{}(\theta _k)`$ \- the $`u`$-component of conjugated bispinor field of “antiquark” $`\overline{q}_k`$. The index $`(k)`$ refers to colour degree of freedom in the case of rotations through the angles $`\theta _{+k}`$ and anticolour degree of freedom in the case of $`\theta _k`$. The $`\eta `$-components of quark fields are plane waves. In both cases of local and global rotations we respectively distinguish two types of quarks: local $`q_k`$ and global $`q_k^c`$. Hence, the quark is a fermion with the half integer spin and certain colour degree of freedom. There are exactly three colours. The rotation through the angle $`\theta _{+k}`$ yields a total quark field defined on the flat manifold $`G=\underset{\eta }{𝐺}\underset{u}{𝐺}`$
$$q_k(\theta )=\mathrm{\Psi }(\theta _{+k})=\underset{\eta }{\Psi }^0\underset{u}{\Psi }(\theta _{+k})$$
(4.3.1)
where $`\underset{\eta }{\Psi }^0`$ is a plane wave defined on $`\underset{\eta }{𝐺}`$. According to eq.(4.2.1), one gets
$$q_k(\theta )=\underset{\eta }{\Psi }^0\underset{u}{𝑞}_k(\theta )=\underset{\eta }{𝑞}_k(\theta )\underset{u}{\Psi }^0,\underset{\eta }{𝑞}_k(\theta )f_{(+)}(\theta _{+k})\underset{\eta }{\Psi }^0,$$
(4.3.2)
where $`\underset{u}{\Psi }^0`$ is a plane wave, $`\underset{u}{𝑞}{}_{k}{}^{}(\theta )`$ and $`\underset{\eta }{𝑞}{}_{k}{}^{}(\theta )`$ may be considered as the quark fields with the same quantum numbers defined respectively on flat manifolds $`\underset{u}{𝐺}`$ and $`\underset{\eta }{𝐺}`$. By making use of the rules stated in subsec.2.1 one may readily return to Minkowski space $`\underset{\eta }{𝐺}M_4`$. In the sequel, the quark field defined on $`M_4`$ will be ensued $`\underset{\eta }{𝑞}{}_{k}{}^{}(\theta )q_k(x)`$, $`xM_4`$. Due to eq.(4.2.7) and eq.(4.3.1) they imply
$$\underset{k}{}q_k\overline{q}_k=\underset{k}{}q_{}^{}{}_{k}{}^{}\overline{q^{}}_k=\mathrm{}=inv,$$
(4.3.3)
namely
$$\underset{k}{}f_{(+)}(\theta _{+k})f_{()}(\theta _k)=\underset{k}{}f_{(+)}^{}(\theta _{+k}^{})f_{()}^{}(\theta _k^{})=\mathrm{}=inv.$$
(4.3.4)
The eq.(4.3.3) utilizes the idea of colour (quark) confinement principle: the quarks emerge in the geometry only in special combinations of colour singlets. Only two colour singlets are available (see below)
$$(q\overline{q})=\frac{1}{\sqrt{3}}\delta _{kk^{}}\widehat{q}_k\overline{\widehat{q}}_k^{}=inv,(qqq)=\frac{1}{\sqrt{6}}\epsilon _{klm}\widehat{q}_k\widehat{q}_l\widehat{q}_m=inv.$$
(4.3.5)
These results will be generalized in the next section where the physically more realistic scheme of the MW geometry with the multi u-channel should be subject for discussion.
### 4.4 Gauge Principle; Internal Symmetries
Following , the principle of identity holds for ordinary regular structures, namely each regular structure in the lowest state can be regarded as a result of transition from an arbitrary state, in which they assumed to be distorted. This is stated in terms of link-functions below
$$\underset{u}{\Psi }{}_{\lambda }{}^{}=f_{(+)}^1(\theta _{+k})\underset{u}{\Psi }{}_{\lambda }{}^{}(\theta _{+k})=f_{(+)}^1(\theta _{+l}^{})\underset{u}{\Psi }{}_{\lambda }{}^{}(\theta _{+l}^{})=\mathrm{}.$$
(4.4.1)
Hence, the following transformations may be implemented upon distorted ordinary structures occurred in the stable system:
$$\begin{array}{c}\underset{u}{\Psi }^\lambda (\theta _{+l}^{})=f_{lk}^{(+)}\underset{u}{\Psi }^\lambda (\theta _{+k})=f(\theta _{+l}^{},\theta _{+k})\underset{u}{\Psi }^\lambda (\theta _{+k}),\hfill \\ \\ \underset{u}{\Psi }{}_{\lambda }{}^{}(\theta _l^{})=\underset{u}{\Psi }{}_{\lambda }{}^{}(\theta _k)f_{kl}^{()}=\underset{u}{\Psi }{}_{\lambda }{}^{}(\theta _k)f^{}(\theta _l^{},\theta _k)|_{\begin{array}{c}\theta _l^{}=\theta _{+l}^{}\hfill \\ \theta _k=\theta _{+k},\hfill \end{array}}\hfill \end{array}$$
(4.4.2)
provided by
$$f_{lk}^{(+)}=\mathrm{exp}\{\chi _{lk}^Ri\chi _{lk}^J\},f_{kl}^{()}=(f_{lk}^{(+)})^{}|_{\begin{array}{c}\theta _l^{}=\theta _{+l}^{}\hfill \\ \theta _k=\theta _{+k},\hfill \end{array}},$$
(4.4.3)
$$\chi _{lk}^R=\chi _R(\theta _{+l}^{})\chi _R(\theta _{+k}),\chi _{lk}^J=\chi _J(\theta _{+l}^{})\chi _J(\theta _{+k}).$$
(4.4.4)
The transformation functions are the operators in the space of internal degrees of freedom labeled by $`(\pm k)`$ corresponding to distortion rotations around the axes $`(\pm k)`$ by the angles $`\theta _{\pm k}`$. We make proposition that the distortion rotations are incompatible, namely the transformation operators $`f_{lk}^{(\pm )}`$ obey the incompatibility relations
$$\begin{array}{c}f_{lk}^{(+)}f_{cd}^{(+)}f_{ld}^{(+)}f_{ck}^{(+)}=f^{(+)}\epsilon _{lcm}\epsilon _{kdn}f_{nm}^{()},\hfill \\ \\ f_{kl}^{()}f_{dc}^{()}f_{dl}^{()}f_{kc}^{()}=f^{()}\epsilon _{lcm}\epsilon _{kdn}f_{mn}^{(+)},\hfill \end{array}$$
(4.4.5)
where $`l,k,c,d,m,n=1,2,3`$. The relations eq.(4.4.5) hold in general for both the local and global rotations. In the following we shall often be concerned with these most important relations. Making use of eq.(4.3.1), eq.(4.3.2) and eq.(4.4.2), one gets the transformations implemented upon the quark field, which in matrix notation take the form $`q^{}(\zeta )=U(\theta (\zeta ))q(\zeta ),\overline{q^{}}(\zeta )=\overline{q}(\zeta )U^+(\theta (\zeta )),`$ where $`q=\{q_k\},U(\theta )=\{f_{lk}^{(+)}\}`$. Under the incompatibility commutation relations (4.4.5), the transformation matrices $`\{U\}`$ generate the unitary group of internal symmetries $`U(1),SU(2),`$ $`SU(3)`$. Since the distorted ordinary structures have made contribution in the realization of geometry $`G`$ instead of regular ones, then, stated somewhat differently the principle of identity of regular structures directly leads to the equivalent principle (the gauge principle): an action integral of any dynamical physical system must be invariant under arbitrary transformations eq.(4.4.2). Below we discuss different possible models.
1. In the simple case of one dimensional local transformations through the local angles $`\theta _{+1}(\zeta )`$ and $`\theta _1(\zeta )onehas`$ $`f^{(+)}=\left(\begin{array}{ccc}f_{11}^{(+)}& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),f^{()}=(f^{(+)})^+.`$ The incompatibility relations eq.(4.4.5) reduce to identity $`f_{11}^{(+)}=f^{(+)}.`$ At $`\chi _R(\theta _{+1})=\chi _R(\theta _1)`$ the transformations
$$f_{11}^{(+)}=U(\theta )=f(\theta _{+1}(\zeta ),\theta _1(\zeta ))=\mathrm{exp}\{i\chi _J^{(+)}(\theta _{+1})+i\chi _J^{()}(\theta _1)\}.$$
(4.4.6)
generate a commutative Abelian unitary local group of electromagnetic interactions realized as the Lie group $`U^{loc}(1)=SO^{loc}(2)`$ with one dimensional trivial algebra $`\widehat{g}_1=R^1`$: $`U(\theta )=e^{i\theta }`$, where $`\theta \chi _J^{(+)}(\theta _{+1})\chi _J^{()}(\theta _1)`$. The strength of interaction is specified by a single coupling $`Q`$ of electric charge. The invariance under the local group $`U^{loc}(1)`$ leads to electromagnetic field, the massless quanta of which - the photons are electrically neutral, because of the condition eq.(4.3.4):
$$f(\theta _{+1},\theta _1)=f(\theta _{+1}^{},\theta _1^{})=\mathrm{}=inv.$$
(4.4.7)
2. Next, we consider the case of two dimensional local transformations through the angles $`\theta _{\pm m}(\zeta )`$ around two axes $`(m=1,2)`$. The matrix function of transformation is written $`f^{(+)}=\left(\begin{array}{ccc}f_{11}^{(+)}& f_{12}^{(+)}& 0\\ f_{21}^{(+)}& f_{22}^{(+)}& 0\\ 0& 0& 1\end{array}\right),f^{()}=(f^{(+)})^+.`$ The incompatibility relations eq.(4.4.5) give rise to nontrivial conditions
$$\begin{array}{cc}f_{11}^{(+)}=f^{(+)}(f_{22}^{(+)})^{},f_{21}^{(+)}=f^{(+)}(f_{12}^{(+)})^{},\hfill & \\ & \\ f_{12}^{(+)}=f^{(+)}(f_{21}^{(+)})^{},f_{22}^{(+)}=f^{(+)}(f_{11}^{(+)})^{},\hfill & \end{array}$$
(4.4.8)
Hence $`f^{(+)}=1.`$ One readily derives the matrix $`U(\theta )`$ of gauge transformations of collection of fundamental fields $`U=e^{i\stackrel{}{T}\stackrel{}{\theta }}=\left(\begin{array}{cc}f_{11}^{(+)}& f_{12}^{(+)}\\ f_{21}^{(+)}& f_{22}^{(+)}\end{array}\right),`$ where $`T_i(i=1,2,3)`$ are the generators of the group $`SU(2)`$. The fields will come in multiplets forming a basis for representations of the isospin group $`SU(2)`$. Meanwhile
$$\begin{array}{c}\frac{\theta _1}{\theta }=\frac{e^{\chi _{12}^R}\mathrm{sin}\chi _{12}^J}{\sqrt{1e^{2\chi _{11}^R}\mathrm{cos}^2\chi _{11}^J}},\frac{\theta _2}{\theta }=\frac{e^{\chi _{12}^R}\mathrm{cos}\chi _{12}^J}{\sqrt{1e^{2\chi _{11}^R}\mathrm{cos}^2\chi _{11}^J}},\frac{\theta _3}{\theta }=\frac{e^{\chi _{11}^R}\mathrm{sin}\chi _{11}^J}{\sqrt{1e^{2\chi _{11}^R}\mathrm{cos}^2\chi _{11}^J}},\hfill \\ \\ \theta =\stackrel{}{\theta }=2\mathrm{arccos}\left(e^{\chi _{11}^R}\mathrm{cos}\chi _{11}^J\right),e^{\chi _{11}^R}1,\hfill \end{array}$$
(4.4.9)
provided by
$$\chi _{11}^R=\chi _{22}^R,\chi _{12}^R=\chi _{21}^R\chi _{11}^J+\chi _{22}^J=0,\chi _{21}^J+\chi _{12}^J=\pi ,\chi _{12}^R=\frac{1}{2}\mathrm{ln}\left(1e^{2\chi _{11}^R}\right).$$
(4.4.10)
That is, three functions $`\chi _{11}^R,\chi _{11}^J`$ and $`\chi _{12}^J`$ or the angles $`\theta _{+1}^{},\theta _{+1}`$ and $`\theta _{+2}`$ are parameters of the group $`SU^{loc}(2)`$
$$\chi _{11}^R=\chi _R(\theta _{+1}^{})\chi _R(\theta _{+1}),\chi _{11}^J=\chi _J(\theta _{+1}^{})\chi _J(\theta _{+1}),\chi _{12}^J=\chi _J(\theta _{+1}^{})\chi _J(\theta _{+2}).$$
(4.4.11)
3. In the case of gauge transformations occurred around all three axes $`(l,k=1,2,3)`$: $`f^{(+)}=\left(\begin{array}{ccc}f_{11}^{(+)}& f_{12}^{(+)}& f_{13}^{(+)}\\ f_{21}^{(+)}& f_{22}^{(+)}& f_{23}^{(+)}\\ f_{31}^{(+)}& f_{32}^{(+)}& f_{33}^{(+)}\end{array}\right),f^{()}=(f^{(+)})^+,`$ the incompatibility relations eq.(4.4.5) yield the unitary condition $`U^1=U^+,f^{(+)}U`$, and also $`U=1`$. Then $`U(\theta )=e^{\frac{i}{2}\stackrel{}{\lambda }\stackrel{}{\theta }}`$, where $`{\displaystyle \frac{\lambda _i}{2}}(i=1,\mathrm{},8)`$ are the matrix representation of generators of the group $`SU(3)`$. Right through differentiation one derives $`\stackrel{}{\lambda }\stackrel{}{d\theta }=2iU^+dU,`$ or $`\stackrel{}{\theta }={\displaystyle Im\left(tr\left(\stackrel{}{\lambda }\left(f^{()}df^{(+)}\right)\right)\right)},`$ provided by $`Re\left(tr\left(\stackrel{}{\lambda }\left(f^{()}df^{(+)}\right)\right)\right)0.`$ For the infinitesimal transformations $`\theta _i1`$ we get
$$\begin{array}{cc}\theta _12e^{\chi _{12}^R}\mathrm{sin}\chi _{12}^J,\theta _3\mathrm{sin}\chi _{33}^J+2\mathrm{sin}\chi _{11}^J,\theta _52(1e^{\chi _{13}^R}\mathrm{cos}\chi _{13}^J),\hfill & \\ \theta _22(1e^{\chi _{12}^R}\mathrm{cos}\chi _{12}^J),\theta _42e^{\chi _{13}^R}\mathrm{sin}\chi _{13}^J,\theta _62e^{\chi _{23}^R}\mathrm{sin}\chi _{23}^J,\hfill & \\ \theta _72(1e^{\chi _{23}^R}\mathrm{cos}\chi _{23}^J),\theta _8\sqrt{3}\mathrm{sin}\chi _{33}^J,\hfill & \end{array}$$
(4.4.12)
provided by
$$\chi _{ll}^R0,\chi _{lk}^R\chi _{kl}^R,\chi _{lk}^J\chi _{kl}^J,(lk)\mathrm{sin}\chi _{11}^J+\mathrm{sin}\chi _{22}^J+\mathrm{sin}\chi _{33}^J0.$$
(4.4.13)
The internal symmetry group $`SU_C^{loc}(3)`$ enables to introduce the gauge theory in colour space, with the colour charges as exactly conserved quantities. While, the local colour transformations are implemented on the coloured quarks right through the $`SU_C^{loc}(3)`$ rotation matrix U in the fundamental representation.
## 5 Operator Multimanifold $`\widehat{G}_N`$
### 5.1 Operator Vector and Covector Fields
The OM formalism of $`\widehat{G}=\underset{\eta }{\widehat{G}}\underset{u}{\widehat{G}}`$ is built up by assuming an existence only of ordinary primordial structures of one sort (one u-channel). Being confronted by our major goal to develop the microscopic approach to field theory based on multiworld geometry, henceforth we generalize the OM formalism via the concept of the OMM. Then, instead of one sort of ordinary structures we are going to deal with different species of ordinary structures. But before proceeding further and to enlarge the previous model it is profitable to assume an existence of infinite number of $`{}_{}{}^{i}u`$-type ordinary structures of different species $`i=1,2,\mathrm{},N`$ (multi-u channel). These structures will be specified by the superscript $`i`$ to the left. This hypothesis, as it will be seen in the subsequent part II, leads to the substantial progress of understanding of the properties of particles. At the very outset we consider the processes of creation and annihilation of regular structures of $`\eta `$\- and $`{}_{}{}^{i}u`$-types in the lowest state ($`s_0`$). The general rules stated in subsec 2.1 regarding to this change apply a substitution of operator basis pseudo vectors and covectors by a new ones $`(i=1,2,\mathrm{},N)`$
$${}_{}{}^{i}\widehat{O}_{\lambda ,\mu }^{r_1r_2}={}_{}{}^{i}\widehat{O}_{\lambda }^{r_1}{}_{}{}^{i}\widehat{O}_{\mu }^{r_2}{}_{}{}^{i}\widehat{O}_{\lambda ,\mu }^{r}={}_{}{}^{i}O_{\lambda ,\mu }^{r}(\alpha _\lambda \alpha _\mu ),$$
(5.1.1)
provided by $`r(r_1,r_2)`$ and
$$\begin{array}{c}{}_{}{}^{i}O_{1,1}^{r}=\frac{1}{\sqrt{2}}(\nu _i\underset{\eta }{𝑂}{}_{+}{}^{r}+{}_{}{}^{i}\underset{u}{𝑂}{}_{+}{}^{r}),{}_{}{}^{i}O_{2,1}^{r}=\frac{1}{\sqrt{2}}(\nu _i\underset{\eta }{𝑂}{}_{+}{}^{r}{}_{}{}^{i}\underset{u}{𝑂}{}_{+}{}^{r}),\hfill \\ \\ {}_{}{}^{i}O_{1,2}^{r}=\frac{1}{\sqrt{2}}(\nu _i\underset{\eta }{𝑂}{}_{}{}^{r}+{}_{}{}^{i}\underset{u}{𝑂}{}_{}{}^{r}),{}_{}{}^{i}O_{2,2}^{r}=\frac{1}{\sqrt{2}}(\nu _i\underset{\eta }{𝑂}{}_{}{}^{r}{}_{}{}^{i}\underset{u}{𝑂}{}_{}{}^{r}),\hfill \end{array}$$
where
$$<\nu _i,\nu _j>=\delta _{ij},<{}_{}{}^{i}\underset{u}{𝑂}{}_{\lambda }{}^{r},{}_{}{}^{j}\underset{u}{𝑂}{}_{\tau }{}^{r^{}}>=\delta _{ij}\delta _{rr^{}}{}_{}{}^{}\delta _{\lambda \tau }^{},<\underset{\eta }{𝑂}{}_{\lambda }{}^{r},{}_{}{}^{i}\underset{u}{𝑂}{}_{\tau }{}^{r^{}}>=0.$$
(5.1.2)
We consider then the operators $`{}_{}{}^{i}\widehat{\gamma }{}_{(\lambda ,\mu ,\alpha )}{}^{r}={}_{}{}^{i}\widehat{O}_{\lambda ,\mu }^{r_1r_2}\widehat{\sigma }_\alpha ^{r_3}.`$ and calculate nonzero matrix elements
$$<\lambda ,\mu {}_{}{}^{i}\widehat{\gamma }{}_{(\tau ,\nu ,\alpha )}{}^{r}\tau ,\nu >={}_{}{}^{}\delta _{\lambda \tau }^{}{}_{}{}^{}\delta _{\mu ,\nu }^{}{}_{}{}^{i}e_{(\tau ,\nu ,\alpha )}^{r},$$
(5.1.3)
where $`{}_{}{}^{i}e_{(\lambda ,\mu ,\alpha )}^{r}={}_{}{}^{i}O_{\lambda ,\mu }^{r}\sigma _\alpha `$. The operators $`\{{}_{}{}^{i}\widehat{\gamma }{}_{}{}^{r}\}`$ are the basis for all the operator vectors $`\widehat{\mathrm{\Phi }}(\zeta )={}_{}{}^{i}\widehat{\gamma }{}_{}{}^{r}{}_{}{}^{i}\mathrm{\Phi }_{r}^{}(\zeta )`$ of tangent section of principle bundle with the base of operator multimanifold $`\widehat{G}_N=({\displaystyle \underset{i}{\overset{N}{}}}{}_{}{}^{}\widehat{\underset{i}{𝑅}}{}_{}{}^{4})\widehat{R}^3`$. Here $`{}_{}{}^{}\widehat{\underset{i}{𝑅}}^4`$ is the $`2\times 2`$ dimensional linear pseudo operator space, with the set of the linear unit operator pseudo vectors eq.(5.1.1) as the basis of tangent vector section, and $`\widehat{R}^3`$ is the three dimensional real linear operator space with the basis consisted of the ordinary unit operator vectors $`\{\widehat{\sigma }_\alpha ^r\}`$. The $`\widehat{G}_N`$ decomposes as follows:
$$\widehat{G}_N=\underset{\eta }{\widehat{G}}\underset{u_1}{\widehat{G}}\mathrm{}\underset{u_N}{\widehat{G}},$$
(5.1.4)
where $`\underset{u_i}{\widehat{G}}`$ is the six dimensional operator manifold of the given species $`(i)`$ with the basis
$`\left\{{}_{}{}^{i}\underset{u}{\widehat{\gamma }}{}_{(\lambda \alpha )}{}^{r}={}_{}{}^{i}\widehat{\underset{u}{𝑂}}{}_{\lambda }{}^{r}\widehat{\sigma }_\alpha ^r\right\}`$. The expansions of operator vectors and covectors are written $`\underset{\eta }{\widehat{\mathrm{\Psi }}}=\underset{\eta }{\widehat{\gamma }}{}_{}{}^{r}\underset{\eta }{\Psi }{}_{r}{}^{},\underset{u}{\widehat{\mathrm{\Psi }}}={}_{}{}^{i}\underset{u}{\widehat{\gamma }}{}_{}{}^{r}{}_{}{}^{i}\underset{u}{\Psi }{}_{r}{}^{},\overline{\underset{\eta }{\widehat{\mathrm{\Psi }}}}=\underset{\eta }{\widehat{\gamma }}{}_{r}{}^{}\underset{\eta }{\Psi }{}_{}{}^{r},\overline{\underset{u}{\widehat{\mathrm{\Psi }}}}={}_{}{}^{i}\underset{u}{\widehat{\gamma }}{}_{r}{}^{}{}_{}{}^{i}\underset{u}{\Psi }{}_{}{}^{r},`$ where the components $`\underset{\eta }{\Psi }{}_{r}{}^{}(\eta )`$ and $`{}_{}{}^{i}\underset{u}{\Psi }{}_{r}{}^{}(u)`$ are respectively the link functions of $`\eta `$-type and $`{}_{}{}^{i}u`$-type structures.
### 5.2 Realization of the Multimanifold $`G_N`$
Now, we consider the special system of the regular structures, which is made of one fundamental structure of $`\eta `$-type and infinite number of $`{}_{}{}^{i}u`$-type ordinary structures of different species $`(i=1,\mathrm{},N)`$. The primordial structures establish the stable linkage to form the stable system
$$p^2=p_\eta ^2\underset{i=1}{\overset{N}{}}p_{u_i}^2=0.$$
(5.2.1)
The free field defined on the multimanifold $`G_N=\underset{\eta }{𝐺}\underset{u_1}{𝐺}\mathrm{}\underset{u_N}{𝐺}`$ is written
$$\mathrm{\Psi }=\underset{\eta }{\Psi }(\eta )\underset{u}{\Psi }(u),\underset{u}{\Psi }(u)=\underset{u_1}{\Psi }(u_1)\mathrm{}\underset{u_N}{\Psi }(u_N),$$
where $`\underset{u_i}{\Psi }`$ is the bispinor defined on the internal manifold $`\underset{u_i}{𝐺}`$. On analogy of subsec.2.2 we make use of localized wave packets by means of superposition of plane wave solutions furnished by creation and annihilation operators in agreement with Pauli’s principle. Straightforward calculations now give the generalization of the relation eq.(2.2.2)
$$\begin{array}{c}\underset{\lambda =\pm }{}<\chi _\lambda \widehat{\mathrm{\Phi }}(\zeta )\overline{\widehat{\mathrm{\Phi }}}(\zeta )\chi _\lambda >=\underset{\lambda =\pm }{}<\chi _\lambda \overline{\widehat{\mathrm{\Phi }}}(\zeta )\widehat{\mathrm{\Phi }}(\zeta )\chi _\lambda >=\hfill \\ \\ i\zeta ^2\underset{\zeta }{𝐺}(0)=i\left(\eta ^2\underset{\eta }{𝐺}(0)\underset{i=1}{\overset{N}{}}u_{i}^{}{}_{}{}^{2}\underset{u_i}{𝐺}(0)\right).\hfill \end{array}$$
(5.2.2)
Along the same line the realization of the multimanifold stems from the condition eq.(2.2.3), which is now imposed upon the matrix element eq.(5.2.2). Let denote $`u^2\underset{u}{𝐺}(0)\underset{u_iu_i^{}}{lim}{\displaystyle \underset{i=1}{\overset{N}{}}}(u_iu_i^{})\underset{u_i}{𝐺}(u_iu_i^{})`$ and consider a stable system eq.(5.2.1). Hence
$$\underset{u}{𝐺}{}_{F}{}^{}(0)=\underset{\eta }{𝐺}{}_{F}{}^{}(0)=\underset{\zeta }{𝐺}{}_{F}{}^{}(0),$$
(5.2.3)
where $`\underset{\eta }{𝐺}{}_{F}{}^{},\underset{u}{𝐺}_F`$ and $`\underset{\zeta }{𝐺}_F`$ are the causal Green’s functions of the $`\eta ,u`$ and $`\zeta `$-type structures, and $`m\left|p_u\right|=\left({\displaystyle \underset{i=1}{\overset{N}{}}}p_{u_i}^{}{}_{}{}^{2}\right)^{1/2}=\left|p_\eta \right|.`$ In the aftermath, the length of each vector $`\zeta ={}_{}{}^{i}e{}_{}{}^{i}\zeta G_N`$ should be equaled zero (subsec.2.2) $`\zeta ^2=\eta ^2u^2=\eta ^2{\displaystyle \underset{i=1}{\overset{N}{}}}(u_i^G)^2=0,`$ where use is made of
$$\left(u_i^G\right)^2u_i^2\underset{\begin{array}{c}u_iu_i^{}\hfill \\ \eta \eta ^{}\hfill \end{array}}{lim}\underset{u_i}{𝐺}{}_{F}{}^{}(u_iu_i^{})/\underset{\eta }{𝐺}{}_{F}{}^{}(\eta \eta ^{})$$
and $`u_i^G={}_{}{}^{i}\underset{u}{\widehat{e}}{}_{(\lambda ,\alpha )}{}^{}u_{i}^{G(\lambda ,\alpha )}`$. Thus, the multimanifold $`G_N`$ comes into being, which decomposes as follows:
$$G_N=\underset{\eta }{𝐺}\underset{u_1}{𝐺}\mathrm{}\underset{u_N}{𝐺}.$$
(5.2.4)
It brings us to the conclusion: the major requirement eq.(2.2.3) provided by stability condition eq.(5.2.1) or eq.(5.2.3) yields the flat multimanifold $`G_N`$. Meanwhile, the Minkowski flat space $`M_4`$ stems from the flat submanifold $`\underset{\eta }{𝐺}`$ (subsec. 2.1), in which the line element turned out to be invariant. That is, the principle of relativity comes into being with the $`M_4`$ ensued from the MW geometry $`G_N`$. In the following we shall use a notion of the $`i`$-th internal world for the submanifold $`\underset{u_i}{𝐺}`$.
### 5.3 Subquarks and Subcolour Confinement
Since our discussion within this section in many respects is similar to that of sec.4, here we will be brief. We assume that the distortion rotations ($`\underset{u_i}{𝐺}\stackrel{\theta }{}\stackrel{~}{\underset{u_i}{𝐺}}`$, for given $`i`$) through the angles $`{}_{}{}^{i}\theta _{+k}^{}`$ and $`{}_{}{}^{i}\theta _{k}^{}k=1,2,3`$ occur separately in the three dimensional internal spaces $`\underset{u_i}{𝑅}_+^3`$ and $`\underset{u_i}{𝑅}_{}^3`$ composing six dimensional distorted submanifold $`\stackrel{~}{\underset{u_i}{𝐺}}=\underset{u_i}{𝑅}{}_{+}{}^{3}\underset{u_i}{𝑅}_{}^3`$ (see eq.(5.2.4)). As it is exemplified in previous section, the laws apply in use the wave packets constructed by superposition of the link functions of distorted ordinary structures furnished by generalized operators of creation and annihilation as the expansion coefficients
$$\underset{u}{\widehat{\mathrm{\Psi }}}(\theta _+)=\underset{\pm s}{}\frac{d^3p_{u_i}}{(2\pi )^{3/2}}\left({}_{}{}^{i}\underset{u}{\widehat{\gamma }}{}_{(+\alpha )}{}^{k}{}_{}{}^{i}\underset{u}{\Psi }_{}^{(+\alpha )}({}_{}{}^{i}\theta _{+k}^{})+{}_{}{}^{i}\underset{u}{\widehat{\gamma }}{}_{(\alpha )}{}^{k}{}_{}{}^{i}\underset{u}{\Psi }_{}^{(\alpha )}({}_{}{}^{i}\theta _{+k}^{})\right),$$
(5.3.1)
etc. The fields $`{}_{}{}^{i}\underset{u}{\Psi }(\theta _{+k})`$ and $`{}_{}{}^{i}\underset{u}{\Psi }(\theta _k)`$ are defined on the distorted internal spaces $`\underset{u_i}{𝑅}_+^3`$ and $`\underset{u_i}{𝑅}_{}^3`$. The generalized expansion coefficients in eq.(5.3.1) imply
$$<\chi _{}\{{}_{}{}^{i}\underset{u}{\widehat{\gamma }}{}_{k}{}^{(+\alpha )}(p_{u_i},s_i),{}_{}{}^{j}\underset{u}{\widehat{\gamma }}{}_{(+\beta )}{}^{k^{}}(p_{u_j}^{},s_j^{})\}\chi _{}>=\delta _{ij}\delta _{kk^{}}\delta _{ss^{}}\delta _{\alpha \beta }\delta ^3(\stackrel{}{p}_{u_i}\stackrel{}{p^{}}_{u_i}).$$
(5.3.2)
The condition of the MW geometry realization eq.(5.2.3) now reduced to be
$$\underset{i=1}{\overset{N}{}}\omega _i\left[\underset{{}_{}{}^{i}\theta _{+}^{}{}_{}{}^{i}\theta _{}^{}}{lim}\underset{u_i}{𝐺}^\theta {}_{F}{}^{}({}_{}{}^{i}\theta _{+}^{}{}_{}{}^{i}\theta _{}^{})\right]=\underset{\eta _f\eta _f^{}}{lim}\underset{\eta }{𝐺}{}_{F}{}^{}(\eta _f\eta _f^{}),$$
(5.3.3)
provided by $`\omega _i={\displaystyle \frac{u_i^2}{u^2}}.`$ Taking into account the expression of causal Green’s function at given $`(i)`$, in the case if
$$\underset{u_{i_1}u_{i_2}}{lim}\underset{u_{i_1}}{𝐺}{}_{F}{}^{}(u_{i_1}u_{i_2})=\underset{u_{i_1}^{}u_{i_2}^{}}{lim}\underset{u_{i_1^{}}}{𝐺}{}_{F}{}^{}(u_{i_1}^{}u_{i_2}^{})=\mathrm{}=inv,$$
one gets
$$\underset{k}{}{}_{}{}^{i}\underset{u}{\Psi }({}_{}{}^{i}\theta _{+k}^{}){}_{}{}^{i}\overline{\underset{u}{\Psi }}({}_{}{}^{i}\theta _{k}^{})=\underset{k}{}{}_{}{}^{i}\underset{u}{\Psi }_{}^{}({}_{}{}^{i}\theta _{+k}^{}){}_{}{}^{i}\overline{\underset{u}{\Psi }}_{}^{}({}_{}{}^{i}\theta _{k}^{})=\mathrm{}=inv.$$
(5.3.4)
Thus, in the context of the physically more realistic MW geometry it is legitimate now to substitute the concept of quark $`(q_k)`$ schematically introduced in sec.4 by the subquark $`({}_{}{}^{i}q_{k}^{})`$. Everything said will then remain valid, provided we make a simple change of quarks into subquarks, the colours into subcolours. Hence, we may think of the function $`{}_{}{}^{i}\underset{u}{\Psi }({}_{}{}^{i}\theta _{+k}^{})`$ as the $`u`$-component of bispinor field of subquark $`({}_{}{}^{i}q_{k}^{})`$ of species $`(i)`$ with subcolour $`k`$, and respectively $`{}_{}{}^{i}\overline{\underset{u}{\Psi }}({}_{}{}^{i}\theta _{k}^{})`$ -the conjugated bispinor field of antisubcolour $`(k)`$. The subquarks and antisubquarks may be local $`({}_{}{}^{i}q_{k}^{})`$ or global $`({}_{}{}^{i}q_{k}^{c})`$. Then, the subquark $`({}_{}{}^{i}q_{k}^{})`$ is the fermion with the half integer spin and subcolour degree of freedom, and, according to eq.(5.3.4), could emerge on the mass shell only in confined phase
$$\underset{k}{}{}_{}{}^{i}q_{k}^{}{}_{}{}^{i}\overline{q}_{k}^{}=\underset{k}{}{}_{}{}^{i}q_{}^{}{}_{k}{}^{}{}_{}{}^{i}\overline{q^{}}_{k}^{}=\mathrm{}=inv.$$
(5.3.5)
To trace a resemblance with the previous section, the internal symmetry group $`{}_{}{}^{i}G=U(1),SU(2),SU(3)`$ enables to introduce the gauge theory in internal world with the subcolour charges as exactly conserved quantities. Furthermore, the subcolour transformation have implemented on subquark fields right through local and global rotation matrices of group $`{}_{}{}^{i}G`$ in fundamental representation. Due to the Noether procedure the conservation of global charges ensued from the global gauge invariance of physical system, meanwhile reinforced requirement of local gauge invariance may be satisfied as well by introducing the gauge fields with the values in Lie algebra $`{}_{}{}^{i}\widehat{g}`$ of the group $`{}_{}{}^{i}G`$.
Part II
Realization of the Particle Physics
Continuing our program based on the OMM formalism, in this part we attempt to develop, further, the microscopic approach to the SM, which enables an insight to the key problems of particle phenomenology. Particularly, we suggest the microscopic theory of the unified electroweak interactions with a small number of free parameters. Besides the microscopic interpretation of all physical parameters the resulting theory has two testable solid implications, which are drastically different from those of conventional models.
## 6 The MW-Structure of the Particles
For our immediate purpose to describe the particles as the composite dynamical systems defined on the MW geometry, we shall consider the collection of matter fields $`\mathrm{\Psi }(\zeta )`$ with nontrivial internal structure $`\mathrm{\Psi }(\zeta )=\underset{\eta }{\Psi }(\eta ){}_{}{}^{1}\underset{u}{\Psi }(\theta _1)\mathrm{}{}_{}{}^{N}\underset{u}{\Psi }(\theta _N).`$ We suppose that the component $`{}_{}{}^{i}\underset{u}{\Psi }(\theta _i)`$ is made of product of some constituent subquarks and antisubquarks, which form the multiplets transformed by fundamental $`{}_{}{}^{i}D(j)`$ and contragradient $`{}_{}{}^{i}\overline{D}(j)`$ irreducible representations of group $`{}_{}{}^{i}G`$. Hereinafter we suppose that the MW index $`(i)`$ will be running only through $`i=Q,W,B,`$ $`s,c,b,t`$ specifying the internal worlds formally taken to denote in following nomenclature: Q-world of electric charge; W-world of weak interactions; B-baryonic world of strong interactions; the s,c,b,t are the worlds of strangeness, charm, bottom and top. We admit also that the distortion rotations in the worlds Q,W and B are local $`{}_{}{}^{i}\theta _{\pm k}^{}(\eta )`$, while they are global in the worlds s,c,b,t. Below we introduce the fields of leptons $`(l)`$ and quarks $`(q_f)`$ with different flavours f=u,d,s,c,b,t. To develop some feeling for this problem and to avoid irrelevant complications, here we may temporarily skeletonize it by taking the leptons to have following MW- structure:
$$l\mathrm{\Psi }_l(\zeta )=\underset{\eta }{\Psi }(\eta )\underset{Q}{\Psi }(u_Q)\underset{W}{\Psi }(u_W),$$
(6.1)
while the quarks are in the form
$$q_f\mathrm{\Psi }_f(\zeta )=\underset{\eta }{\Psi }(\eta )\underset{Q}{\Psi }(u_Q)\underset{W}{\Psi }(u_W)\underset{B}{\Psi }(u_B)q_f^c,$$
(6.2)
where the superscript $`(c)`$ specified the worlds in which rotations are global
$$\begin{array}{c}q_u^c=q_d^c=1,q_s^c=\underset{s}{\Psi }^c(u_s),q_c^c\underset{c}{\Psi }(u_c),q_b^c\underset{b}{\Psi }(u_b),q_t^c\underset{t}{\Psi }(u_t).\hfill \end{array}$$
(6.3)
We can take this scheme as a starting point for our considerations, which will become clearer in the next sections where the explicit forms of the $`\underset{i}{\Psi }(u_i)`$ will be subject for discussion. We assign a scale $`1/3`$ to each distortion rotational mode in the three dimensional spaces $`\underset{u_i}{𝑅}_+^3`$ and $`\underset{u_i}{𝑅}{}_{}{}^{3},`$ namely the subquark arisen after the rotation around the given axis of the given world carries the $`1/3`$ charge of corresponding species; while, the antisubquark carries respectively the $`(1/3)`$ charge. In the case of the worlds C=s,c,b,t, where distortion rotations are global and diagonal with respect to axes 1,2,3, the physical system of corresponding subquarks is invariant under the global transformations $`f_C^{(3)}(\theta ^c)`$ of the global unitary group $`SU_3^c`$:
$$f_C^{(3)}=\left(\begin{array}{ccc}f_{11}^c& 0& 0\\ 0& f_{22}^c& 0\\ 0& 0& f_{33}^c\end{array}\right)=\mathrm{exp}\left\{\frac{i}{3}\left(\begin{array}{ccc}\theta _1^c& 0& 0\\ 0& \theta _2^c& 0\\ 0& 0& \theta _3^c\end{array}\right)\right\},$$
where $`f_C^{(3)}\left(f_C^{(3)}\right)^+=1,f_C^{(3)}=1`$. That is $`\theta _1^c+\theta _2^c+\theta _3^c=0.`$ The simplest possibility gives $`\theta ^c\theta _1^c=\theta _2^c`$. Hence, one gets
$$f_C^{(3)}=\mathrm{exp}\left\{\frac{i}{3}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right)\theta ^c\right\}=\mathrm{exp}\left(i\frac{\lambda _8}{\sqrt{3}}\theta ^c\right)=e^{iY^c\theta ^c},$$
(6.4)
provided by the operator of hypercharge $`Y^c`$ of diagonal group $`SU_i^c`$. If one has all the worlds involved, then $`Y^c=s+c+b+t.`$ The Q-world, which has an essential role for realization of the condition of the MW connections, will be discussed in detail in the next section. We only note that conservation of each rotation mode in Q- and B- worlds, where the distortion rotations are local, means that corresponding subquarks carry respectively the conserved charges Q and B in the scale $`1/3`$, and antisubquarks - $`(1/3)`$ charges. It can be provided by including the matrix $`\lambda _8`$ as the generator with the others in the symmetries of corresponding worlds (Q, B), and expressed in the invariance of the system of corresponding subquarks under the transformations of these symmetries. The incompatibility relations eq.(4.4.5) for global distortion rotations in the worlds C=s,c,b,t reduced to
$$f_{11}^cf_{22}^c=\overline{f}_{33}^c,f_{22}^cf_{33}^c=\overline{f}_{11}^c,f_{33}^cf_{11}^c=\overline{f}_{22}^c,$$
where $`f_C^{(3)}=f_{11}^cf_{22}^cf_{33}^c=1,f_{ii}^c\overline{f}_{ii}^c=1\text{for}i=1,2,3.`$ It means that the two subcolour singlets are available: $`\left(q\overline{q}\right)_i^c=inv,\left(q_1q_2q_3\right)^c=inv,`$ carried the charges $`C_{\left(q\overline{q}\right)_i^c}=0C_{\left(q_1q_2q_3\right)^c}=1,`$ respectively, where we denote $`\left(q\overline{q}\right)_i^c{}_{}{}^{c}q_{i}^{}{}_{}{}^{c}\overline{q}_{i}^{},\left(q_1q_2q_3\right)^c{}_{}{}^{c}q_{1}^{}{}_{}{}^{c}q_{2}^{}{}_{}{}^{c}q_{3}^{}.`$ Including the baryonic charge into strong hypercharge
$$Y=B+s+c+b+t,$$
(6.5)
we conclude that the hypercharge Y is the sum of all the conserved quantum numbers associated with the corresponding rotational modes of the internal worlds B,s,c,b,t involved in the MW geometry realization condition eq.(5.2.3).
### 6.1 Realization of Q-World and Gell-Mann-Nishijima Relation
The symmetry of the Q-world is assumed to be a local unitary symmetry $`diag\left(SU^{loc}(3)\right)`$ is diagonal with respect to axes 1,2,3. The unitary unimodular matrix $`f_Q^{(3)}`$ of local distortion rotations takes the form
$$f_Q^{(3)}=\left(\begin{array}{ccc}f_{11}^Q& 0& 0\\ 0& f_{22}^Q& 0\\ 0& 0& f_{33}^Q\end{array}\right)=Q_1e^{i\theta _1}+Q_2e^{i\theta _2}+Q_3e^{i\theta _3}=e^{i\stackrel{}{Q}\stackrel{}{\theta }}=e^{i\lambda _Q\theta _Q},$$
where $`f_Q^{(3)}\left(f_Q^{(3)}\right)^+=1,f_Q^{(3)}=1`$, provided
$$Q_1=\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right),Q_2=\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right),Q_3=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 0& 1\end{array}\right),$$
and $`\theta _1+\theta _2+\theta _3=0.`$ Taking into account the scale of rotation mode, in the other simple case than those of eq.(6.4): $`\theta _2=\theta _3={\displaystyle \frac{1}{3}}\theta _Q,`$ it follows that $`\theta _1={\displaystyle \frac{2}{3}}\theta _Q.`$ Among the generators of the group $`SU(3)`$ only the matrices $`\lambda _3`$ and $`\lambda _8`$ are diagonal. Therefore, the matrix $`\lambda _Q`$ may be written $`\lambda _Q={\displaystyle \frac{1}{2}}\lambda _3+{\displaystyle \frac{1}{2\sqrt{3}}}\lambda _8.`$ Making use of the corresponding operators of the group $`SU(3)`$ we arrive at Gell-Mann-Nishijima relation
$$Q=T_3+\frac{1}{2}Y,$$
(6.1.1)
where $`Q=\lambda _Q`$ is the generator of electric charge, $`T_3={\displaystyle \frac{1}{2}}\lambda _3`$ is the third component of isospin $`\stackrel{}{T}`$, and $`Y={\displaystyle \frac{1}{\sqrt{3}}}\lambda _8`$ is the hypercharge. The eigenvalues of these operators will be defined in due course considering the concrete symmetries and microscopic structures of fundamental fields. We think of operators $`T_3`$ and $`Y`$ as the MW connection charges and of relation eq.(6.1.1) as the condition of the realization of the MW connections. Thus, during realization of the MW- structure the symmetries of corresponding internal worlds must be unified into more higher symmetry including also the $`\lambda _3`$ and $`\lambda _8`$. The realization conditions of the MW- structure are embodied in eq.(5.2.3) and eq.(6.1.1), provided by the conservation law of each rotational mode eq.(6.4) of the given internal worlds involved into the MW geometry realization condition. For example, in the case of quarks eq.(6.2), one has
$$\underset{i=B,s,c,b,t}{}\omega _i\underset{i}{𝐺}^\theta {}_{F}{}^{}(0)=\underset{\eta }{𝐺}{}_{F}{}^{}(0),$$
(6.1.2)
and the Gell-Mann-Nishijima relation is written down
$$Q=T_3+\frac{1}{2}(B+s+c+b+t).$$
(6.1.3)
The other case of leptons eq.(6.1) related closely to the realization of W-world of weak interactions will be discussed in detail in subsec.6.8, the realization condition for which reduces to following:
$$\underset{Q}{𝐺}^\theta {}_{F}{}^{}(0)=\underset{\eta }{𝐺}{}_{F}{}^{}(0),uu_Q,(iQ)$$
(6.1.4)
and
$$Q=T_3^w+\frac{1}{2}Y^w,$$
(6.1.5)
where $`T_3^w`$ and $`Y^w`$ are respectively the operators of third component of weak isospin $`\stackrel{}{T}^w`$ and weak hypercharge (subsec.6.2, 6.8). The incompatibility relations eq.(4.4.5) for the local distortion transformations in Q-world lead to
$$f_{11}^Qf_{22}^Q=\overline{f}_{33}^Q,f_{22}^Qf_{33}^Q=\overline{f}_{11}^Q,f_{33}^Qf_{11}^Q=\overline{f}_{22}^Q,$$
where $`f_{ii}^Q\overline{f}_{ii}^Q=1,\text{for}i=1,2,3,f_Q^{(3)}=f_{11}^Qf_{22}^Qf_{33}^Q=1.`$ This in turn suggests two subcolour singlets $`\left(q\overline{q}\right)_i^Q=inv,\left(q_1q_2q_3\right)^Q=inv,`$ with the electric charges $`Q_{\left(q\overline{q}\right)_i^Q}=0,Q_{\left(q_1q_2q_3\right)^Q}=1,`$ respectively. The singlets $`\left(q\overline{q}\right)_i^Q`$ for given $`i`$ allow us to think of the $`\left(q\overline{q}\right)^Q`$ system as the mixed ensemble, such that a fraction of the members with relative population $`L_1`$ are characterized by the $`(q_1)^Q`$, some other fraction with relative population $`L_2`$, by $`(q_2)^Q`$, and so on. Namely, the $`\left(q\overline{q}\right)^Q`$ ensemble can be regarded as a mixture of pure ensembles. The fractional populations are constrained to satisfy the normalization condition
$$\underset{i}{}L_i=1,L_i\left(q\overline{q}\right)_i^Q/\left(q\overline{q}\right)^Q.$$
(6.1.6)
The $`L_i`$ also imply the orthogonality condition ensued from the symmetry of the $`Q`$-world
$$<L_i,L_j>=0\text{if}ij.$$
(6.1.7)
This prompts us to define the usual quantum mechanical density operator
$$\rho _1^Q=\underset{i}{}L_i(q_i)^Q)(q_i)^{Q+},tr(\rho _1^Q)=1.$$
(6.1.8)
The eq.(6.1.6) suggests another singlets as well
$$\left(q_1q_2q_3\right)_i^QL_i\left(q_1q_2q_3\right)^Q,$$
(6.1.9)
which will be used to build up the MW-structures of the leptons.
### 6.2 The Symmetries of the W- and B-Worlds
$``$ The W-World
Invoking local group of weak hypercharge $`Y^w`$ ($`U^{loc}(1)`$), it will be seen in subsec. 6.8 that the symmetry of W-world of weak interactions rather is $`SU^{loc}(2)U^{loc}(1)`$. However, for the present it is worthwhile to restrict oneself by admitting that the symmetry of W-world is simply expressed by the group of weak isospin $`SU^{loc}(2)`$. Namely, from the very first we consider a case of two dimensional distortion transformations through the angles $`\theta _\pm `$ around two arbitrary axes in the W-world. In accordance with the results of sec.4, the fields of subquarks and antisubquarks will come in doublets, which form the basis for fundamental representation of weak isospin group $`SU^{loc}(2)`$ often called a “custodial” symmetry . The doublet states are complex linear combinations of up and down states of weak isotopic spin. Three possible doublets of six subquark states are $`\left(\begin{array}{c}q_1\\ q_2\end{array}\right)^w,\left(\begin{array}{c}q_2\\ q_3\end{array}\right)^w,\left(\begin{array}{c}q_3\\ q_1\end{array}\right)^w.`$
$``$ The B-World
The B-world is responsible for strong interactions. The internal symmetry group is $`SU_c^{loc}(3)`$ enabling to introduce gauge theory in subcolour space with subcolour charges as exactly conserved quantities (sec.4). The local distortion transformations are implemented on the subquarks $`(q_i)^B,i=1,2,3`$ through the $`SU_c^{loc}(3)`$ rotation matrix $`U`$ in the fundamental representation. Taking into account a conservation of rotation mode, each subquark carries $`(1/3)`$ baryonic charge, while the antisubquark carries the $`(1/3)`$ baryonic charge.
### 6.3 The Microscopic Structure of Leptons: <br>Lepton Generations
After a quantitative discussion of the properties of symmetries of internal worlds, below we will attempt to show how the known fermion fields of leptons and quarks fit into this scheme. In this section we start with the leptons. Taking into account the eq.(6.1.1), eq(6.1.4), we may consider six possible lepton fields forming three doublets of lepton generations $`\left(\begin{array}{c}\nu _e\\ e\end{array}\right),\left(\begin{array}{c}\nu _\mu \\ \mu \end{array}\right),\left(\begin{array}{c}\nu _\tau \\ \tau \end{array}\right),`$ where
$$\begin{array}{c}\{\begin{array}{c}\nu _e\underset{\eta }{\Psi }{}_{\nu _e}{}^{}(\eta )(q_1\overline{q_1})^Q(q_1)^w=L_e\underset{\eta }{\Psi }{}_{\nu _e}{}^{}(\eta )(q\overline{q})^Q(q_1)^w,\hfill \\ \\ e\underset{\eta }{\Psi }{}_{e}{}^{}(\eta )(\overline{q_1q_2q_3})_1^Q(q_2)^w=L_e\underset{\eta }{\Psi }{}_{e}{}^{}(\eta )(\overline{q_1q_2q_3})^Q(q_2)^w,\hfill \end{array}\hfill \\ \\ \{\begin{array}{c}\nu _\mu \underset{\eta }{\Psi }{}_{\nu _\mu }{}^{}(\eta )(q_2\overline{q_2})^Q(q_2)^w=L_\mu \underset{\eta }{\Psi }{}_{\nu _\mu }{}^{}(\eta )(q\overline{q})^Q(q_2)^w,\hfill \\ \\ \mu \underset{\eta }{\Psi }{}_{\mu }{}^{}(\eta )(\overline{q_1q_2q_3})_2^Q(q_3)^w=L_\mu \underset{\eta }{\Psi }{}_{\mu }{}^{}(\eta )(\overline{q_1q_2q_3})^Q(q_3)^w,\hfill \end{array}\hfill \\ \\ \{\begin{array}{c}\nu _\tau \underset{\eta }{\Psi }{}_{\nu _\tau }{}^{}(\eta )(q_3\overline{q_3})^Q(q_3)^w=L_\tau \underset{\eta }{\Psi }{}_{\nu _\tau }{}^{}(\eta )(q\overline{q})^Q(q_3)^w,\hfill \\ \\ \tau \underset{\eta }{\Psi }{}_{\tau }{}^{}(\eta )(\overline{q_1q_2q_3})_3^Q(q_1)^w=L_\tau \underset{\eta }{\Psi }{}_{\tau }{}^{}(\eta )(\overline{q_1q_2q_3})^Q(q_1)^w\hfill \end{array}.\hfill \end{array}$$
(6.3.1)
Here $`e,\mu ,\tau `$ are the electron, the muon and the tau meson, $`\nu _e,\nu _\mu ,\nu _\tau `$ are corresponding neutrinos, $`L_eL_1,L_\mu L_2,L_\tau L_3,`$ are leptonic charges. The leptons carry leptonic charges as follows: $`L_e:(e,\nu _e),`$ $`L_\mu :(\mu ,\nu _\mu )`$ and $`L_\tau :(\tau ,\nu _\tau ),`$ which are conserved in all interactions. The leptons carry also the weak isospins: $`T_3^w={\displaystyle \frac{1}{2}}`$ for $`\nu _e,\nu _\mu ,\nu _\tau `$; and $`T_3^w={\displaystyle \frac{1}{2}}`$ for $`e,\mu ,\tau `$, respectively, and following electric charges: $`Q_{\nu _e}=Q_{\nu _\mu }=Q_{\nu _\tau }=0,Q_e=Q_\mu =Q_\tau =1.`$ The Q-components $`\underset{Q}{\Psi }(u_Q)`$ of lepton fields eq.(6.3.1) are made of singlet combinations of subquarks in Q-world. They imply subcolour confinement eq.(6.1.4). Then, the MW geometry realization condition is already satisfied and leptons may emerge in free combinations without any constraint. Thus, in suggested scheme there are only three possible generations of six leptons with integer electric and leptonic charges have being free of confinement.
### 6.4 The Microscopic Structure of Quarks: <br>Quark Generations
The only possible MW- structures of 18 quark fields read
$$\begin{array}{c}\{\begin{array}{c}u_i\underset{\eta }{\Psi }{}_{u}{}^{}(\eta )(q_2q_3)^Q(q_1)^w(q_i^B),\hfill \\ \\ d_i\underset{\eta }{\Psi }{}_{d}{}^{}(\eta )(\overline{q}_1)^Q(q_2)^w(q_i^B),\hfill \end{array}\{\begin{array}{c}c_i\underset{\eta }{\Psi }{}_{c}{}^{}(\eta )(q_3q_1)^Q(q_2)^w(q_i^B)(q_c^c),\hfill \\ \\ s_i\underset{\eta }{\Psi }{}_{s}{}^{}(\eta )(\overline{q}_2)^Q(q_3)^w(q_i^B)(\overline{q}_s^c),\hfill \end{array}\hfill \\ \\ \{\begin{array}{c}t_i\underset{\eta }{\Psi }{}_{t}{}^{}(\eta )(q_1q_2)^Q(q_3)^w(q_i^B)(q_t^c),\hfill \\ \\ b_i\underset{\eta }{\Psi }{}_{b}{}^{}(\eta )(\overline{q}_3)^Q(q_1)^w(q_i^B)(\overline{q}_b^c),\hfill \end{array},\hfill \end{array}$$
(6.4.1)
where the subcolour index $`(i)`$ runs through $`i=1,2,3`$, the $`(q_f^c)`$ are given in eq.(6.3). Henceforth the subcolour index will be left implicit, but always a summation must be extended over all subcolours in B-world. These fields form three possible doublets of weak isospin in the W-world $`\left(\begin{array}{c}u\\ d\end{array}\right),\left(\begin{array}{c}c\\ s\end{array}\right),\left(\begin{array}{c}t\\ b\end{array}\right).`$ The quark flavour mixing and similar issues are left for treatment in sec.6.10. The corresponding electric charges of quarks read $`Q_u=Q_c=Q_t={\displaystyle \frac{2}{3}},Q_d=Q_s=Q_b={\displaystyle \frac{1}{3}},`$ in agreement with the rules governing the MW connections eq.(6.1.1), where the electric charge difference of up and down quarks implies $`\mathrm{\Delta }Q=\mathrm{\Delta }T_3^w=1.`$ The explicit form of structure of $`(q_f^c)`$ will be discussed in next section. Here we only note that all components of $`(q_f^c)`$ are made of singlet combinations of global subquarks in corresponding internal worlds. They obey a condition of subcolour confinement. According to eq.(6.1.2), the subcolour confinement condition for B-world still remains to be satisfied such that the total quark fields obey to confinement. Then quarks would not be free particles and unwanted states (since not seen) like quarks or diquarks etc. are eliminated by construction at the very beginning. Thus, three quark generations of six possible quark fields exist. They carry fractional electric and baryonic charges and imply a confinement. Their other charges are left to be discussed below. Although within considered schemes the subquarks are defined on the internal worlds, however the resulting $`\eta `$-components , which we are going to deal with to describe the leptons and quarks defined on the spacetime continuum, are affected by them. Actually, as it is seen in subsec.4.3 the rotation through the angle $`\theta _{+k}`$ yields a total subquark field
$$q_k(\theta )=\mathrm{\Psi }(\theta _{+k})=\underset{\eta }{\Psi }^0\underset{u}{\Psi }(\theta _{+k})$$
where $`\underset{\eta }{\Psi }^0`$ is the plane wave defined on $`\underset{\eta }{𝐺}`$. Hence, one gets
$$q_k(\theta (\eta ))=\underset{\eta }{\Psi }^0\underset{u}{𝑞}{}_{k}{}^{}(\theta (\eta ))=\underset{\eta }{𝑞}{}_{k}{}^{}(\theta (\eta ))\underset{u}{\Psi }^0,\underset{\eta }{𝑞}{}_{k}{}^{}(\theta (\eta ))f_{(+)}(\theta _{+k}(\eta ))\underset{\eta }{\Psi }^0,$$
where $`\underset{u}{\Psi }^0`$ is a plane wave defined on $`\underset{u}{𝐺}`$. The $`\underset{\eta }{𝑞}{}_{k}{}^{}(\theta (\eta ))`$ can be considered as the subquark field defined on the flat manifold $`\underset{\eta }{𝐺}`$ with the same quantum numbers of $`\underset{u}{𝑞}{}_{k}{}^{}(\theta (\eta ))`$. Thus, instead of the eq.(6.3.1) and eq.(6.4.1) we may consider on equal footing only the resulting $`\eta `$-components of leptons and quarks implying the given same structures. This enables to pass back to the Minkowski spacetime continuum $`\underset{\eta }{𝐺}M_4`$ (subsec.2.1).
### 6.5 The Flavour Group $`SU_f(6)`$
We adopt a simplified view-point on the field component $`(q_f^c)`$ (f=u,d,s,c,b,t) eq.(6.3) associated with the global distortion rotations in the given worlds s,c,b,t, such that they have following microscopic structure with corresponding global charges:
$$\begin{array}{c}q_u^c=q_d^c=1,\overline{q}_s^c=\left(\overline{q_1^cq_2^cq_3^c}\right)^s,s=1;q_c^c=\left(q_1^cq_2^cq_3^c\right)^c,c=1;\hfill \\ \\ \overline{q}_b^c=\left(\overline{q_1^cq_2^cq_3^c}\right)^b,b=1;q_t^c=\left(q_1^cq_2^cq_3^c\right)^t,t=1.\hfill \end{array}$$
(6.5.1)
To realize the MW-structure the global symmetries of internal worlds have unified into more higher symmetry including the generators $`\lambda _3`$ and $`\lambda _8`$ (subsec.6.1). This global group is the flavour group $`SU_f(6)`$ unifying all the symmetries $`SU_i^c`$ of the worlds Q,B,s,c,b,t: $`SU_f(6)SU_f(2)SU_B^cSU_s^cSU_c^cSU_b^cSU_t^c.`$ The total symmetry reads $`G_{tot}G^{loc}G^{glob}=G^{loc}SU_f(6),`$ provided $`G^{loc}SU^{loc}(3)G_w^{loc},`$ where $`G_w^{loc}`$ is the local symmetry of the electroweak interactions (subsec.6.8). The other important aspects of standard model are left for investigation in the next sections. However, below we proceed at once with further exposition of our approach to consider a gauge invariant Lagrangian of primary field with the MW- structure and nonlinear fermion interactions of the components.
### 6.6 The Primary Field
All the fields including the leptons eq.(6.3.1) and quarks eq.(6.4.1), along with the spacetime components have also the MW internal components made of the various constituent subquarks defined on the given internal worlds, such that the internal components are consisted of distorted ordinary structures (sec.5)
$$\mathrm{\Psi }(\theta )=\underset{\eta }{\Psi }(\eta )\underset{Q}{\Psi }(\theta _Q)\underset{W}{\Psi }(\theta _W)\underset{B}{\Psi }(\theta _B)\underset{C}{\Psi }(\theta ^c).$$
(6.6.1)
The components $`\underset{Q}{\Psi }(\theta _Q),\underset{W}{\Psi }(\theta _W),\underset{B}{\Psi }(\theta _B)`$ are primary massless bare Fermi fields. We assume that this field has arisen from primary field in the lowest state $`(s_0)`$ with the same field components consisted of regular ordinary structures, which is motivated by the argument given in the sec.5 that the regular ordinary structures directly could not take part in link exchange processes with the $`\eta `$-type regular structure. Therefore, the primary field defined on $`G_N`$
$$\mathrm{\Psi }(0)=\underset{\eta }{\Psi }(\eta )\underset{Q}{\Psi }(0)\underset{W}{\Psi }(0)\underset{B}{\Psi }(0)\underset{C}{\Psi }(0)$$
(6.6.2)
serves as the ready made frame into which the distorted ordinary structures of the same species should be involved. We apply the Lagrangian of this field possessed local gauge invariance written in the notations of App. D:
$$\stackrel{~}{L}_0(D)=\frac{i}{2}\{\overline{\mathrm{\Psi }}_e(\zeta ){}_{}{}^{i}\gamma \underset{i}{𝐷}\mathrm{\Psi }_e(\zeta )\underset{i}{𝐷}\overline{\mathrm{\Psi }}_e(\zeta ){}_{}{}^{i}\gamma \mathrm{\Psi }_e(\zeta )\},$$
(6.6.3)
with the vector indices contracted to form scalars, where $`\underset{i}{𝐷}=\underset{i}{}ig\underset{𝐢}{𝐁}(\zeta ),`$ $`\underset{𝐢}{𝐁}`$ are gauge fields. Since the components $`\underset{B}{\Psi }`$ and $`\underset{C}{\Psi }`$ will be of no consequence for a discussion, then we temporarily leave them implicit, namely $`i=\eta ,Q,W`$. The equation of primary field of the MW- structure with nonlinear fermion interactions of the components may be derived from an invariant action in terms of local gauge invariant Lagrangian, which looks like Heisenberg theory
$$\stackrel{~}{L}(D)=\stackrel{~}{L}_0(D)+\stackrel{~}{L}_I+\stackrel{~}{L}_B,$$
(6.6.4)
provided by the Lagrangians of nonlinear fermion interactions of the components $`\stackrel{~}{L}_I=\sqrt{2}\stackrel{~}{O}_1L_I,`$ and gauge field $`\stackrel{~}{L}_B=\sqrt{2}\stackrel{~}{O}_1L_B.`$ The binding interactions are in the form
$$\begin{array}{c}L_I=\underset{Q}{𝐿}{}_{I}{}^{}+\underset{W}{𝐿}{}_{I}{}^{},\underset{Q}{𝐿}{}_{I}{}^{}=\frac{\lambda }{4}(\underset{Q}{𝐽}{}_{L}{}^{}\underset{Q}{𝐽}{}_{R}{}^{+}+\underset{Q}{𝐽}{}_{R}{}^{}\underset{Q}{𝐽}{}_{L}{}^{+}),\underset{W}{𝐿}{}_{I}{}^{}=\frac{\lambda }{2}S_WS_W^+,\hfill \\ \\ L_B=\frac{1}{2}Tr(𝐁\overline{𝐁})=\frac{\mathrm{𝟏}}{\mathrm{𝟐}}\mathrm{𝐓𝐫}\left(\underset{𝐢}{𝐁}\underset{𝐢}{𝐁}\right),\hfill \end{array}$$
(6.6.5)
where
$$\begin{array}{c}\underset{Q}{𝐽}{}_{L,R}{}^{}=\underset{Q}{𝑉}\underset{Q}{𝐴},\underset{Q}{𝑉}{}_{(\lambda \alpha )}{}^{}=\underset{Q}{\overline{\mathrm{\Psi }}}\gamma _{(\lambda \alpha )}\underset{Q}{\Psi },\underset{Q}{𝑉}{}_{(\lambda \alpha )}{}^{+}=\underset{Q}{𝑉}{}_{}{}^{(\lambda \alpha )}=\underset{Q}{\overline{\mathrm{\Psi }}}\gamma ^{(\lambda \alpha )}\underset{Q}{\Psi }\hfill \\ \\ \underset{Q}{𝐴}{}_{(\lambda \alpha )}{}^{}=\underset{Q}{\overline{\mathrm{\Psi }}}\gamma _{(\lambda \alpha )}\gamma _5\underset{Q}{\Psi },\underset{Q}{𝐴}{}_{(\lambda \alpha )}{}^{+}=\underset{Q}{𝐴}^{(\lambda \alpha )}=\underset{Q}{\overline{\mathrm{\Psi }}}\gamma ^5\gamma ^{(\lambda \alpha )}\underset{Q}{\Psi },S_W=\underset{W}{\overline{\mathrm{\Psi }}}\underset{W}{\Psi },\hfill \end{array}$$
$`\gamma _\mu `$ and $`\gamma _5=i\gamma _0\gamma _1\gamma _2\gamma _3`$ are Dirac matrices. According to Fiertz theorem the interaction Lagrangian $`\underset{Q}{𝐿}{}_{I}{}^{}={\displaystyle \frac{\lambda }{2}}(VV^+AA^+)`$ may be written $`\underset{Q}{𝐿}{}_{I}{}^{}=\lambda (S_QS_Q^+P_QP_Q^+),`$ provided by $`S_Q=\underset{Q}{\overline{\mathrm{\Psi }}}\underset{Q}{\Psi },P_Q=\underset{Q}{\overline{\mathrm{\Psi }}}\gamma _5\underset{Q}{\Psi }.`$ Hence
$$\stackrel{~}{L}(D)=\sqrt{2}\stackrel{~}{O}_1L(D),L(D)=\underset{\eta }{𝐿}(\underset{\eta }{𝐷})\underset{Q}{𝐿}(\underset{Q}{𝐷})\underset{W}{𝐿}(\underset{W}{𝐷}),$$
(6.6.6)
where
$$\begin{array}{c}\underset{\eta }{𝐿}(\underset{\eta }{𝐷})=\underset{\eta }{𝐿}^{}{}_{0}{}^{(0)}(\underset{\eta }{𝐷})\frac{1}{2}Tr(\underset{\eta }{𝐁}\overline{\underset{\eta }{𝐁}}),\underset{Q}{𝐿}(\underset{Q}{𝐷})=\underset{Q}{𝐿}^{}{}_{0}{}^{(0)}(\underset{Q}{𝐷})\underset{Q}{𝐿}{}_{I}{}^{}\frac{1}{2}Tr(\underset{𝐐}{𝐁}\overline{\underset{𝐐}{𝐁}}),\hfill \\ \\ \underset{W}{𝐿}(\underset{W}{𝐷})=\underset{W}{𝐿}^{}{}_{0}{}^{(0)}(\underset{W}{𝐷})\underset{W}{𝐿}{}_{I}{}^{}\frac{1}{2}Tr(\underset{𝐖}{𝐁}\overline{\underset{𝐖}{𝐁}}).\hfill \end{array}$$
Here
$$\begin{array}{c}\underset{\eta }{𝐿}^{}{}_{0}{}^{(0)}=\frac{i}{2}\{\overline{\mathrm{\Psi }}\stackrel{}{\gamma \underset{\eta }{𝐷}}\mathrm{\Psi }\overline{\mathrm{\Psi }}\stackrel{}{\gamma \underset{\eta }{\stackrel{}{D}}}\mathrm{\Psi }\}=\underset{u}{\Psi }^+\underset{\eta }{𝐿}{}_{0}{}^{(0)}\underset{u}{\Psi },\hfill \\ \\ \underset{u}{𝐿}^{}{}_{0}{}^{(0)}=\frac{i}{2}\{\overline{\mathrm{\Psi }}\stackrel{}{\gamma \underset{u}{𝐷}}\mathrm{\Psi }\overline{\mathrm{\Psi }}\stackrel{}{\gamma \underset{u}{\stackrel{}{D}}}\mathrm{\Psi }\}=\underset{\eta }{\Psi }^+\underset{u}{𝐿}{}_{0}{}^{(0)}\underset{\eta }{\Psi },\hfill \end{array}$$
and
$$\underset{\eta }{𝐿}{}_{0}{}^{(0)}=\frac{i}{2}\{\underset{\eta }{\overline{\mathrm{\Psi }}}\stackrel{}{\gamma \underset{\eta }{𝐷}}\underset{\eta }{\Psi }\underset{\eta }{\overline{\mathrm{\Psi }}}\stackrel{}{\gamma \underset{\eta }{\stackrel{}{D}}}\underset{\eta }{\Psi }\},\underset{u}{𝐿}{}_{0}{}^{(0)}=\frac{i}{2}\{\underset{u}{\overline{\mathrm{\Psi }}}\stackrel{}{\gamma \underset{u}{𝐷}}\underset{u}{\Psi }\underset{u}{\overline{\mathrm{\Psi }}}\stackrel{}{\gamma \underset{u}{\stackrel{}{D}}}\underset{u}{\Psi }\}.$$
The Lagrangian eq.(6.6.6) has the global $`\gamma _5`$ and local gauge symmetries. We consider only $`\gamma _5`$ symmetry in Q-world, namely $`\underset{𝐐}{𝐁}0`$.
According to the OMM formalism, it is important to fix the mass shell of the stable MW- structure (eq.(5.2.1). It means that we must take at first the variation of the Lagrangian eq.(6.6.3) with respect to primary field eq.(6.6.2), then have switched on nonlinear fermion interactions of the components. In other words we take the variation of the Lagrangian eq.(14.6) with respect to the components on the fixed mass shell. The equations of free field ($`𝐁=0`$) of the MW-structure follow at once
$$\widehat{p}\mathrm{\Psi }_e(\zeta )=i\stackrel{}{\gamma }\mathrm{\Psi }_e(\zeta )=i{}_{}{}^{i}\gamma \underset{i}{}\mathrm{\Psi }_e(\zeta )=0,\overline{\mathrm{\Psi }}_e\stackrel{}{\widehat{p}}=i\underset{i}{}\overline{\mathrm{\Psi }}_e{}_{}{}^{i}\gamma =0,$$
(6.6.7)
which lead to separate equations for the massless components $`\underset{\eta }{\Psi }`$, $`\underset{Q}{\Psi }`$ and $`\underset{W}{\Psi }`$:
$$\gamma ^{(\lambda \alpha )}\underset{\eta }{𝑝}{}_{(\lambda \alpha )}{}^{}\underset{\eta }{\Psi }=i\gamma \underset{\eta }{}\underset{\eta }{\Psi }=0,\gamma \underset{Q}{𝑝}\underset{Q}{\Psi }=i\gamma \underset{Q}{}\underset{Q}{\Psi }=0,\gamma \underset{W}{𝑝}\underset{W}{\Psi }=i\gamma \underset{W}{}\underset{W}{\Psi }=0.$$
(6.6.8)
The important feature is that the field equations (6.6.7) remain invariant under the substitution $`\underset{Q}{\Psi }^{(0)}\underset{Q}{\Psi }^{(m)}`$, where $`\underset{Q}{\Psi }^{(0)}`$ and $`\underset{Q}{\Psi }^{(m)}`$ are respectively the massless and massive $`Q`$-component fields, to which merely the substitution $`\underset{\eta }{\Psi }^{(0)}\underset{\eta }{\Psi }^{(m)}`$ is corresponded. In free state the massless field components $`\underset{\eta }{\Psi }`$, $`\underset{Q}{\Psi }`$ and $`\underset{W}{\Psi }`$ are independent and due to eq.(6.6.8), the Lagrangian
$$L^{}{}_{0}{}^{(0)}=\underset{u}{\Psi }^+\underset{\eta }{𝐿}{}_{0}{}^{(0)}\underset{u}{\Psi }\underset{\eta }{\Psi }{}_{}{}^{+}\underset{u}{𝐿}_{0}^{(0)}\underset{\eta }{\Psi }=\underset{u}{\Psi }^+\underset{\eta }{𝐿}{}_{0}{}^{(0)}\underset{u}{\Psi }\underset{\eta }{\Psi }^+(\underset{Q}{𝐿}{}_{0}{}^{(0)}+\underset{W}{𝐿}{}_{0}{}^{(0)})\underset{\eta }{\Psi }$$
(6.6.9)
reduces to the following:
$$L^{}{}_{0}{}^{(0)}=\underset{\eta }{𝐿}{}_{0}{}^{(0)}\underset{u}{𝐿}{}_{0}{}^{(0)}=\underset{\eta }{𝐿}{}_{0}{}^{(0)}\underset{Q}{𝐿}{}_{0}{}^{(0)}\underset{W}{𝐿}{}_{0}{}^{(0)}.$$
(6.6.10)
Hence, we implement our scheme as follows: starting with the reduced Lagrangian $`L^{}_0^{(0)}`$ of free field we shall switch on nonlinear fermion interactions of the components. After a generation of nonzero mass of the $`\underset{Q}{\Psi }`$ component in Q-world (next sec.) we shall look for the corresponding corrections via the eq.(6.6.9) to the reduced Lagrangian eq.(6.6.10) of free field. These corrections mean the interaction between the components governed by the eq.(6.6.7) and eq.(6.6.9), and do not imply at all the mass acquiring process for the $`\eta `$-component (see eq.(6.7.6)).
### 6.7 A Generation of the Fermion Mass in the Q-World
We apply now a well known Nambu-Jona-Lasinio model to generate a fermion mass in the Q-world and start from the chirality invariant total Lagrangian of the field $`\underset{Q}{\Psi }:`$ $`\underset{Q}{𝐿}=\underset{Q}{𝐿}_0^{(0)}\underset{Q}{𝐿}_I,`$ where a primary field $`\underset{Q}{\Psi }`$ is the massless bare spinor implying $`\gamma _5`$ invariance. However, due to interaction the rearrangement of vacuum state has caused a generation of nonzero mass of fermion such like to appearance of energy gap in superconductor \[53-55\]. Actually, one may say that particle physicists have always shown greater openness to adopt creatively concepts of condensed matter physics. In this case as well pursuing the analogy with the BCS-Bogoliubov theory of superconductivity, wherein the energy gap is created by the electron-electron interaction of Cooper pairs, in the it was assumed that the mass of Dirac quasi-particle excitation is due to some interaction between massless bare fermions, which may be considered as a self-consistent (Hartree-Fock) representation of it. This approach based on the main idea that due to a dynamical instability the field theory in general may admit also nontrivial solutions with less symmetry than the initial symmetry of Lagrangian. Hence, it is considered such possibility that the field equations may possess higher symmetry, while their solutions may reflect some asymmetries arisen due to fact that nonperturbative solutions to nonlinear equations do not in general possess the symmetry of the equations themselves. In the solution of massive fermion is obtained which lack the initial $`\gamma _5`$ symmetry of the Lagrangian. On the analogy of Gor’kov’s theory it is shown that if one takes into account only the qualitative dynamical effects connected with rearrangement of vacuum state, in addition to the trivial solution of equation of massless fermion a real Dirac quasi-particle will satisfy the equation with non-zero self-energy operator $`\mathrm{\Sigma }(p,m,\stackrel{~}{\lambda },\mathrm{\Lambda })`$ depending on mass $`(m)`$, coupling constant $`(\stackrel{~}{\lambda })`$ and cut-off $`(\mathrm{\Lambda })`$. In the mean time $`\stackrel{~}{\lambda }=\lambda \mathrm{\Gamma }(m,\stackrel{~}{\lambda },\mathrm{\Lambda })`$, where $`\lambda `$ is a bare coupling, $`\mathrm{\Gamma }`$ is the vertex function. This theory leads to the expression of self-energy operator $`\mathrm{\Sigma }_Q`$ for the field $`\underset{Q}{\Psi }`$. In lowest order it is quadratically divergent, but with a cutoff can be made finite. Making use of passage $`\underset{Q}{𝐺}\underset{Q}{𝑀}_4`$ (subsec.2.1), one shall proceed directly with the calculation. In momentum space one gets
$$\begin{array}{c}\mathrm{\Sigma }_Q=m_Q=\frac{8\lambda i}{(2\pi )^4}\frac{m_Q}{p_Q^2+m_Q^2i\epsilon }F(p_Q,\mathrm{\Lambda })d^4p_Q,\hfill \end{array}$$
(6.7.1)
where $`F(p_Q,\mathrm{\Lambda })`$ is a cutoff factor, $`m_Q=\mathrm{\Delta }_Q`$, $`\mathrm{\Delta }_Q=4\lambda <\underset{Q}{\Psi }_R,\underset{Q}{\Psi }_L^+>`$, $`\underset{Q}{\Psi }_{L,R}={\displaystyle \frac{1\gamma _5}{2}}\underset{Q}{\Psi }`$, $`<\mathrm{}>`$ specifies the physical vacuum averaging. Besides of trivial solution $`m_Q=0`$, this equation has also nontrivial solution determining $`m_Q`$ in terms of $`\lambda `$ and $`\mathrm{\Lambda }`$. Straightforward calculations with invariant cutoff yield the relation $`{\displaystyle \frac{2\pi ^2}{\lambda \mathrm{\Lambda }^2}}=1{\displaystyle \frac{m_Q^2}{\mathrm{\Lambda }^2}}\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{m_Q^2}}+1\right).`$ The latter is valid only if $`{\displaystyle \frac{\lambda \mathrm{\Lambda }^2}{2\pi ^2}}1`$. After a vacuum rearrangement the total Lagrangian of initial massless bare field $`\underset{Q}{\Psi }^0`$ gives rise to corresponding Lagrangian $`\underset{Q}{𝐿}^{(m)}`$ of massive field $`\underset{Q}{\Psi }^{(m)}:`$ $`\underset{Q}{𝐿}=\underset{Q}{𝐿}_0^{(0)}\underset{Q}{𝐿}_I=\underset{Q}{𝐿}^{(m)}`$ describing Dirac particle $`(\gamma p_Q\mathrm{\Sigma }_Q)\underset{Q}{\Psi }^{(m)}=0.`$ In lowest order $`\mathrm{\Sigma }_Q=m_Q\lambda ^{1/2}.`$ Within the refined theory of superconductivity, the collective excitations of quasi-particle pairs arise in addition to the individual quasi-particle excitations when a quasi-particle accelerated in the medium \[55, 60-64\]. This leads to the conclusion given in that, in general, the Dirac quasi-particle is only an approximate description of an entire system with the collective excitations as the stable or unstable bound quasi-particle pairs. In a simple approximation there arise CP-odd excitations of zero mass as well as CP-even massive bound states of nucleon number zero and two. Along the same line we must substitute in eq.(6.6.7) the massless field $`\mathrm{\Psi }^{(0)}\underset{\eta }{\Psi }\underset{Q}{\Psi }^{(0)}\underset{W}{\Psi }`$ by massive field $`\mathrm{\Psi }^{(m)}\underset{\eta }{\Psi }\underset{Q}{\Psi }^{(m)}\underset{W}{\Psi }`$. We obtain
$$\gamma p_Q\underset{Q}{\Psi }^{(m)}=\mathrm{\Sigma }_Q\underset{Q}{\Psi }^{(m)},\gamma p_W\mathrm{\Psi }^{(m)}=0,\gamma p_\eta \mathrm{\Psi }^{(m)}=(\gamma p_Q+\gamma p_W)\mathrm{\Psi }^{(m)}=\mathrm{\Sigma }_Q\mathrm{\Psi }^{(m)}.$$
(6.7.2)
This applies following corrections to eq.(6.6.9):
$$\begin{array}{c}\underset{\eta }{𝐿}^{}{}_{0}{}^{(m)}=\underset{u}{\Psi }{}_{}{}^{+}\underset{\eta }{𝐿}_{0}^{(m)}\underset{u}{\Psi }=\underset{u}{\Psi }^+(\underset{\eta }{𝐿}{}_{0}{}^{(0)}\mathrm{\Sigma }_Q\underset{\eta }{\overline{\mathrm{\Psi }}}\underset{\eta }{\Psi })\underset{u}{\Psi }\underset{\eta }{𝐿}{}_{0}{}^{(0)}\mathrm{\Sigma }_Q\overline{\mathrm{\Psi }}\mathrm{\Psi },\hfill \\ \\ \underset{Q}{𝐿}^{}{}_{0}{}^{(m)}=(\underset{\eta }{\Psi }\underset{W}{\Psi })^+\underset{Q}{𝐿}{}_{0}{}^{(m)}(\underset{\eta }{\Psi }\underset{W}{\Psi })=(\underset{\eta }{\Psi }\underset{W}{\Psi })^+(\underset{Q}{𝐿}{}_{0}{}^{(0)}\mathrm{\Sigma }_Q\underset{Q}{\overline{\mathrm{\Psi }}}\underset{Q}{\Psi })(\underset{\eta }{\Psi }\underset{W}{\Psi })\underset{Q}{𝐿}{}_{0}{}^{(0)}\mathrm{\Sigma }_Q\overline{\mathrm{\Psi }}\mathrm{\Psi },\hfill \end{array}$$
(6.7.3)
where suffix $`(m)`$ in $`\mathrm{\Psi }^{(m)}`$ is left implicit. A redefinition $`\underset{Q}{\Psi }^{(0)}\underset{Q}{\Psi }^{(m)}`$ leaves the structure of the piece of the Lagrangian eq.(12.10) involving only the fields $`\underset{\eta }{\Psi }`$ and $`\underset{W}{\Psi }`$ unchanged
$$\begin{array}{c}L_0=\underset{\eta }{𝐿}{}_{0}{}^{(0)}\underset{W}{𝐿}{}_{0}{}^{(0)}=\left(\underset{\eta }{𝐿}{}_{0}{}^{(0)}\mathrm{\Sigma }_Q\overline{\mathrm{\Psi }}\mathrm{\Psi }\right)\left(\underset{W}{𝐿}{}_{0}{}^{(0)}\mathrm{\Sigma }_Q\overline{\mathrm{\Psi }}\mathrm{\Psi }\right)=\underset{\eta }{𝐿}{}_{0}{}^{(m)}\underset{W}{𝐿}{}_{0}{}^{(m)},\hfill \end{array}$$
(6.7.4)
where the component $`\underset{Q}{\Psi }`$ is left implicit. The gauge invariant Lagrangian eq.(6.6.6) takes the form
$$L(D)=\underset{\eta }{𝐿}(\underset{\eta }{𝐷})\underset{W}{𝐿}(\underset{W}{𝐷}),$$
(6.7.5)
where upon combining and rearranging relevant terms we separate the Lagrangians
$$\underset{\eta }{𝐿}(\underset{\eta }{𝐷})=\frac{i}{2}\{\underset{\eta }{\overline{\mathrm{\Psi }}}\stackrel{}{\gamma \underset{\eta }{𝐷}}\underset{\eta }{\Psi }\underset{\eta }{\overline{\mathrm{\Psi }}}\stackrel{}{\gamma \underset{\eta }{\stackrel{}{D}}}\underset{\eta }{\Psi }\}f_Q\overline{\mathrm{\Psi }}\mathrm{\Psi }\frac{1}{2}Tr(\underset{\eta }{𝐁}\overline{\underset{\eta }{𝐁}})$$
(6.7.6)
$$\underset{W}{𝐿}(\underset{W}{𝐷})=\frac{i}{2}\{\underset{W}{\overline{\mathrm{\Psi }}}\stackrel{}{\gamma \underset{W}{𝐷}}\underset{W}{\Psi }\underset{W}{\overline{\mathrm{\Psi }}}\stackrel{}{\gamma \underset{W}{\stackrel{}{D}}}\underset{W}{\Psi }\}\mathrm{\Sigma }_Q\overline{\mathrm{\Psi }}\mathrm{\Psi }\frac{\lambda }{2}S_WS_{W}^{}{}_{}{}^{+}\frac{1}{2}Tr(\underset{𝐖}{𝐁}\overline{\underset{𝐖}{𝐁}}),$$
(6.7.7)
provided by $`f_Q\mathrm{\Sigma }_Q(p_Q,m_Q,\lambda ,\mathrm{\Lambda }),\mathrm{\Psi }=\underset{\eta }{\Psi }\underset{W}{\Psi }.`$
The eq.(6.7.6) and eq.(6.7.7) are the Lagrangians that will be further evaluated and we shall be concerned within the following.
### 6.8 The Electroweak Interactions: the P-Violation
The microscopic approach creates a particular incentive for the pertinent concepts and ideas of the unified electroweak interactions. We admit that the local rotations in the W-world are occurred at very beginning around two arbitrary axes (subsec.6.2), namely $`DimW_{(2)}^{loc}=N_{(q_1^w,q_2^w)}=2,`$ where N is a subquark number, Dim is a dimension of local rotations. The subquarks come up in doublets forming the basis of fundamental representation of weak isospin group $`SU(2)`$ $`q_{(2)}^w=\left(\begin{array}{c}q_1\\ q_2\end{array}\right)^w.`$ The transformations $`U`$ of local group $`SU^{loc}(2)`$ are implemented upon the left- and right-handed fields $`q_{L,R}(\stackrel{}{T}^w)={\displaystyle \frac{1\gamma _5}{2}}q_{(2)}^w(\stackrel{}{T}^w).`$ If P-symmetry holds, one has $`q_{L,R}^{}=Uq_{L,R},USU^{loc}(2).`$ But, under such circumstances the weak interacting particles could not be realized, because of the condition of the MW connections eq.(6.1.5), which is not satisfied yet, i.e. the Q- and W-worlds could not be realized separately. A simple way of effecting a reconciliation is to assume that during a realization of weak interacting charged fermions, under the action of the Q-world, instead the spanning of the initial world $`W_{(2)}^{loc}`$ into the world of unified electroweak interaction $`W_{(3)}^{loc}`$ took place, where the local rotations always occur around all the three axes: $`W_{(2)}^{loc}W_{(3)}^{loc}`$ provided by $`DimW_{(3)}^{loc}=3N_{(q_1^w,q_2^w)}=2.`$ As far as at the very beginning all the subquark fields in W-world are massless, we cannot rule out the possibility that they are transformed independently. On the other hand, when this situation prevails the spanning $`W_{(2)}^{loc}W_{(3)}^{loc}`$ must be occurred compulsory in order to provide a necessary background for the condition eq.(6.1.5) to be satisfied. The most likely attitude here is that doing away this shortage the subquark fields $`q_{L_1},q_{L_2},q_{R_1},`$ and $`q_{R_2}`$ tend to give rise to triplet. The three dimensional effective space $`W_{(3)}^{loc}`$ will then arise
$$W_{(2)}^{loc}q_{(2)}^w(\stackrel{}{T}^w=\frac{1}{2})q_{(3)}^w=\left(\begin{array}{c}q_R(\stackrel{}{T}^w=0)\\ \\ q_L(\stackrel{}{T}^w=\frac{1}{2})\end{array}\right)=\left(\begin{array}{c}q_3^w\\ q_1^w\\ q_2^w\end{array}\right)\left(\begin{array}{c}q_{R_2}\\ q_{L_1}\\ q_{L_2}\end{array}\right)W_{(3)}^{loc}.$$
(6.8.1)
The latter holds if violating initial P-symmetry the components $`q_{R_1},q_{R_2}`$ have remained in the isosinglet states, i.e. the components $`q_L`$ form isodoublet while $`q_R`$ is a isosinglet: $`q_L(\stackrel{}{T}^w={\displaystyle \frac{1}{2}}),q_R(\stackrel{}{T}^w=0).`$ Hence, the mirror symmetry is broken. Corresponding local transformations are implemented upon triplet $`q_{}^{w}{}_{(3)}{}^{}=f_W^{(3)}q_{(3)}^w,`$ where the unitary matrix of three dimensional local rotations reads
$$f_W^{(3)}=\left(\begin{array}{ccc}f_{33}& 0& 0\\ 0& f_{11}& f_{12}\\ 0& f_{21}& f_{22}\end{array}\right)^w.$$
Making use of incompatibility relations eq.(4.4.5) one gets
$$f_W^{(3)}=f_{33}(f_{11}f_{22}f_{12}f_{21})=f_{33}\epsilon _{123}\epsilon _{123}f_W^{(3)}f_{33}^{},$$
(6.8.2)
or $`f_{33}f_{33}^{}=1.`$ That is $`f_{33}=e^{i\beta },`$ and
$$f_W^{(2)}=f_{11}f_{22}f_{12}f_{21}=f_W^{(3)}f_{33}^{}=f_W^{(3)}e^{i\beta }.$$
Due to condition $`f_W^{(3)}=1`$ it reads $`f_W^{(2)}=e^{i\beta }1,`$ thus, the initial symmetry $`SU^{loc}(2)`$ is broken. Restoring it the fields $`q_L`$ should be undergone to additional transformations
$$f_W^{(2)}f_{W}^{(2)}{}_{}{}^{}=\left(\begin{array}{cc}f_{11}e^{i\frac{\beta }{2}}& f_{12}e^{i\frac{\beta }{2}}\\ f_{21}e^{i\frac{\beta }{2}}& f_{22}e^{i\frac{\beta }{2}}\end{array}\right)^w$$
in order to satisfy the unimodularity condition of the matrix of the group $`SU^{loc}(2)`$: $`f_{W}^{(2)}{}_{}{}^{}=f_W^{(2)}e^{i\beta }=1,f_{W}^{(2)}{}_{}{}^{}SU^{loc}(2).`$ While, the expanded group of local rotations in W-world has arisen
$$f_{exp}^{(3)}=\left(\begin{array}{ccc}e^{i\beta }& 0& 0\\ 0& f_{11}e^{i\frac{\beta }{2}}& f_{12}e^{i\frac{\beta }{2}}\\ 0& f_{21}e^{i\frac{\beta }{2}}& f_{22}e^{i\frac{\beta }{2}}\end{array}\right)SU^{loc}(2)_LU^{loc}(1),$$
(6.8.3)
where $`U=e^{i\stackrel{}{T}^w\stackrel{}{\theta ^w}}SU^{loc}(2)_L,U_1=e^{iY^w\theta _1}U^{loc}(1).`$ Here $`U^{loc}(1)`$ is the group of weak hypercharge $`Y^w`$ taking the following values for left- and right-handed subquark fields: $`q_R:Y^w=0,2,q_L:Y^w=1.`$ Whence $`q_{(3)}^w=f_{exp}^{(3)}q_{(3)}^w,`$ and
$$q_L^{}=e^{i\stackrel{}{T}^w\stackrel{}{\theta }^wiY_L^w\theta _1}q_L,q_R^{}=e^{iY_R^w\theta _1}q_R.$$
### 6.9 The Reduction Coefficient and the Weinberg Mixing Angle
The realization of weak interacting particles has always incorporated with the spanning eq.(6.8.1). This implies P-violation in W-world expressed in the reduction of initial symmetry group of local transformations of right-handed components $`q_R`$:
$$\left[SU(2)\right]_R\left[U(1)\right]_R,$$
(6.9.1)
The invariance of physical system of the fields $`q_R`$ under initial group $`\left[SU(2)\right]_R`$ may be realized as well by introducing non-Abelian massless vector gauge fields $`𝐀=\stackrel{}{T}^w\stackrel{}{A}`$ with the values in Lie algebra of the group $`\left[SU(2)\right]_R`$. Under the reduction eq.(6.9.1) the coupling constant $`(g)`$ changed into $`(g^{})`$ specifying the interaction strength between $`q_R`$ and the Abelian gauge field $`B`$ associated with the local group $`\left[U(1)\right]_R`$. While $`g=g^{}\mathrm{tan}\theta _w,`$ where $`\theta _w`$ is the Weinberg mixing angle, in terms of which the reduction coefficient reads $`r_p={\displaystyle \frac{gg^{}}{g+g^{}}}={\displaystyle \frac{1\mathrm{tan}\theta _w}{1+\mathrm{tan}\theta _w}}.`$ To define the $`r_p`$ we consider the interaction vertices corresponding to the groups $`\left[SU(2)\right]_R:`$ $`g𝐀\overline{q}_R\gamma {\displaystyle \frac{\tau }{2}}q_R`$ and $`\left[U(1)\right]_R:`$ $`g^{}B\overline{q}_R\gamma {\displaystyle \frac{Y^w}{2}}q_R.`$ Notifying that the matrix $`{\displaystyle \frac{\lambda _8}{2}}`$ is in the same normalization scale as each of the matrices $`{\displaystyle \frac{\lambda _i}{2}}(i=1,2,3):`$ $`Tr\left({\displaystyle \frac{\lambda _8}{2}}\right)^2=Tr\left({\displaystyle \frac{\lambda _i}{2}}\right)^2={\displaystyle \frac{1}{2}}`$ the vertex scale reads $`(\text{Scale})_{SU(2)}=g{\displaystyle \frac{\lambda _3}{2}},`$ which is equivalent to $`g{\displaystyle \frac{\lambda _8}{2}}`$. It is obvious that per generator scale should not be changed at the reduction eq.(6.9.1), i.e. $`{\displaystyle \frac{(\text{Scale})_{SU(2)}}{N_{SU(2)}}}={\displaystyle \frac{(\text{Scale})_{U(1)}}{N_{U(1)}}},`$ where $`N_{SU(2)}`$ and $`N_{U(1)}`$ are the numbers of generators respectively in the groups $`SU(2)`$ and $`U(1)`$. Hence $`(Scale)_{U(1)}={\displaystyle \frac{1}{3}}(Scale)_{SU(2)}.`$ Stated somewhat differently, the normalized vertex for the group $`\left[U(1)\right]_R`$ reads $`{\displaystyle \frac{1}{3}}gB\overline{q}_R\gamma {\displaystyle \frac{\lambda _8}{2}}q_R.`$ In comparing the coefficients can then be equated $`{\displaystyle \frac{g^{}}{g}}=\mathrm{tan}\theta _w={\displaystyle \frac{1}{\sqrt{3}}},`$ and $`r_p0.27.`$ We may draw a statement that during the realization of the MW-structure the spanning eq.(6.8.1) compulsory occurred, which underlies the P-violation in W-world incorporated with the reduction eq.(6.9.1). The latter is characterized by the Weinberg mixing angle with the value fixed at $`30^0`$.
### 6.10 Emergence of Composite Isospinor-Scalar Bosons
The field $`q_{(2)}^w`$ is the W-component of total field $`q_{(2)}=\underset{\eta }{𝑞}{}_{(2)}{}^{}\underset{W}{𝑞}{}_{(2)}{}^{}(\underset{\eta }{\Psi }{}_{(2)}{}^{}\underset{W}{\Psi }{}_{(2)}{}^{}),`$ where the field component $`\underset{Q}{𝑞}(\underset{Q}{\Psi })`$ is left implicit. Instead of it, below we introduce the additional suffix $`(Q=0,\pm )`$ specifying electric charge of the field. At the very beginning there is an absolute symmetry between the components $`q_1=\underset{\eta }{𝑞}{}_{1}{}^{}\underset{W}{𝑞}_1`$ and $`q_2=\underset{\eta }{𝑞}{}_{2}{}^{}\underset{W}{𝑞}{}_{2}{}^{}.`$ Hence, left- and right-handed components of fields may be written
$$q_{1L}=\underset{\eta }{𝑞}{}_{1L}{}^{(0)}\underset{W}{𝑞}{}_{1L}{}^{()},q_{2L}=\underset{\eta }{𝑞}{}_{2L}{}^{()}\underset{W}{𝑞}{}_{2L}{}^{(0)},q_{1R}=\underset{\eta }{𝑞}{}_{1R}{}^{(0)}\underset{W}{𝑞}{}_{1R}{}^{()},q_{2R}=\underset{\eta }{𝑞}{}_{2R}{}^{()}\underset{W}{𝑞}{}_{2R}{}^{(0)}.$$
(6.10.1)
On the example of one lepton generation $`e`$ and $`\nu `$, without loss of generality, we shall exploit the properties of these fields. A further implication of other fermion generations will be straightforward. One has
$$\begin{array}{c}\underset{\eta }{𝑞}{}_{L}{}^{}=\left(\begin{array}{c}\underset{\eta }{𝑞}_{1L}^{(0)}\\ \\ \underset{\eta }{𝑞}_{2L}^{()}\end{array}\right)L=\left(\begin{array}{c}\nu _L\\ e_L^{}\end{array}\right),\underset{\eta }{𝑞}{}_{R}{}^{}=(\underset{\eta }{𝑞}{}_{1R}{}^{(0)},\underset{\eta }{𝑞}{}_{2R}{}^{()})R=(\nu _R,e_R^{}),\hfill \\ \\ \underset{W}{𝑞}_L=\left(\begin{array}{c}\underset{W}{𝑞}_{1L}^{()}\\ \\ \underset{W}{𝑞}_{2L}^{(0)}\end{array}\right),\underset{W}{𝑞}{}_{R}{}^{}=(\underset{W}{𝑞}{}_{1R}{}^{()},\underset{W}{𝑞}{}_{2R}{}^{(0)}).\hfill \end{array}$$
We evaluate the term $`f_Q\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ in the Lagrangian eq.(6.7.6)
$$\overline{\mathrm{\Psi }}\mathrm{\Psi }=\mathrm{\Psi }_L^+\mathrm{\Psi }_R+\mathrm{\Psi }_R^+\mathrm{\Psi }_Lq_L^+q_R+q_R^+q_L,$$
provided by
$$\begin{array}{c}q_L^+q_R=\underset{\eta }{𝑞}{}_{L}{}^{+}\underset{W}{𝑞}{}_{L}{}^{+}\underset{W}{𝑞}{}_{R}{}^{}\underset{\eta }{𝑞}{}_{R}{}^{}=\overline{\underset{\eta }{𝑞}}{}_{L}{}^{}(\gamma ^0\underset{W}{𝑞}{}_{L}{}^{+}\underset{W}{𝑞}{}_{R}{}^{})\underset{\eta }{𝑞}{}_{R}{}^{}=\overline{L}\phi R,\hfill \\ \\ q_R^+q_L=\underset{\eta }{𝑞}{}_{R}{}^{+}\underset{W}{𝑞}{}_{R}{}^{+}\underset{W}{𝑞}{}_{L}{}^{}\underset{\eta }{𝑞}{}_{L}{}^{}=\overline{\underset{\eta }{𝑞}}{}_{R}{}^{}(\gamma ^0\underset{W}{𝑞}{}_{R}{}^{+}\underset{W}{𝑞}{}_{L}{}^{})\underset{\eta }{𝑞}{}_{L}{}^{}=\overline{R}\phi ^+L.\hfill \end{array}$$
For appropriate values of the parameters this term causes
$$\overline{\mathrm{\Psi }}\mathrm{\Psi }=\overline{q}_{(2)}q_{(2)}=\overline{L}\phi R+\overline{R}\phi ^+L,$$
where the isospinor-scalar meson field $`\phi `$ reads
$$\phi \gamma ^0\underset{W}{𝑞}{}_{L}{}^{+}\underset{W}{𝑞}{}_{R}{}^{},\phi ^+\gamma ^0\underset{W}{𝑞}{}_{R}{}^{+}\underset{W}{𝑞}{}_{L}{}^{}.$$
A calculation gives
$$\phi =\left(\begin{array}{c}\phi _1\\ \phi _2\end{array}\right),\phi _1(\underset{W}{𝑞}{}_{1L}{}^{()})^+\underset{W}{𝑞}{}_{R}{}^{},\phi _2(\underset{W}{𝑞}{}_{2L}{}^{(0)})^+\underset{W}{𝑞}{}_{R}{}^{},\phi ^+=(\phi _1^+,\phi _2^+).$$
Hence, the possible two doublets of the composite isospinor-scalar bosons read
$$\phi _u=\left(\begin{array}{c}\phi _{1u}^{(+)}\\ \\ \phi _{2u}^{(0)}\end{array}\right),\phi _d=\left(\begin{array}{c}\phi _{1d}^{(0)}\\ \\ \phi _{2d}^{()}\end{array}\right).$$
where
$$\phi _{1u}^+(\underset{W}{𝑞}{}_{1L}{}^{()})^+\underset{W}{𝑞}{}_{2R}{}^{(0)},\phi _{2u}^0(\underset{W}{𝑞}{}_{2L}{}^{(0)})^+\underset{W}{𝑞}{}_{2R}{}^{(0)},$$
and
$$\phi _{1d}^0(\underset{W}{𝑞}{}_{1L}{}^{()})^+\underset{W}{𝑞}{}_{1R}{}^{()},\phi _{2d}^{}(\underset{W}{𝑞}{}_{2L}{}^{(0)})^+\underset{W}{𝑞}{}_{1R}{}^{()}.$$
In accordance with eq.(6.1.5), the isospinor-scalar meson carries following weak hypercharge $`\phi :Y^w=1.`$ Thus, the term $`f_Q\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ arisen in the total Lagrangian of fundamental fermion field eq.(6.7.6) accommodates the Yukawa couplings between the fermions and corresponding isospinor-scalar bosons in fairy conventional form
$$f_Q\overline{\mathrm{\Psi }}\mathrm{\Psi }=f_e\left(\overline{L}\phi e_R+\overline{e_R}\phi ^+L\right)f_\nu \left(\overline{L}\phi _c\nu _R+\overline{\nu _R}\phi _c^+L\right),$$
(6.10.2)
where the charge conjugated field $`\phi _c`$ is defined $`\left(\phi _c\right)_i=\phi ^k\epsilon _{ik}`$. To compute the coupling constants $`f_e`$ and $`f_\nu `$ for the leptons one must retrieve their implicit field-components $`\underset{Q}{\Psi }`$. Hence
$$f_i=tr(\rho _i^Q\mathrm{\Sigma }_Q),f_i^\nu =tr(\rho _i^{Q\nu }\mathrm{\Sigma }_Q),$$
(6.10.3)
where the density operators $`\rho _i^Q`$ and $`\rho _i^{Q\nu }`$ for given $`i`$ of the pure ensembles are used
$$\begin{array}{c}\rho _i^Q=\left(q_1q_2q_3\right)_i^{Q+}\left(q_1q_2q_3\right)_i^Q,\rho _i^{Q\nu }=\left(q_i\overline{q}_i\right)^{Q+}\left(q_i\overline{q}_i\right)^Q,\hfill \\ \\ tr(\rho _i^Q)^2=tr(\rho _i^Q)=1,tr(\rho _i^{Q\nu })^2=tr(\rho _i^{Q\nu })=1.\hfill \end{array}$$
(6.10.4)
According to eq.(6.1.6), one gets
$$f_i=L_i^2\overline{\mathrm{\Sigma }}_Q,f_i^\nu =L_i^2\overline{\mathrm{\Sigma }}_Q^\nu ,\overline{\mathrm{\Sigma }}_Q\mathrm{\Sigma }_Q(\lambda ,L)\rho ^Q,\overline{\mathrm{\Sigma }}_Q{}_{}{}^{\nu }\mathrm{\Sigma }_Q(\lambda ,L)\rho ^{Q\nu },$$
(6.10.5)
where
$$\rho ^Q=\left(q_1q_2q_3\right)^{Q+}\left(q_1q_2q_3\right)^Q,\rho ^{Q\nu }=\left(q\overline{q}\right)^{Q+}\left(q\overline{q}\right)^Q.$$
(6.10.6)
An implication of the quarks into this scheme is straightforward if one retrieves their implicit field-components $`\underset{Q}{\Psi },\underset{B}{\Psi },\underset{C}{\Psi },(C=s,c,b,t)`$ (subsec.6.6). On the analogy of previous case the coupling constants read
$$f_i=tr(\rho _i\mathrm{\Sigma }_Q),$$
(6.10.7)
where $`i=u,d,s,c,b,t.`$ Taking into account the MW structure of the quarks eq.(6.4.1), we may write down the corresponding density operators
$$\rho _i=\rho _i^Q\rho _i^B\rho _i^C$$
(6.10.8)
given in a convention
$$\rho _i^A=\underset{A}{\Psi }{}_{i}{}^{+}\underset{A}{\Psi }{}_{i}{}^{},$$
(6.10.9)
where $`\rho _u^C=\rho _d^C=1.`$
## 7 The Higgs Boson
To break the gauge symmetry down and leading to masses of the fields , one needs in general, several kinds of spinless Higgs bosons , with conventional Yukawa couplings to fermion currents and transforming by an irreducible representation of gauge group. The Higgs theory like \[76-79\] involves these bosons as the ready made fundamental elementary fields, which entails various difficulties. Within outlined here microscopic approach the self-interacting isospinor-scalar Higgs bosons arise in the $`W`$-world as the collective modes of excitations of the bound quasi-particle iso-pairs.
### 7.1 The Bose Condensate of Iso-Pairs
The ferromagnetism , Bose superfluid and BCS-Bogoliubov model of superconductivity \[53-55\] are characterized by the condensation phenomenon leading to the symmetry-breaking ground state. It is particularly helpful to remember that in BCS-Bogoliubov theory the importance of this phenomenon resides in the possibility suggested by Cooper that in the case of an arbitrary weak interaction the pair, composed of two mutually interacting electrons above the quiescent Fermi sea, remains in a bound state. The electrons filling the Fermi sea do not interact with the pair and in the same time they block the levels below the Fermi surface. The superconductive phase arises due to effective attraction between electrons occurred by exchange of virtual phonons . In BCS microscopic theory of superconductivity instead of bound states, with inception by Cooper, one has a state with strongly correlated electron pairs or condensed state in which the pairs form the condensate. The energy of a system in the superconducting state is smaller than the energy in the normal state described by the Bloch individual-particle model. The energy gap arises is due to existence of the binding energy of a pair as a collective effect, the width of which is equal to twice the binding energy. According to Pauli exclusion principle, only the electrons situated in the spherical thin shell near the Fermi surface can form bound pairs in which they have opposite spin and momentum. The binding energy is maximum at absolute zero and decreases along the temperature increasing because of the disintegration of pairs. Pursuing the analogy with these ideas in outlined here approach a serious problem is to find out the eligible mechanism leading to the formation of pairs, somewhat like Cooper mechanism, but generalized for relativistic fermions, of course in absence of any lattice. We suggest this mechanism in the framework of gauge invariance incorporated with the P-violation phenomenon in W-world. To trace a maximum resemblance to the superconductivity theory, within this section it will be convenient to describe our approach in terms of four dimensional Minkowski space $`\underset{W}{𝑀}_4`$ corresponding to the internal W-world: $`\underset{W}{𝐺}\underset{W}{𝑀}_4`$(subsec.2.1). Although we shall leave the suffix $`(W)`$ implicit, but it goes without saying that all results obtained within this section refer to the W-world. According to previous section, we consider the isospinor-scalar $`\phi `$-meson arisen in the W-world
$$\phi (x)=\gamma ^0\mathrm{\Psi }_{L}^{}{}_{}{}^{+}(x)\mathrm{\Psi }_R(x),$$
where $`xM_4`$ is a point of the W-world. The following notational conventions will be employed throughout
$$\underset{W}{𝑞}{}_{L}{}^{}\underset{W}{\Psi }{}_{L}{}^{}(\underset{W}{𝑥})\mathrm{\Psi }_L(x),\underset{W}{𝑀}_4M_4,\underset{W}{𝑞}{}_{R}{}^{}\underset{W}{\Psi }{}_{R}{}^{}(\underset{W}{𝑥})\mathrm{\Psi }_R(x),$$
where $`\mathrm{\Psi }_R(x)=\gamma (1+\stackrel{}{\sigma }\stackrel{}{\beta })\mathrm{\Psi }_L(x),\stackrel{}{\beta }={\displaystyle \frac{\stackrel{}{v}}{c}},\mathrm{\Psi }_L(x)=\gamma (1\stackrel{}{\sigma }\stackrel{}{\beta })\mathrm{\Psi }_R(x),\gamma ={\displaystyle \frac{E}{m}},`$ provided by the spin $`\stackrel{}{\sigma }`$, energy $`E`$ and velocity $`\stackrel{}{v}`$ of particle. In terms of Fourier integrals
$$\mathrm{\Psi }_L(x)=\frac{1}{(2\pi )^4}\mathrm{\Psi }_L(p_L)e^{ip_Lx}d^4p_L,\mathrm{\Psi }_R(x)=\frac{1}{(2\pi )^4}\mathrm{\Psi }_R(p_R)e^{ip_Rx}d^4p_R,$$
(7.1.1)
it is readily to get
$$\phi (k)=\phi (x)e^{ikx}d^4x=\gamma ^0\frac{d^4p_L}{(2\pi )^4}\mathrm{\Psi }_{L}^{}{}_{}{}^{+}(p_L)\mathrm{\Psi }_R(p_L+k)=\gamma ^0\frac{d^4p_R}{(2\pi )^4}\mathrm{\Psi }_{L}^{}{}_{}{}^{+}(p_Rk)\mathrm{\Psi }_R(p_R)$$
(7.1.2)
provided by the conservation law of fourmomentum $`k=p_Rp_L,`$ where $`k=k(\omega ,\stackrel{}{k})`$, $`p_{L,R}=p_{L,R}(E_{L,R},\stackrel{}{p}_{L,R})`$. Our arguments on Bose-condensation are based on the local gauge invariance of the theory incorporated with the P-violation in weak interactions. The rationale for this approach is readily forthcoming from the consideration of gauge transformations of the fields eq.(7.1.2) under the P-violation in the W-world
$$\mathrm{\Psi }_L^{}(x)=U_L(x)\mathrm{\Psi }_L(x),\mathrm{\Psi }_R^{}(x)=U_R(x)\mathrm{\Psi }_R(x),$$
where the Fourier expansions carried out over corresponding gauge quanta with wave fourvectors $`q_L`$ and $`q_R`$
$$U_L(x)=\frac{d^4q_L}{(2\pi )^4}e^{iq_Lx}U_L(q_L),U_R(x)=\frac{d^4q_R}{(2\pi )^4}e^{iq_Rx}U_L(q_R),$$
(7.1.3)
and $`U_L(x)U_R(x).`$ They induce the gauge transformations implemented upon the $`\phi `$-field $`\phi ^{}(x)=U(x)\phi (x).`$ The matrix of induced gauge transformations may be written down in terms of induced gauge quanta
$$U(x)U_L^+(x)U_R(x)=\frac{d^4q}{(2\pi )^4}e^{iqx}U(q),$$
(7.1.4)
where $`q=q_L+q_R,q(q^0,\stackrel{}{q})`$. In momentum space one gets
$$\begin{array}{c}\phi ^{}(k^{})=\frac{d^4q}{(2\pi )^4}U(q)\phi (k^{}q)=\frac{d^4k}{(2\pi )^4}U(k^{}k)\phi (k).\hfill \end{array}$$
(7.1.5)
Conservation of the fourmomentum requires that $`k^{}=k+q.`$ According to eq.(7.1.2) and eq.(7.1.5), we have
$$p_L^{}+p_R^{}=p_L+p_R+q=p_L^{\prime \prime }+p_R=p_L+p_R^{\prime \prime },$$
where $`p_L^{\prime \prime }=p_Lq,p_R^{\prime \prime }=p_R+q.`$ Whence the wave vectors of fermions imply the conservation law $`\stackrel{}{p}_L+\stackrel{}{p}_R=\stackrel{}{p}_L^{\prime \prime }+\stackrel{}{p}_R^{\prime \prime },`$ characterizing the scattering process of two fermions with effective interaction caused by the mediating induced gauge quanta. We suggest the mechanism for the effective attraction between the fermions in the following manner: Among all induced gauge transformations with miscellaneous gauge quanta we distinguish a special subset with the induced gauge quanta of the frequencies belonged to finite region characterized by the maximum frequency $`{\displaystyle \frac{\stackrel{~}{q}}{\mathrm{}}}(\stackrel{~}{q}=max\{q^0\})`$ greater than the frequency of inducing oscillations fermion force $`{\displaystyle \frac{\overline{E_L}\overline{E_L^{\prime \prime }}}{\mathrm{}}}<{\displaystyle \frac{\stackrel{~}{q}}{\mathrm{}}}.`$ To the extent that this is a general phenomenon, we can expect under this condition the effective attraction (negative interaction) arisen between the fermions caused by exchange of virtual induced gauge quanta if only the forced oscillations of these quanta occur in the same phase with the oscillations of inducing force (the oscillations of fermion density). In view of this we may think of isospinor $`\mathrm{\Psi }_L`$ and isoscalar $`\mathrm{\Psi }_R`$ fields as the fermion fields composing the iso-pairs with the same conserving net momentum $`\stackrel{}{p}=\stackrel{}{p}_L+\stackrel{}{p}_R`$ and opposite spin, for which the maximum number of negative matrix elements of operators composed by corresponding creation and annihilation operators $`a_{\stackrel{}{p}_R^{\prime \prime }}^+a_{\stackrel{}{p}_R}a_{\stackrel{}{p}_L^{\prime \prime }}^+a_{\stackrel{}{p}_L}`$ (designated by the pair wave vector $`\stackrel{}{p}`$) may be obtained for coherent ground state with $`\stackrel{}{p}=\stackrel{}{p}_L+\stackrel{}{p}_R=0.`$ In the mean time the interaction potential reads
$$V=\underset{\stackrel{}{p}_R^{\prime \prime },\stackrel{}{p}_L^{\prime \prime },\stackrel{}{p}_R,\stackrel{}{p}_L}{}\left(a_{\stackrel{}{p}_L}^+\right)^+a_{\stackrel{}{p}_R^{\prime \prime }}^+\left(a_{\stackrel{}{p}_L^{\prime \prime }}^+\right)a_{\stackrel{}{p}_R}=\underset{\stackrel{}{p}_R^{\prime \prime },\stackrel{}{p}_L^{\prime \prime },\stackrel{}{p}_R,\stackrel{}{p}_L}{}a_{\stackrel{}{p}_R^{\prime \prime }}^+a_{\stackrel{}{p}_L^{\prime \prime }}^+a_{\stackrel{}{p}_R}a_{\stackrel{}{p}_L},$$
(7.1.6)
implying the attraction between the fermions situated in the spherical thin shall near the Fermi surface
$$V_{\stackrel{}{p}\stackrel{}{p}^{\prime \prime }}=\{\begin{array}{cc}V\text{at}E_\stackrel{}{p}E_F\stackrel{~}{q},E_{\stackrel{}{p}^{\prime \prime }}E_F\stackrel{~}{q},\hfill & \\ 0\text{otherwise}\hfill & \end{array}.$$
(7.1.7)
The fermions filled up the Fermi sea block the levels below Fermi surface. Hence, the fermions are in superconducting state if the condition eq.(7.1.7) holds. Otherwise, they are in normal state described by Bloch individual particle model. Hence, the Bose-condensate arises in the W-world as the collective mode of excitations of bound quasi-particle iso-pairs described by the same wave function in the superconducting phase $`\mathrm{\Psi }=<\mathrm{\Psi }_L\mathrm{\Psi }_R>,`$ where $`<\mathrm{}>`$ is taken to denote the vacuum averaging. The vacuum of the W-world is filled up by such iso-pairs at absolute zero $`T=0`$.
We make a final observation that $`\mathrm{\Psi }_R\mathrm{\Psi }_R^+=n_R`$ is a scalar density number of right-handed particles. It readily follows that:
$$(\mathrm{\Psi }_L\mathrm{\Psi }_R)^+(\mathrm{\Psi }_L\mathrm{\Psi }_R)=\mathrm{\Psi }_R^+\mathrm{\Psi }_L^+\mathrm{\Psi }_L\mathrm{\Psi }_R=\frac{1}{n_R}\mathrm{\Psi }_R^+\gamma ^0(\gamma ^0\mathrm{\Psi }_L^+\mathrm{\Psi }_R)(\mathrm{\Psi }_R^+\mathrm{\Psi }_L\gamma ^0)\gamma ^0\mathrm{\Psi }_R=\phi \phi ^+,$$
(7.1.8)
where $`\mathrm{\Psi }^2=<\phi \phi ^+>=<\phi >^2\phi ^2.`$ It is convenient to abbreviate the $`<\phi >`$ by the symbol $`\phi `$. The eq.(7.1.8) indeed shows that the $`\phi `$-meson actually arises as the collective mode of excitations of bound quasi-particle iso-pairs.
### 7.2 The Non-Relativistic Approximation
In the approximation to non-relativistic limit $`(\beta 1,\mathrm{\Psi }_L\mathrm{\Psi }_R,\gamma ^01)`$ by making use of Ginzburg-Landau’s (GL) phenomenological theory it is straightforward to write down the free-energy functional for the order parameter in equilibrium superconducting phase in presence of magnetic field. The self-consistent coupled GL-equations are differential equations like Schrödinger and Maxwell equations, which relate the spatial variation of the order parameter $`\mathrm{\Psi }`$ to the vector potential $`\stackrel{}{A}`$ and the current $`\stackrel{}{j}_s`$. In the papers , by means of thermodynamic Green’s functions in well defined limit Gor’kov was able to show that GL-equations are a consequence of the BCS-Bogoliubov microscopic theory of superconductivity. The theoretical significance of these works resides in the microscopic interpretation of all physical parameters of GL-theory. Subsequently these ideas were extended to lower temperatures by others \[85-87\] using a requirement that the order parameter and vector potential vary slowly over distances of the order of the coherence length and that the electrodynamics be local (London limit). Namely, the validity of derived GLG-equations is restricted to the temperature $`T`$, such $`T_cTT_c`$ and to the local electrodynamics region $`q\xi _01`$, where $`T_c`$ is transition temperature, $`\xi _0`$ is coherent length characterizing the spatial extent of the electron pair correlations, $`q`$ are the wave numbers of magnetic field $`\stackrel{}{A}`$. The most important order parameter $`\mathrm{\Psi }`$, the mass $`m_\mathrm{\Psi }`$ and the coupling constant $`\lambda _\mathrm{\Psi }`$ figured in GLG-equations read
$$\begin{array}{c}\mathrm{\Psi }(\stackrel{}{r})=\frac{(7\zeta (3)N)^{1/2}}{4\pi k_BT_c}\mathrm{\Delta }(\stackrel{}{r}),\mathrm{\Delta }(T)=3.1k_BT_c(1\frac{T}{T_c})^{1/2},\xi _0=0.18\frac{\mathrm{}v_F}{k_BT_c},\hfill \\ \\ m_\mathrm{\Psi }^2=1.83\frac{\mathrm{}^2}{m}\frac{1}{\xi _0^2}\left(1\frac{T}{T_c}\right),\lambda _\mathrm{\Psi }^2=1.4\frac{1}{N(0)}\left(\frac{\mathrm{}^2}{2m\xi _0^2}\right)^2\frac{1}{(k_BT_c)^2}.\hfill \end{array}$$
(7.2.1)
Reviewing the notation $`\mathrm{\Delta }(\stackrel{}{r})`$ is the energy gap, $`e^{}=2e`$ is the effective charge, $`N(0)`$ is the state density at Fermi surface, $`N`$ is the number of particles per unit volume in normal mode, $`v_F`$ is the Fermi velocity, $`m\mathrm{\Sigma }_Q=f_Q`$ is the mass of fermion field. The transition temperature relates to gap at absolute zero $`\mathrm{\Delta }_0`$ . The estimate for the pair size at $`v_F10^8cm/s,T_c1`$ gives $`\xi _010^4cm.`$
### 7.3 The Relativistic Treatment
We start with the Lagrangian eq.(6.7.7) of self-interacting fermion field in W-world, which is arisen from the Lagrangian eq.(6.6.6) of primary fundamental field after the rearrangement of the vacuum of the Q-world
$$\begin{array}{c}\underset{W}{𝐿}(x)=\frac{i}{2}\{\underset{W}{\overline{\mathrm{\Psi }}}(x)\gamma ^\mu \underset{W}{}_\mu \underset{W}{\Psi }(x)\underset{W}{\overline{\mathrm{\Psi }}}(x)\gamma ^\mu \underset{W}{\stackrel{}{}}_\mu \underset{W}{\Psi }(x)\}m\underset{W}{\overline{\mathrm{\Psi }}}(x)\underset{W}{\Psi }(x)\hfill \\ \\ \frac{\lambda }{2}\underset{W}{\overline{\mathrm{\Psi }}}(x)\left(\underset{W}{\overline{\mathrm{\Psi }}}(x)\underset{W}{\Psi }(x)\right)\underset{W}{\Psi }(x).\hfill \end{array}$$
(7.3.1)
Here, $`m=\mathrm{\Sigma }_Q`$ is the self-energy operator of the fermion field component in Q-world, the suffix $`(W)`$ just was put forth in illustration of a point at issue. For the sake of simplicity, we also admit $`\underset{𝐖}{𝐁}(x)=0`$, but, of course, one is free to restore the gauge field $`\underset{𝐖}{𝐁}(x)`$ whenever it will be needed. In lowest order the relation $`mm_Q\lambda ^{1/2}`$ holds. The Lagrangian eq.(7.3.1) leads to the field equations
$$\begin{array}{c}(\gamma pm)\mathrm{\Psi }(x)\lambda \left(\overline{\mathrm{\Psi }}(x)\mathrm{\Psi }(x)\right)\mathrm{\Psi }(x)=0,\hfill \\ \\ \overline{\mathrm{\Psi }}(x)(\gamma \stackrel{}{p}+m)+\lambda \overline{\mathrm{\Psi }}(x)\left(\overline{\mathrm{\Psi }}(x)\mathrm{\Psi }(x)\right)=0,\hfill \end{array}$$
(7.3.2)
where the indices have been suppressed as usual. At non-relativistic limit the function $`\mathrm{\Psi }`$ reads $`\mathrm{\Psi }e^{imc^2t}\mathrm{\Psi },`$ and Lagrangian eq.(7.3.1) leads to Hamiltonian used in . In the following our discussion will be in close analogy with the latter. We make use of the Gor’kov’s technique and evaluate the equations (7.3.2) in following manner: The spirit of the calculation will be to treat interaction between the particles as being absent everywhere except the thin spherical shell $`2\stackrel{~}{q}`$ near the Fermi surface. The Bose condensate of bound particle iso-pairs occurred at zero momentum. The scattering processes between the particles are absent. We consider the matrix elements
$$\begin{array}{c}<T\left(\mathrm{\Psi }_\alpha (x_1)\mathrm{\Psi }_\beta (x_2)\overline{\mathrm{\Psi }}_\gamma (x_3)\overline{\mathrm{\Psi }}_\delta (x_4)\right)>=<T\left(\mathrm{\Psi }_\alpha (x_1)\overline{\mathrm{\Psi }}_\gamma (x_3)\right)><T\left(\mathrm{\Psi }_\beta (x_2)\overline{\mathrm{\Psi }}_\delta (x_4)\right)>\hfill \\ \\ +<T\left(\mathrm{\Psi }_\alpha (x_1)\overline{\mathrm{\Psi }}_\delta (x_4)\right)><T\left(\mathrm{\Psi }_\beta (x_2)\overline{\mathrm{\Psi }}_\gamma (x_3)\right)>+\hfill \\ \\ <T\left(\mathrm{\Psi }_\alpha (x_1)\mathrm{\Psi }_\beta (x_2)\right)><T\left(\overline{\mathrm{\Psi }}_\gamma (x_3)\overline{\mathrm{\Psi }}_\delta (x_4)\right)>,\hfill \end{array}$$
where
$$\begin{array}{c}<T\left(\mathrm{\Psi }_\alpha (x_1)\mathrm{\Psi }_\beta (x_2)\right)><T\left(\overline{\mathrm{\Psi }}_\gamma (x_3)\overline{\mathrm{\Psi }}_\delta (x_4)\right)>=\hfill \\ \\ <T\left(\mathrm{\Psi }_\gamma ^+(x_3)\gamma ^0\mathrm{\Psi }_\delta ^+(x_4)\right)><T\left(\gamma ^0\mathrm{\Psi }_\alpha (x_1)\mathrm{\Psi }_\beta (x_2)\right)>,\hfill \end{array}$$
also introduce the functions
$$\begin{array}{c}<NT\left(\gamma ^0\mathrm{\Psi }(x)\mathrm{\Psi }(x^{})\right)N+2>=e^{2i\mu ^{}t}F(xx^{}),\hfill \\ \\ <N+2T\left(\mathrm{\Psi }^+(x)\gamma ^0\mathrm{\Psi }^+(x^{})\right)N>=e^{2i\mu ^{}t}F^+(xx^{}).\hfill \end{array}$$
(7.3.3)
Here, $`\mu ^{}=\mu +m`$, $`\mu `$ is a chemical potential. We omit a prime over $`\mu `$, but should understand under it $`\mu +m`$. Let us now make use of Fourier integrals
$$G_{\alpha \beta }(xx^{})=\frac{d\omega d\stackrel{}{p}}{(2\pi )^4}G_{\alpha \beta }(p)e^{i\stackrel{}{p}(\stackrel{}{x}\stackrel{}{x^{}})i\omega (tt^{})}$$
etc, which render the equation (7.3.2) easier to handle in momentum space
$$\begin{array}{c}(\gamma pm)G(p)i\lambda \gamma ^0F(0+)\overline{F}(p)=1,\hfill \\ \\ \overline{F}(p)(\gamma p+m2\mu \gamma ^0)+i\lambda \overline{F}(0+)G(p)=0,\hfill \end{array}$$
(7.3.4)
where $`F_{\alpha \beta }(0+)=e^{2i\mu t}<T\left(\gamma ^0\mathrm{\Psi }_\alpha (x)\mathrm{\Psi }_\beta (x)\right)>=\underset{xx^{}(tt^{})}{lim}F_{\alpha \beta }(xx^{}).`$ Next we substitute
$$(\gamma pm)=(\omega ^{}\xi _p)\gamma ^0,(\gamma p+m2\mu \gamma ^0)=\gamma ^0(\omega ^{}+\xi _p^+),$$
where
$$\begin{array}{c}\omega ^{}=\omega \mu ^{}=\omega m\mu ,\xi _p=(\stackrel{}{\gamma }\stackrel{}{p}+m)\gamma ^0\mu ^{}=(\stackrel{}{\gamma }\stackrel{}{p}+m)\gamma ^0m\mu ,\hfill \\ \xi _p^+=\gamma ^0(\stackrel{}{\gamma }^+\stackrel{}{p}+m)m\mu ,\hfill \end{array}$$
and omit a prime over $`\omega ^{}`$ for the rest of this section. We employ
$$F(0+)=JI,I=\left(\begin{array}{c}01\\ 10\end{array}\right),F^+(0+)F(0+)=J^2I^2=J^2,$$
and $`\widehat{\omega }+\widehat{\xi }_p=\gamma ^0(\omega +\xi _p)`$. The gap function $`\mathrm{\Delta }`$ reads $`\mathrm{\Delta }^2=\lambda ^2J^2,`$ where $`J={\displaystyle \frac{d\omega d\stackrel{}{k}}{(2\pi )^4}F^+(p)}.`$ Making use of standard rules , one may pass over the poles. This method allows oneself to extend the study up to limit of temperatures, such that $`T_cTT_c`$, by making use of thermodynamic Green’s function. Hence
$$\begin{array}{c}F^+(p)=i\lambda J(\omega \xi _p+i\delta )^1(\omega +\xi _pi\delta )^1\frac{\pi \mathrm{\Delta }}{\epsilon _p}n(\epsilon _p)\{\delta (\omega \epsilon _p)+\delta (\omega +\epsilon _p)\},\hfill \\ \\ G(p)=\gamma ^0\{u_p^2(\omega \xi _p+i\delta )^1+v_p^2(\omega +\xi _pi\delta )^1+2\pi in(\epsilon _p)[u_p^2\delta (\omega \epsilon _p)\hfill \\ \\ v_p^2\delta (\omega +\epsilon _p)]\},\hfill \end{array}$$
(7.3.5)
where $`u_p^2={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\xi _p}{\epsilon _p}}\right),v_p^2={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\xi _p}{\epsilon _p}}\right),\epsilon _p=(\xi _p^2+\mathrm{\Delta }^2(T))^{1/2}.`$ and $`n(\epsilon _p)`$ is the usual Fermi function $`n(\epsilon _p)=\left(\mathrm{exp}{\displaystyle \frac{\epsilon _p}{T}}+1\right)^1.`$ Then
$$1=\frac{\lambda }{2(2\pi )^3}𝑑\stackrel{}{k}\frac{12n(\epsilon _k)}{\epsilon _k(T)}\left(\xi _p<\stackrel{~}{q}\right),$$
(7.3.6)
determining the energy gap $`\mathrm{\Delta }`$ as a function of $`T`$. According to eq.(17.3.6), the $`\mathrm{\Delta }(T)0`$ at $`TT_c\mathrm{\Delta }(0)`$ .
### 7.4 Self-Interacting Potential of Bose-Condensate
To go any further in exploring the form and significance of obtained results it is entirely feasible to include the generalization of the equations (6.1.3.4) in presence of spatially varying magnetic field with vector potential $`\stackrel{}{A}(\stackrel{}{r})`$, which is straightforward $`(t\tau =it)`$
$$\begin{array}{c}\left\{\gamma ^0\frac{}{\tau }i\stackrel{}{\gamma }\left(\frac{}{\stackrel{}{r}}ie\stackrel{}{A}(\stackrel{}{r})\right)m+\gamma ^0\mu \right\}G(x,x^{})+\gamma ^0\mathrm{\Delta }(\stackrel{}{r})\overline{F}(x,x^{})=\delta (xx^{}),\hfill \\ \\ \overline{F}(x,x^{})\left\{\gamma ^0\frac{}{\tau }+i\stackrel{}{\gamma }\left(\frac{}{\stackrel{}{r}}+ie\stackrel{}{A}(\stackrel{}{r})\right)m+\gamma ^0\mu \right\}\mathrm{\Delta }^{}(\stackrel{}{r})\gamma ^0G(x,x^{})=0,\hfill \end{array}$$
(7.4.1)
where the thermodynamic Green’s function is used and the energy gap function is in the form $`\mathrm{\Delta }^{}(\stackrel{}{r})=\lambda F^+(\tau ,\stackrel{}{r};\tau ,\stackrel{}{r}).`$ This function is logarithmically divergent, but with a cutoff of energy of interacting fermions at the spatial distances in order of $`{\displaystyle \frac{\mathrm{}v}{\stackrel{~}{\omega }}}`$ can be made finite, where $`\stackrel{~}{\omega }{\displaystyle \frac{\stackrel{~}{q}}{\mathrm{}}}`$. If one uses the Fourier components of functions $`G(x,x^{})`$ and $`F(x,x^{})`$
$$G(\stackrel{}{r},\stackrel{}{r}^{};u)=T\underset{n}{}e^{i\omega u}G_\omega (\stackrel{}{r},\stackrel{}{r}^{}),G_\omega (\stackrel{}{r},\stackrel{}{r}^{})=\frac{1}{2}_{1/T}^{1/T}e^{i\omega u}G(\stackrel{}{r},\stackrel{}{r}^{};u)𝑑u,$$
(7.4.2)
where $`u=\tau \tau ^{}`$, $`\omega `$ is the discrete index $`\omega =\pi T(2n+1),n=0,\pm 1,\mathrm{}`$, then the eq.(7.4.1) reduces to
$$\begin{array}{c}\left\{i\omega \gamma ^0i\stackrel{}{\gamma }\left(\stackrel{}{}_\stackrel{}{r}ie\stackrel{}{A}(\stackrel{}{r})\right)m+\gamma ^0\mu \right\}G_\omega (\stackrel{}{r},\stackrel{}{r}^{})+\gamma ^0\mathrm{\Delta }(\stackrel{}{r})\overline{F_\omega }(\stackrel{}{r},\stackrel{}{r}^{})=\delta (\stackrel{}{r}\stackrel{}{r}^{}),\hfill \\ \\ \overline{F_\omega }(\stackrel{}{r},\stackrel{}{r}^{})\left\{i\omega \gamma ^0+i\stackrel{}{\gamma }\left(\stackrel{}{}_\stackrel{}{r}+ie\stackrel{}{A}(\stackrel{}{r})\right)m+\gamma ^0\mu \right\}\mathrm{\Delta }^{}(\stackrel{}{r})\gamma ^0G_\omega (\stackrel{}{r},\stackrel{}{r}^{})=0,\hfill \end{array}$$
(7.4.3)
where the gap function is defined by $`\mathrm{\Delta }^{}(\stackrel{}{r})=\lambda T{\displaystyle \underset{n}{}}F_\omega ^+(\stackrel{}{r},\stackrel{}{r}^{}).`$ The Bloch individual particle Green’s function $`\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})`$ for the fermion in normal mode is written
$$\left\{i\omega \gamma ^0i\stackrel{}{\gamma }\left(\stackrel{}{}_\stackrel{}{r}ie\stackrel{}{A}(\stackrel{}{r})\right)m+\gamma ^0\mu \right\}\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})=\delta (\stackrel{}{r}\stackrel{}{r}^{}),$$
(7.4.4)
or the adjoint equation
$$\left\{i\omega \gamma ^0+i\stackrel{}{\gamma }\left(\stackrel{}{}_\stackrel{}{r}^{}+ie\stackrel{}{A}(\stackrel{}{r}^{})\right)m+\gamma ^0\mu \right\}\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})=\delta (\stackrel{}{r}\stackrel{}{r}^{}).$$
(7.4.5)
By means of eq.(7.4.5) the eq.(7.4.3) gives rise to
$$G_\omega (\stackrel{}{r},\stackrel{}{r}^{})=\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{s})\gamma ^0\mathrm{\Delta }(\stackrel{}{s})\overline{F_\omega }(\stackrel{}{s},\stackrel{}{r}^{})d^3s,$$
(7.4.6)
and
$$\overline{F_\omega }(\stackrel{}{r},\stackrel{}{r}^{})=\stackrel{~}{G}_\omega (\stackrel{}{s},\stackrel{}{r}^{})\mathrm{\Delta }^{}(\stackrel{}{s})\gamma ^0\stackrel{~}{G}_\omega (\stackrel{}{s},\stackrel{}{r})d^3s.$$
(7.4.7)
The gap function $`\mathrm{\Delta }(\stackrel{}{r})`$ as well as $`\overline{F_\omega }(\stackrel{}{r},\stackrel{}{r}^{})`$ are small ones at close neighbourhood of transition temperature $`T_c`$ and varied slowly over a coherence distance. This approximation, which went into the derivation of equations, meets our interest in eq.(7.4.6), eq.(7.4.7). Using standard procedure one may readily express them in power series of $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ by keeping only the terms in $`\overline{F_\omega }(\stackrel{}{r},\stackrel{}{r}^{})`$ up to the cubic and in $`G_\omega (\stackrel{}{r},\stackrel{}{r}^{})`$ \- quadratic order in $`\mathrm{\Delta }`$. After averaging over the polarization of particles the following equation coupling $`\mathrm{\Delta }(\stackrel{}{r})`$ and $`\stackrel{}{A}(\stackrel{}{r})`$ ensued:
$$\begin{array}{c}\overline{\mathrm{\Delta }^{}(\stackrel{}{r})}=\lambda T\underset{n}{}\overline{\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})}\mathrm{\Delta }^{}(\stackrel{}{r}^{})d^3r^{}\hfill \\ \\ \lambda T\underset{n}{}\overline{\stackrel{~}{G}_\omega (\stackrel{}{s},\stackrel{}{r})\stackrel{~}{G}_\omega (\stackrel{}{s},\stackrel{}{l})\stackrel{~}{G}_\omega (\stackrel{}{m},\stackrel{}{l})\stackrel{~}{G}_\omega (\stackrel{}{m},\stackrel{}{r})}\mathrm{\Delta }(\stackrel{}{s})\mathrm{\Delta }^{}(\stackrel{}{l})\mathrm{\Delta }^{}(\stackrel{}{m})d^3sd^3ld^3m.\hfill \end{array}$$
(7.4.8)
It is worthwhile to determine the function $`\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})`$. In the absence of applied magnetic field, i.e. $`\stackrel{}{A}(\stackrel{}{r})=0`$ the eq.(7.4.4) reduces to
$$\left\{(i\omega +\mu )\gamma ^0\stackrel{}{\gamma }\stackrel{}{p}m\right\}\stackrel{~}{G}_\omega ^0(\stackrel{}{r},\stackrel{}{r}^{})=\delta (\stackrel{}{r}\stackrel{}{r}^{}).$$
(7.4.9)
It is well to study at this point certain properties of the solution which we shall continually encounter
$$\stackrel{~}{G}_\omega ^0(\stackrel{}{r},\stackrel{}{r}^{})=\frac{1}{2m}\left\{(i\omega +\mu )\gamma ^0\stackrel{}{\gamma }\stackrel{}{p}+m\right\}\stackrel{~}{G}_{G\omega }^0(\stackrel{}{r}\stackrel{}{r}^{}),$$
(7.4.10)
where the function $`\stackrel{~}{G}_{G\omega }^0(\stackrel{}{r}\stackrel{}{r}^{})`$ satisfies the equation
$$\frac{1}{2m}\left(q^2+\mathrm{\Delta }\right)\stackrel{~}{G}_{G\omega }^0(\stackrel{}{r}\stackrel{}{r}^{})=\delta (\stackrel{}{r}\stackrel{}{r}^{}),$$
(7.4.11)
provided by $`q^2=(i\omega +\mu )^2m^2=2im\mathrm{\Omega }+p_0^2`$ and $`\mathrm{\Omega }=\omega {\displaystyle \frac{\mu }{m}},p_0^2=\mu ^2m^2.`$ At $`\mu \omega `$ one has $`{\displaystyle \frac{p_0^2}{2m}}\mathrm{\Omega }`$, and
$$\frac{1}{2m}\left\{2im\mathrm{\Omega }+p_0^2+\mathrm{\Delta }\right\}\stackrel{~}{G}_{G\omega }^0(\stackrel{}{r}\stackrel{}{r}^{})=\delta (\stackrel{}{r}\stackrel{}{r}^{}),$$
(7.4.12)
the solution of which reads
$$\stackrel{~}{G}_{G\omega }^0(\stackrel{}{r}\stackrel{}{r}^{})=\frac{m}{2\pi R}\mathrm{exp}(iqR),$$
where
$$q=sgn\mathrm{\Omega }p_0+i\frac{\mathrm{\Omega }}{v},R=\stackrel{}{r}\stackrel{}{r}^{},sgn\mathrm{\Omega }=\frac{\mathrm{\Omega }}{\mathrm{\Omega }}.$$
In approximation to non-relativistic limit $`\stackrel{}{p}0`$ this Green’s function reduces to
$`\stackrel{~}{G}_{G\omega }^0(\stackrel{}{r}\stackrel{}{r}^{})`$ used in . Making use of Fourier integrals we readily get
$$\stackrel{~}{G}_\omega ^0(\stackrel{}{p})=\frac{\mu \gamma ^0+\stackrel{}{\gamma }\stackrel{}{p}+m}{q^2\stackrel{}{p}^2+i0}=\widehat{\gamma }\stackrel{~}{G}_{G\omega }^0(\stackrel{}{p}),$$
where $`\widehat{\gamma }={\displaystyle \frac{1}{2m}}(\mu \gamma ^0+\stackrel{}{\gamma }\stackrel{}{p}+m).`$ One has
$$\stackrel{~}{G}_{G\omega }^0(\stackrel{}{p})=\stackrel{~}{G}_{G\mathrm{\Omega }}^0(\stackrel{}{p})=\frac{1}{i\mathrm{\Omega }\xi }=\stackrel{~}{G}_{G\mathrm{\Omega }}^0(\stackrel{}{p}),$$
where as usual $`\xi ={\displaystyle \frac{\stackrel{}{p}}{2m}}{\displaystyle \frac{\stackrel{}{p}_0}{2m}}`$. The Green’s function $`\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})`$ in presence of magnetic field differs from $`\stackrel{~}{G}_\omega ^0(\stackrel{}{r}\stackrel{}{r}^{})`$ only by phase multiplier
$$\stackrel{~}{G}_\omega (\stackrel{}{r},\stackrel{}{r}^{})=\mathrm{exp}\left\{\frac{ie}{c}(\stackrel{}{A}(\stackrel{}{r}),\stackrel{}{r}\stackrel{}{r}^{})\right\}\stackrel{~}{G}_\omega ^0(\stackrel{}{r}\stackrel{}{r}^{}).$$
The technique now is to expand a second term in right-hand side of the eq.(7.4.8) up to the terms quadratic in $`(\stackrel{}{r}\stackrel{}{r}^{})`$. After calculations it transforms
$$\begin{array}{c}\{(i\mathrm{}\stackrel{}{}+\frac{e^{}}{c}\stackrel{}{A})^2+\frac{2m}{\nu }[\frac{2\pi ^2}{\lambda mp_0}\left(\frac{\mu }{m}\right)^2(\frac{\mu }{m}1)+\left(\frac{\mu }{m}\right)^2(\frac{T}{T_{c\mu }}1)+\hfill \\ \\ \frac{2}{N}\mathrm{\Psi }(\stackrel{}{r})^2]\}\mathrm{\Psi }(\stackrel{}{r})=0,\hfill \end{array}$$
(7.4.13)
where $`\nu ={\displaystyle \frac{7\zeta (3)mv_F^2}{24(\pi k_BT_c)^2}}`$ and $`T_{c\mu }={\displaystyle \frac{m}{\mu }}T_c`$. Succinctly
$$\left\{\stackrel{}{p}_A^2\frac{1}{2}m_\mathrm{\Psi }^2+\frac{1}{4}\lambda _\mathrm{\Psi }^2\mathrm{\Psi }(\stackrel{}{r})^2\right\}\mathrm{\Psi }(\stackrel{}{r})=0,$$
(7.4.14)
provided by
$$\begin{array}{c}m_\mathrm{\Psi }^2(\lambda ,T,T_{c\mu })=\frac{24}{7\zeta (3)}\left(\frac{\mathrm{}}{\xi _0}\right)^2\left(\frac{\mu }{m}\right)^2\left[1\frac{T}{T_{c\mu }}\left(\frac{\mu }{m}1\right)\mathrm{ln}\frac{2\stackrel{~}{\omega }}{\mathrm{\Delta }_0}\right],\hfill \\ \\ \lambda _\mathrm{\Psi }^2(\lambda ,T_c)=\frac{96}{7\zeta (3)}\left(\frac{\mathrm{}}{\xi _0}\right)^2\frac{1}{N},\mathrm{\Psi }(\stackrel{}{r})=\mathrm{\Delta }(\stackrel{}{r})\frac{\left(7\zeta (3)N\right)^{1/2}}{4\pi k_BT_c}.\hfill \end{array}$$
(7.4.15)
According to the eq.(7.4.15), the magnitude of the relativistic effects, however, is found to be greater to account for the large contribution to the values $`m_\mathrm{\Psi }^2`$ and $`\lambda _\mathrm{\Psi }^2`$. While the transition temperature decreases inversely by the relativistic factor $`{\displaystyle \frac{\mu }{m}}`$. A spontaneous breakdown of symmetry of ground state occurs at $`\eta _\mathrm{\Psi }^2(\lambda ,T<T_{c\mu })>0,`$ where $`\eta _\mathrm{\Psi }^2(\lambda ,T,T_{c\mu })={\displaystyle \frac{m_\mathrm{\Psi }^2}{\lambda _\mathrm{\Psi }^2}}.`$ The eq.(7.4.14) splits into the couple of equations for $`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_R`$. Subsequently, a Lagrangian of the $`\phi `$ will be arisen with the corresponding values of mass $`m_\mathrm{\Psi }^2m_\phi ^2`$ and coupling constant $`\lambda _\mathrm{\Psi }^2\lambda _\phi ^2`$.
### 7.5 The Four-Component Bose-Condensate in Magnetic Field
Now we are going to derive the equation of four-component bispinor field of Bose-condensate, which may sound strange, but due to self-interaction the spin part of it is vanished. We start with the nonsymmetric state $`\mathrm{\Delta }_L\mathrm{\Delta }_R`$, where $`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_R`$ are two eigenstates of chirality operator $`\gamma _5`$. In standard representation
$$\begin{array}{c}\mathrm{\Psi }=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)\left(\begin{array}{c}\mathrm{\Psi }_L\\ \mathrm{\Psi }_R\end{array}\right)=\frac{1}{\sqrt{2}}\left(\begin{array}{c}\mathrm{\Psi }_L+\mathrm{\Psi }_R\\ \mathrm{\Psi }_L\mathrm{\Psi }_R\end{array}\right)\left(\begin{array}{c}\mathrm{\Psi }_1\\ \mathrm{\Psi }_2\end{array}\right),\hfill \\ \\ \mathrm{\Delta }=\left(\begin{array}{c}\mathrm{\Delta }_1\\ \mathrm{\Delta }_2\end{array}\right)=\frac{1}{\sqrt{2}}\left(\begin{array}{c}\mathrm{\Delta }_L+\mathrm{\Delta }_R\\ \mathrm{\Delta }_L\mathrm{\Delta }_R\end{array}\right),\mathrm{\Delta }_L\mathrm{\Delta }_R.\hfill \end{array}$$
(7.5.1)
The eq.(7.4.14) enables to postulate the equation of four-component Bose-condensate in magnetic field and equilibrium state
$$i\mathrm{}\frac{\mathrm{\Psi }}{t}=\left\{c\stackrel{}{\alpha }\left(\stackrel{}{p}+\frac{e^{}}{c}\stackrel{}{A}\right)+\beta mc^2+M(F)+L(F)\mathrm{\Psi }^2\right\}\mathrm{\Psi }=0,$$
(7.5.2)
or succinctly
$$\left(\gamma p_Am\right)\mathrm{\Psi }=0.$$
(7.5.3)
This is in standard notations
$$\begin{array}{c}\gamma ^0=\beta =\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\stackrel{}{\alpha }=\gamma ^0\stackrel{}{\gamma }=\left(\begin{array}{cc}0& \stackrel{}{\sigma }\\ \stackrel{}{\sigma }& 0\end{array}\right),F=F_{\mu \nu }\sigma ^{\mu \nu }=inv,\hfill \\ \sigma ^{\mu \nu }=\frac{1}{2}[\gamma ^\mu ,\gamma ^\nu ],\frac{i}{2}e^{}F=e^{}\stackrel{}{\mathrm{\Sigma }}\stackrel{}{H},\stackrel{}{H}=rot\stackrel{}{A},F_{\mu \nu }=(0,\stackrel{}{H}),\stackrel{}{\mathrm{\Sigma }}=\left(\begin{array}{cc}\stackrel{}{\sigma }& 0\\ 0& \stackrel{}{\sigma }\end{array}\right),\hfill \\ p_A=(p_{A0},\stackrel{}{p}_A),\stackrel{}{p}_A=i\mathrm{}\stackrel{}{}+\frac{e^{}}{c}\stackrel{}{A},p_{A0}=\left(M(F)+L(F)\mathrm{\Psi }^2\right).\hfill \end{array}$$
The $`M(F)`$ and $`L(F)`$ are some functions depending upon the invariant F, which will be determined under the requirement that the second-order equations ensued from the eq.(7.5.3) must match onto eq.(7.4.14). Defining the functions $`M(F)`$ and $`L(F)`$:
$$\begin{array}{c}M(F)=\left(M_0^2+\frac{i}{2}e^{}F\right)^{1/2},M_0=\left(m^2+\frac{1}{2}m_\phi ^2\right)^{1/2},L(F)=\frac{\lambda _\phi ^2}{8M(F)},\hfill \\ \\ L_0=\frac{\lambda _\phi ^2}{8M_0},M(F)L(F)=M_0L_0=\frac{1}{8}\lambda _\phi ^2,\hfill \end{array}$$
and taking into account an approximation fitting our interest that the gap function is small at close neighbourhood of transition temperature, one gets
$$\left\{\stackrel{}{p}_A^2+m^2\left(M_0+L_0\mathrm{\Psi }^2\right)^2\right\}\mathrm{\Psi }(\stackrel{}{r})\left\{\stackrel{}{p}_A^2\frac{1}{2}m_\phi ^2+\frac{1}{4}\lambda _\phi ^2\mathrm{\Psi }^2\right\}\mathrm{\Psi }(\stackrel{}{r})=0.$$
(7.5.4)
This has yet another important consequence that at $`\mathrm{\Delta }_L0`$ and imposed constraint
$`\left(m+M(F)+L(F)\mathrm{\Psi }^2\right)_{F0}0`$ we have
$$\mathrm{\Delta }_2=\frac{1}{\sqrt{2}}(\mathrm{\Delta }_L\mathrm{\Delta }_R)=0,\mathrm{\Psi }_2=0,$$
(7.5.5)
where $`\mathrm{\Psi }_0`$ is the gap function symmetry-restoring value
$$\mathrm{\Delta }_2\left(\mathrm{\Psi }_0^2=\frac{m+M_0}{L_0}\right)=0,\mathrm{\Delta }_L\left(\mathrm{\Psi }_0^2\right)=\mathrm{\Delta }_R\left(\mathrm{\Psi }_0^2\right),$$
where, according to eq.(7.5.4), one has
$$V\left[m^2\left(M_0+L_0\mathrm{\Psi }^2\right)^2\right]\mathrm{\Psi }^2=\left[\frac{1}{2}m_\phi ^2+\frac{1}{4}\lambda _\phi ^2\mathrm{\Psi }^2\right]\mathrm{\Psi }^2,$$
(7.5.6)
and
$$V\left(\mathrm{\Psi }_0^2=\frac{m+M_0}{L_0}=\frac{1}{2}\eta _\phi ^2(\lambda ,T,T_{c\mu })\right)=0.$$
(7.5.7)
It leads us to the conclusion that the field of symmetry-breaking Higgs boson must be counted off from the $`\mathrm{\Delta }_L=\mathrm{\Delta }_R`$ symmetry-restoring value of Bose-condensate $`\mathrm{\Psi }_0={\displaystyle \frac{1}{\sqrt{2}}}\eta _\phi (\lambda ,T,T_{c\mu })`$ as the point of origin describing the excitation in the neighbourhood of stable vacuum eq.(7.5.7). We may write down the Lagrangian corresponding to the eq.(7.5.3)
$$L_\mathrm{\Psi }=\frac{1}{2}\left\{\overline{\mathrm{\Psi }}\gamma p_A\mathrm{\Psi }\overline{\mathrm{\Psi }}\gamma \stackrel{}{p}_A\mathrm{\Psi }\right\}m\overline{\mathrm{\Psi }}\mathrm{\Psi }.$$
The gauge invariant Lagrangian eq.(6.7.7) takes the form
$$\underset{W}{𝐿}(\underset{W}{𝐷})=\frac{i}{2}\left\{\underset{W}{\overline{\mathrm{\Psi }}}\gamma \underset{W}{𝐷}\underset{W}{\Psi }\underset{W}{\overline{\mathrm{\Psi }}}\gamma \stackrel{}{\underset{W}{𝐷}}\underset{W}{\Psi }\right\}\underset{W}{\overline{\mathrm{\Psi }}}\left\{m+\gamma ^0\left[M(F)+L(F)\mathrm{\Psi }^2\right]\right\}\underset{W}{\Psi }.$$
(7.5.8)
At the symmetry-restoring point, this Lagrangian can be replaced
$$\underset{W}{𝐿}(\underset{W}{𝐷})\underset{W}{𝐿}{}_{1}{}^{}(\underset{W}{𝐷})=\frac{1}{2}(\underset{W}{𝐷}\underset{W}{\Psi }{}_{1}{}^{})^2\underset{W}{𝑉}(\underset{W}{\Psi }{}_{1}{}^{}_{}^{2}),$$
provided by
$$\underset{W}{𝑉}(\underset{W}{\Psi }{}_{1}{}^{}_{}^{2})=\frac{1}{2}m_\phi ^2\underset{W}{\Psi }{}_{1}{}^{2}+\frac{1}{4}\lambda _\phi ^2\underset{W}{\Psi }{}_{1}{}^{4}.$$
Taking into account the eq.(7.1.8), in which $`\underset{W}{\Psi }{}_{1}{}^{}_{}^{2}=\phi ^2={\displaystyle \frac{1}{2}}\eta _\phi +\chi ^2,`$ one gets
$$\underset{W}{𝐿}{}_{\phi }{}^{}(\underset{W}{𝐷})=\frac{1}{2}\left(\underset{W}{𝐷}\phi \right)^2\underset{W}{𝑉}(\phi ^2),\underset{W}{𝑉}(\phi ^2)=\frac{1}{2}m_\phi ^2\phi ^2+\frac{1}{4}\lambda _\phi ^2\phi ^4.$$
(7.5.9)
The average value of well-defined current source expressed in terms of spinless field $`\mathrm{\Psi }`$ eq.(7.5.3) is given by $`\stackrel{}{j}(\stackrel{}{r})|_{\stackrel{}{\mathrm{\Sigma }}=0}.`$ According to eq.(7.4.6) and eq.(7.4.7), one has
$$\begin{array}{c}\delta G(x,x^{})|_{x=x^{}}=T\underset{\omega }{}\stackrel{~}{G}_{G\mathrm{\Omega }}(\stackrel{}{r},\stackrel{}{s})\stackrel{~}{G}_{G\mathrm{\Omega }}(\stackrel{}{l},\stackrel{}{r})\stackrel{~}{G}_{G\mathrm{\Omega }}(\stackrel{}{l},\stackrel{}{s})\mathrm{\Delta }^{}(\stackrel{}{l})\mathrm{\Delta }(\stackrel{}{s})d^3sd^3l.\hfill \end{array}$$
(7.5.10)
Hence
$$\stackrel{}{j}(\stackrel{}{r})|_{\stackrel{}{\mathrm{\Sigma }}=0}=\frac{m}{4(mM_0L_0\mathrm{\Psi }^2)}(\gamma ^0+\frac{\mu }{m})\{(\gamma ^0+\frac{\mu }{m})^2+3\beta _F^2\}\stackrel{}{j}_G(\stackrel{}{r}),$$
(7.5.11)
provided by
$$\stackrel{}{j}_G(\stackrel{}{r})=\left(\frac{m}{\mu }\right)^3\left\{\frac{ie^{}\mathrm{}}{2m}\left(\mathrm{\Psi }\frac{\mathrm{\Psi }^{}}{\stackrel{}{r}}\mathrm{\Psi }^{}\frac{\mathrm{\Psi }}{\stackrel{}{r}}\right)\frac{e_{}^{}{}_{}{}^{2}}{mc}\stackrel{}{A}|\mathrm{\Psi }|^2\right\}.$$
(7.5.12)
Below we write the eq.(7.5.3) in the form on close analogy of the elementary excitations in superconductivity model described by coherent mixture of electrons and holes near the Fermi surface . In chirial representation $`\mathrm{\Psi }=\left(\begin{array}{c}\mathrm{\Psi }_L\\ \mathrm{\Psi }_R\end{array}\right)`$ one has
$$p_{A0}\mathrm{\Psi }_L=\stackrel{}{\sigma }\stackrel{}{p}_A\mathrm{\Psi }_L+m\mathrm{\Psi }_R,p_{A0}\mathrm{\Psi }_R=\stackrel{}{\sigma }\stackrel{}{p}_A\mathrm{\Psi }_R+m\mathrm{\Psi }_L,p_{A0}=\pm \left(\stackrel{}{p}_A^2+m^2\right)^{1/2}.$$
The two states of quasi-particle are separated in energy by $`2p_{A0}`$. In the ground state all quasi-particles should be in lower (negative) energy states. It would take a finite energy $`2p_{A0}2m`$ to excite a particle to the upper state (the case of Dirac particle). Thus, one may assume that the energy gap parameter $`m`$ is also due to some interaction between massless bare fermions .
### 7.6 Extension to Lower Temperatures
It is worth briefly recording the question of whether or not it is possible to extend the ideas of former approach to lower temperatures as it was investigated in the case of Gor’kov’s theory by others \[85-87\]. Here, as usual we admit that the order parameter and vector potential vary slowly over distances of the order of the coherence length. We restrict ourselves to the London limit and the derivation of equations will be proceeded by iterating to a low order giving only the leading terms. Taking into account the eq.(7.4.1), eq.(7.4.3), eq.(7.4.6) and eq.(7.4.7), it is straightforward to derive the separate integral equations for $`G`$ and $`F^+`$ in terms of $`\mathrm{\Delta }`$, $`\mathrm{\Delta }^{}`$ and $`\stackrel{~}{G}`$. We introduce
$$K_\mathrm{\Omega }(\stackrel{}{r},\stackrel{}{s})=\delta (\stackrel{}{r}\stackrel{}{s})\left\{i\mathrm{\Omega }+\frac{1}{2m}\left(\frac{}{\stackrel{}{s}}i\frac{e}{c}\stackrel{}{A}(\stackrel{}{s})\right)^2+\mu _0\right\}+\mathrm{\Delta }(\stackrel{}{r})\widehat{\gamma }_A^2(\stackrel{}{s})\stackrel{~}{G}_\mathrm{\Omega }(\stackrel{}{s},\stackrel{}{r})\mathrm{\Delta }^{}(\stackrel{}{s}),$$
(7.6.1)
provided by $`\mu _0={\displaystyle \frac{p_0^2}{2m}}`$, and
$$F_\omega ^+(\stackrel{}{s},\stackrel{}{r}^{})=\widehat{\gamma }_A(\stackrel{}{s})F_\mathrm{\Omega }^+(\stackrel{}{s},\stackrel{}{r}^{}),\widehat{\gamma }_A(\stackrel{}{s})=\frac{1}{2m}\{\gamma ^0(i\omega +\mu )\stackrel{}{\gamma }\stackrel{}{p}_A(\stackrel{}{s})+m\}.$$
We write down the coupled equations in the form
$$d^3sK_\mathrm{\Omega }(\stackrel{}{r},\stackrel{}{s})G_\mathrm{\Omega }(\stackrel{}{s},\stackrel{}{r}^{})=\delta (\stackrel{}{r}\stackrel{}{r}^{}),$$
(7.6.2)
and
$$d^3sF_\mathrm{\Omega }^+(\stackrel{}{s},\stackrel{}{r}^{})K_\mathrm{\Omega }(\stackrel{}{s},\stackrel{}{r})=\mathrm{\Delta }^{}(\stackrel{}{r})\widehat{\gamma }_A(\stackrel{}{r}^{})\stackrel{~}{G}_\mathrm{\Omega }(\stackrel{}{r},\stackrel{}{r}^{}),$$
(7.6.3)
The mathematical structure of obtained equations is closely similar to that studied by in somewhat different context. So, adopting their technique we introduce sum and difference coordinates, and Fourier transform with respect to the difference coordinates as follows:
$$K_\mathrm{\Omega }(\stackrel{}{p},\stackrel{}{R})d^3(rs)e^{i\stackrel{}{p}(\stackrel{}{r}\stackrel{}{s})}K_\mathrm{\Omega }(\stackrel{}{r},\stackrel{}{s})$$
(7.6.4)
with $`\stackrel{}{R}{\displaystyle \frac{1}{2}}(\stackrel{}{r}+\stackrel{}{s})`$. We also involve similar expansions for all other functions. Then, the eq.(7.6.2) and eq.(7.6.3) reduce to following:
$$\begin{array}{c}\mathrm{\Theta }\left[K_\mathrm{\Omega }(\stackrel{}{p},\stackrel{}{R})G_\mathrm{\Omega }(\stackrel{}{p}^{},\stackrel{}{R}^{})\right]=1,\hfill \\ \\ \mathrm{\Theta }\left[F_\mathrm{\Omega }^+(\stackrel{}{p}^{},\stackrel{}{R}^{})K_\mathrm{\Omega }(\stackrel{}{p},\stackrel{}{R})\right]=\mathrm{\Theta }\left[\mathrm{\Delta }^{}(\stackrel{}{R})\widehat{\gamma }_A(\stackrel{}{p}^{},\stackrel{}{R}^{})\stackrel{~}{G}_\mathrm{\Omega }(\stackrel{}{p}^{},\stackrel{}{R}^{})\right]\hfill \end{array}$$
(7.6.5)
provided by the standard differential operator of finite order defined under the requirement that it produces the Fourier transform of the matrix product of two functions when it operates on the transforms of the individual functions
$$\mathrm{\Theta }\underset{\begin{array}{c}\stackrel{}{R}^{}\stackrel{}{R}\\ \stackrel{}{p}^{}\stackrel{}{p}\end{array}}{lim}\mathrm{exp}\left[\frac{i}{2}\left(\frac{}{\stackrel{}{R}}\frac{}{\stackrel{}{p}^{}}\frac{}{\stackrel{}{p}}\frac{}{\stackrel{}{R}^{}}\right)\right].$$
(7.6.6)
One gets
$$\begin{array}{c}K_\mathrm{\Omega }(\stackrel{}{p},\stackrel{}{R})=i\mathrm{\Omega }ϵ(\stackrel{}{p},\stackrel{}{R})+\underset{\stackrel{}{R}^{},\stackrel{}{R}^{\prime \prime }\stackrel{}{R}}{lim}\mathrm{exp}\left[\frac{i}{2}\frac{}{\stackrel{}{p}}(\frac{}{\stackrel{}{R}^{}}\frac{}{\stackrel{}{R}^{\prime \prime }})\right]\times \hfill \\ \\ \times \mathrm{\Delta }^{}(\stackrel{}{R}^{})\widehat{\gamma }_A(\stackrel{}{p},\stackrel{}{R}^{})\widehat{\gamma }_A(\stackrel{}{p},\stackrel{}{R})\mathrm{\Delta }^{}(\stackrel{}{R}^{\prime \prime }),\hfill \end{array}$$
(7.6.7)
where it is denoted $`ϵ(\stackrel{}{p},\stackrel{}{R}){\displaystyle \frac{1}{2m}}\left(\stackrel{}{p}{\displaystyle \frac{e}{c}}\stackrel{}{A}(\stackrel{}{R})\right)^2\mu _0.`$ To obtain resulting expressions we shall proceed with further calculations, but shall forbear to write them out as they are so standard. There is only one thing to be noticed about the integration. That is, due to the angular integration in momentum space, as mentioned above, the terms linear in the vector $`\stackrel{}{p}`$ will be vanished , as well as the integration over the energies removes the linear terms in $`ϵ(\stackrel{}{p})`$. So, we may expand the quantities in eq.(7.6.5) according to the degree of inhomogenity somewhat like it we have done in equation (7.4.8) of gap function $`\mathrm{\Delta }^{}(\stackrel{}{r})`$, which in mixed representation transforms to the following:
$$\mathrm{\Delta }^{}(\stackrel{}{R})=T\underset{\omega }{}\frac{d^3p}{(2\pi )^3}F_\omega ^+(\stackrel{}{p},\stackrel{}{R})=T\underset{\omega }{}\frac{d^3p}{(2\pi )^3}\widehat{\gamma }_A(\stackrel{}{p},\stackrel{}{R})F_\mathrm{\Omega }^+(\stackrel{}{p},\stackrel{}{R}).$$
(7.6.8)
The approximation was used to obtain the function $`F_\mathrm{\Omega }^+`$ must be of one order higher $`F_\mathrm{\Omega }^+F_\mathrm{\Omega }^{(0)+}+F_\mathrm{\Omega }^{(1)+}+F_\mathrm{\Omega }^{(2)+}`$ than that for function $`\stackrel{~}{G}_\mathrm{\Omega }\stackrel{~}{G}_\mathrm{\Omega }^{(0)}+\stackrel{~}{G}_\mathrm{\Omega }^{(1)}`$. Employing an iteration method of solution one replaces $`K\stackrel{~}{K},G\stackrel{~}{G}`$ in eq.(7.6.5) and puts $`\mathrm{\Theta }^{(0)}=1,\stackrel{~}{K}^{(1)}=0,\stackrel{~}{G}^{(1)}=0`$. Hence $`\stackrel{~}{G}=\stackrel{~}{G}^{(0)}`$. The resulting equation for gap function is similar to those occurring in , although not identical. The sole difference is that in the resulting equation we use the expressions of $`\mathrm{\Omega }`$ and $`\xi `$. With this replacement the equation reads
$$\begin{array}{c}\mathrm{\Delta }^{}=T\underset{\omega }{}\frac{d^3p}{(2\pi )^3}\{\frac{\mathrm{\Delta }^{}}{\mathrm{\Omega }^2+\xi ^2}+[(\frac{}{\stackrel{}{R}}+\frac{2ie}{c}\stackrel{}{A})^2\mathrm{\Delta }^{}+\mathrm{\Delta }\left((\frac{}{\stackrel{}{R}}+\frac{2ie}{c}\stackrel{}{A})\mathrm{\Delta }^{}\right)^2\frac{}{|\mathrm{\Delta }|^2}+\hfill \\ \\ +\frac{\mathrm{\Delta }^{}}{3}\frac{^2|\mathrm{\Delta }|^2}{\stackrel{}{R}^2}\frac{}{|\mathrm{\Delta }|^2}+\frac{\mathrm{\Delta }^{}}{6}\left(\frac{|\mathrm{\Delta }|^2}{\stackrel{}{R}}\right)^2\frac{^2}{\left(|\mathrm{\Delta }|^2\right)^2}]\frac{p^2/6m^2}{\left(\mathrm{\Omega }^2+\xi ^2\right)^2}\},\hfill \end{array}$$
(7.6.9)
where $`\xi (\stackrel{}{p},\stackrel{}{R})\left[ϵ^2(\stackrel{}{p})+|\mathrm{\Delta }(\stackrel{}{R})|^2\right]^{1/2}.`$ The average value of the operator of current density at $`TT_c`$ follows at once
$$\begin{array}{c}\stackrel{}{j}_G(\stackrel{}{R})=\left(\overline{\gamma ^0\widehat{\gamma }^3}\right)_{\mathrm{\Sigma }=0}\frac{2e}{m}\{\frac{i}{2}(\mathrm{\Delta }^{}(\stackrel{}{R})\frac{\mathrm{\Delta }(\stackrel{}{R})}{\stackrel{}{R}}\frac{\mathrm{\Delta }^{}(\stackrel{}{R})}{\stackrel{}{R}}\mathrm{\Delta }(\stackrel{}{R}))\hfill \\ \\ \frac{2e}{c}|\mathrm{\Delta }(\stackrel{}{R})|^2\stackrel{}{A}(\stackrel{}{R})\}T_\omega \frac{d^3p}{(2\pi )^3}\frac{p^2/3m}{\left(\mathrm{\Omega }^2+\xi ^2(\stackrel{}{p},\stackrel{}{R})\right)^2},\hfill \end{array}$$
(7.6.10)
where $`\left(\overline{\gamma ^0\widehat{\gamma }^3}\right)_{\mathrm{\Sigma }=0}={\displaystyle \frac{1}{8}}\left(\gamma ^0+{\displaystyle \frac{\mu }{m}}\right)\left\{\left(\gamma ^0+{\displaystyle \frac{\mu }{m}}\right)^2+3\beta _F^2\right\}`$. At $`\mathrm{\Delta }\pi k_BT`$ and $`\stackrel{}{A}`$ is independent of position the eq.(7.6.9) and eq.(7.6.10) lead back to the equations (7.4.13) and (7.5.11). Actually, from such results it is then easy by ordinary manipulations to investigate the pertinent physical problem in several particular cases, but a separate calculation for each case would be needed.
## 8 Lagrangian of Electroweak Interactions; The Transmission of the Electroweak Symmetry Breaking From the $`W`$-World to Spacetime Continuum
The results obtained within the previous subsections enable us to trace unambiguously rather general scheme of unified electroweak interactions, where the self-interacting isospinor scalar Higgs bosons have arisen as the collective modes of excitations of bound quasi-particle iso-pairs on the internal $`W`$-world. But, at the very first we remind some features allowing us to write down the final Lagrangian of electroweak interactions.
1. During the realization of the MW connections of weak interacting fermions under the action of the $`Q`$-world the P-violation compulsory occurred in the W-world incorporated with the symmetry reduction eq.(6.7.1) characterized by the Weinberg mixing angle with the fixed value at $`30^0`$. This gives rise to the local symmetry $`SU(2)U(1)`$, under which the left-handed fermions transformed as six independent doublets, while the right-handed fermions transformed as twelve independent singlets.
2. Due to vacuum rearrangement in Q-world the Yukawa couplings arise between the fermion fields and corresponding isospinor-scalar $`\phi `$\- meson in conventional form.
3. In the framework of suggested mechanism providing the effective attraction between the relativistic fermions caused by the exchange of the mediating induced gauge quanta in W-world, the self-interacting isospinor-scalar Higgs bosons arise as Bose-condensate, namely the $`SU(2)`$ multiplets of spinless $`\phi `$-meson fields coupled to the gauge fields in a gauge invariant way. Thus, in the Lagrangian of $`\phi `$-meson with the degenerate vacuum of the W-world the symmetry-breaking Higgs boson is counted off from the gap symmetry restoring value as the point of origin. In view of this the total Lagrangian ensues from the eq.(6.7.5)- eq.(6.7.7), which is now invariant under the local symmetry $`SU(2)U(1)`$, where a set of gauge fields are coupled to various multiplets of fields among which is also a multiplet of Higgs boson. Subsequently, we separate a piece of Lagrangian containing only the fields defined on four dimensional Minkowski flat spacetime continuum $`M_4`$. To facilitate writing we shall forbear here to write out the piece of Lagrangian containing the terms of other fermion generations than one, as it is a somewhat lengthy and so standard. But, in the mean time, we shall retain the explicit terms of Higgs bosons arisen on the internal $`W`$world to emphasize the specific mechanism of the electroweak symmetry breakdown discussed below. The resulting Lagrangian reads
$$\begin{array}{c}L=\frac{1}{2}TrG_{\mu \nu }G^{\mu \nu }\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+i\overline{L}\widehat{D}L+i\overline{e}_R\widehat{D}e_R+i\overline{\nu }_R\widehat{D}\nu _R+|\underset{W}{𝐷}{}_{\mu }{}^{}\phi |^2\hfill \\ \\ \frac{1}{2}\lambda _\phi ^2\left(|\phi |^2\frac{1}{2}\eta _\phi ^2\right)^2f_e\left(\overline{L}\phi e_R+\overline{e}_R\phi ^+L\right)f_\nu \left(\overline{L}\phi _c\nu _R+\overline{\nu }_R\phi _c^+L\right)+\hfill \\ \\ \text{similar terms for other fermion generations}.\hfill \end{array}$$
(8.1)
The gauge fields $`𝐀_\mu (x)`$ and $`B_\mu (x)`$ associate respectively with the groups $`SU(2)`$ and $`U(1)`$, where the gauge covariant curls are $`F_{\mu \nu },G_{\mu \nu }.`$ The corresponding gauge covariant derivatives are in standard form. One took into account corresponding values of the operators $`𝐓`$ and $`Y`$ for left- and right-handed fields, and for isospinor $`\phi `$-meson. The Yukawa coupling constants $`f_e`$ and $`f_\nu `$ are inserted in subsec.6.10. Since the electroweak symmetry is at any rate only approximate, the test of the theory will depend on its ability to account for its breaking as well. here the MSM creates a particular incentive for the study of such a breaking. Then, just it remains to see how can such Higgs bosons arisen on the internal $`W`$-world break the gauge symmetry down in $`M_4`$ and lead to masses of the spacetime-components of the MW-fields? It is remarkable to see that the suggested MSM, in contrast to the SM, predicts the transmission of electroweak symmetry breaking from the $`W`$world to the $`M_4`$ spacetime continuum. Actually, in standard scenario for the simplest Higgs sector, a gauge invariance of the Lagrangian is broken when the $`\phi `$-meson fields eq.(7.5.9) acquire a VEV $`\eta _\phi 0`$ in the $`W`$-world. While the mass $`m_\phi `$ and coupling constant $`\lambda _\phi `$ are in the form eq.(7.4.15). The spontaneous breakdown of symmetry is vanished at $`\eta _\phi ^2(\lambda ,T>T_{c\mu })<0`$. When this doublet obtains a VEV, three of the gauge fields $`\underset{W}{𝑍}{}_{\mu }{}^{0},\underset{W}{𝑊}_\mu ^\pm `$ acquire masses. These fields are the $`W`$-components of the mesons mediating the weak interactions. Certainly, the derivative
$$\underset{W}{𝐷}{}_{\mu }{}^{}\phi (\underset{W}{}{}_{\mu }{}^{}\frac{i}{2}g\tau \underset{𝐖}{𝐖}{}_{\mu }{}^{}\frac{i}{2}g^{}\underset{W}{𝑋}{}_{\mu }{}^{})\phi $$
arisen in the eq.(8.1) leads to the masses
$$M_W=\frac{g\eta _\phi }{2},M_Z^2=\frac{(g^2+g^{}{}_{}{}^{2})^{1/2}\eta _\phi }{2},\mathrm{cos}\theta _W=\frac{M_W}{M_Z},$$
respectively of the gauge field components
$$\underset{W}{𝑊}{}_{\mu }{}^{\pm }=\frac{1}{\sqrt{2}}(\underset{W}{𝑊}{}_{\mu }{}^{1}\pm \underset{W}{𝑊}{}_{\mu }{}^{2}),\underset{W}{𝑍}{}_{\mu }{}^{}=\frac{g\underset{W}{𝑊}{}_{\mu }{}^{3}g^{}\underset{W}{𝑋}_\mu }{(g^2+g^{}{}_{}{}^{2})^{1/2}}\mathrm{cos}\theta _W\underset{W}{𝑊}{}_{\mu }{}^{3}\mathrm{sin}\theta _W\underset{W}{𝑋}{}_{\mu }{}^{}.$$
Consequently, a massless gauge field
$$\underset{W}{𝐴}{}_{\mu }{}^{}=\frac{g^{}\underset{W}{𝑊}{}_{\mu }{}^{3}+g\underset{W}{𝑋}_\mu }{(g^2+g^{}{}_{}{}^{2})^{1/2}}\mathrm{sin}\theta _W\underset{W}{𝑊}{}_{\mu }{}^{3}+\mathrm{cos}\theta _W\underset{W}{𝑋}{}_{\mu }{}^{}.$$
may be identified as the $`W`$-component of the photon field coupled to the electric current. Therewith the $`x`$-components of the fields above simultaneously acquire corresponding masses too, since, according to the specific MW scheme (eq.(6.6.7)), all the components of corresponding frame fields are defined on the MW mass shells, i.e.,
$$\underset{x}{}\underset{x}{𝑊}{}_{\mu }{}^{}=M_W^2\underset{x}{𝑊}{}_{\mu }{}^{},\underset{x}{}\underset{x}{𝑍}{}_{\mu }{}^{}=M_Z^2\underset{x}{𝑍}{}_{\mu }{}^{},\underset{x}{}\underset{x}{𝐴}{}_{\mu }{}^{}=M_A^2\underset{x}{𝐴}{}_{\mu }{}^{},$$
provided by
$$M_W^2\underset{W}{𝑊}{}_{\mu }{}^{}\underset{W}{}\underset{W}{𝑊}{}_{\mu }{}^{},M_Z^2\underset{W}{𝑍}{}_{\mu }{}^{}\underset{W}{}\underset{W}{𝑍}{}_{\mu }{}^{},M_A^2\underset{W}{𝐴}{}_{\mu }{}^{}\underset{W}{}\underset{W}{𝐴}{}_{\mu }{}^{}=0.$$
The microscopic structure of these fields reads
$$\begin{array}{c}W^+\varphi {}_{W}{}^{}(\eta )(q_1q_2q_3)^Q(q\overline{q})^W,W^{}\varphi {}_{W}{}^{}(\eta )(\overline{q_1q_2q_3})^Q(\overline{q}q)^W,\hfill \\ \\ Z^0\varphi {}_{Z}{}^{}(\eta )(q\overline{q})^Q(q\overline{q})^W,A\varphi {}_{A}{}^{}(\eta )(q\overline{q})^Q\underset{W}{𝐴}(0).\hfill \end{array}$$
The values of the masses $`M_W`$ and $`M_Z`$ are changed if the Higgs sector is built up more compoundly. Due to Yukawa couplings the fermions acquire the masses after symmetry-breaking. The mass of electron reads $`m_e={\displaystyle \frac{\eta _\phi }{\sqrt{2}}}f_e`$ etc. One gets for the leptons $`f_e:f_\mu :f_\tau =m_e:m_\mu :m_\tau .`$ This mechanism does not disturb the renormalizability of the theory . In approximation to lowest order $`f=\mathrm{\Sigma }_Qm_Q\lambda ^{1/2}\left(\lambda ^1={\displaystyle \frac{mp_0}{2\pi ^2}}\mathrm{ln}{\displaystyle \frac{2\stackrel{~}{\omega }}{\mathrm{\Delta }_0}}\right),`$ the Lagrangian eq.(8.1) produce the Lagrangian of phenomenological SM, where at $`f10^6`$ one gets $`\lambda 10^{12}`$.
## 9 The Two Solid Phenomenological Implications of the MSM
Discussing now the relevance of our present approach to the physical realities we should attempt to provide some ground for checking the predictions of the MSM against experimental evidence. It is remarkable that the resulting theory makes plausible following testable implications for the current experiments at LEP2, at the Tevatron and LHC discussed below, which are drastically different from the predictions of conventional models:
1. Due to the specific mechanism of the electroweak symmetry breakdown given in the previous subsection, the first important phenomenological implication of the MSM ensued at once that:
$``$ the Higgs bosons never could emerge in spacetime continuum and, thus, could not be discovered in experiments nor at any energy range.
2. According to previous subsection, the lowest pole $`m_Q`$ of the self-energy operator $`\mathrm{\Sigma }_Q`$ in eq.(8.1) has fixed the whole mass spectrum of the SM particles. But, in general, one must also take into account the mass spectrum of expected various collective excitations of bound quasi-particle pairs produced by higher-order interactions as a “superconductive” solution obtained from a nonlinear spinor field Lagrangian of the $`Q`$-component possessed $`\gamma _5`$ invariance (subsec.6.7). These states must be considered as a direct effect of the same primary nonlinear fermion interaction which provides the mass of the $`Q`$-component of Fermi field, which itself is a collective effect. They would manifest themselves as stable or unstable states. The general features of mass spectrum of the collective excitations and their coupling with the fermions are discussed in through the use of the Bethe-Salpeter equation handled in the simplest ladder approximation incorporated with the self-consistency conditions, when one is still left with unresolved divergence problem. One can reasonably expect that these results for the bosons of small masses at low energy compared to the unbound fermion states are essentially correct in spite of the very simple approximations. Therein, some bound states are predicted too the obtained mass values of which are rather high, and these states should decay very quickly. The high-energy poles may in turn determine the low-energy resonances. All this prompt us to expect that the other poles different from those of lowest one in turn will produce the new heavy SM family partners. Hence, one would expect a second important phenomenological implication of the MSM that:
$``$ for each of the three SM families of quarks and leptons there are corresponding heavy family partners with the same $`SU(3)_cSU(2)_LU(1)_Y`$ quantum numbers at the energy scales related to next poles with respect to lowest one.
To see its nature, now we may estimate the energy threshold of creation of such heavy family partners using the results far obtained in . It is therefore necessary under the simplifying assumption to consider in the $`Q`$-world a composite system of dressed fermion $`(N_{})`$ made of the unbound fermion $`(N)`$ coupled with the different kind two-fermion bound states $`(N\overline{N})`$ at low energy, which all together represent the primary manifestation of the fundamental interaction. Such a dressed fermion would have a total mass $`m_{}m_Q+\mu `$, where $`m_Q`$ and $`\mu `$ are the masses, respectively, of the unbound fermion and the bound state. According to the general discussion of the mass spectrum of the collective excitations given in , here we are interested only in the following low-energy bound states written explicitly in spectroscopic notation $`({}_{}{}^{1}S_{0}^{}){}_{N=0}{}^{},({}_{}{}^{1}S_{0}^{}){}_{N=\pm 2}{}^{},({}_{}{}^{3}P_{1}^{})_{N=0}`$ and $`({}_{}{}^{3}P_{0}^{})_{N=0}`$ with the expected masses $`\mu =0,>\sqrt{2}m_Q,\sqrt{{\displaystyle \frac{8}{3}}}m_Q`$ and $`2m_Q`$, respectively, where the subscript $`N`$ indicates the nucleon number. One notes the peculiar symmetry existing between the pseudoscalar and the scalar states that the first has zero mass and binding energy $`2m_Q`$, while the opposite holds for the scalar state. When the next pole $`m_{}`$ to the lowest one $`m_Q`$ will be switched on, then due to the Yukawa couplings in the eq.(8.1) the all fermions will acquire the new masses with their common shift $`{\displaystyle \frac{m_{}}{m_Q}}1+k`$ held upwards along the energy scale. To fix the energy threshold value all we have to do then is choose the heaviest member among the SM fermions, which is the top quark, and to set up the quite obvious formula
$$EE_0m_t^{}c^2=(1+k)m_tc^2,$$
where $`m_t`$ is the mass of the top quark. The top quark observed firstly in the two FNAL $`p\overline{p}`$ collider experiments in 1995, has the mass turned out to be startlingly large $`m_t=(173.8\pm 5.0)GeV/c^2`$ compared to all the other SM fermion masses . Thus, we get the most important energy threshold scale estimate where the heavy partners of the SM extra families of quarks and leptons likely would reside at: $`E_1>(419.6\pm 12.0)GeV,E_2=(457.6\pm 13.2)GeV`$ and $`E_3=(521.4\pm 15.0)GeV,`$ corresponded to the next nontrivial poles are written: $`k_1>\sqrt{2},k_2=\sqrt{8/3}`$ and $`k_3=2`$, respectively. We recognize well that the general results obtained in , however, model-dependent and may be considerably altered, especially in the high energy range by using better approximation. In present state of the theory it seemed to be a bit premature to get exact high energy results, which will be important subject for the future investigations. But, in the same time we believe that the approximation used in has enough accuracy for the low-energy estimate made above . Anyhow, it is for one thing, the new scale where the family partners reside will be much higher than the electroweak scale and thus these heavy partners lie far above the electroweak scale.
## 10 Quark Flavour Mixing and the Cabibbo Angles
An implication of quark generations into general scheme will be carried out in the same way of the leptons. But before proceeding further that it is profitable to enlarge it by the additional assumption without asking the reason behind it:
$``$ The MW components imply
$${}_{}{}^{i}\overline{\underset{u}{\Psi }}_{}^{A}(\mathrm{},\theta _{i_1},\mathrm{}\theta _{i_n},\mathrm{}){}_{}{}^{j}\underset{u}{\Psi }_{}^{B}(\mathrm{},\theta _{i_1},\mathrm{}\theta _{i_n},\mathrm{})=\delta _{ij}\underset{l=i_1,\mathrm{},i_n}{}f_{il}^{AB}{}_{}{}^{i}\left(\overline{q}_lq_l\right),$$
(10.1)
namely, the contribution of each individual subquark $`{}_{}{}^{i}q_{l}^{}`$, into the component of given world ($`i`$) is determined by the partial formfactor $`f_{il}^{AB}`$. Under the group $`SU(2)U(1)`$ the left-handed quarks transform as three doublets, while the right-handed quarks transform as independent singlets except of following differences:
1. The values of weak-hypercharge of quarks are changed due to their fractional electric charges $`q_L:Y^w={\displaystyle \frac{1}{3}},u_R:Y^w={\displaystyle \frac{4}{3}},d_R:Y^w={\displaystyle \frac{2}{3}},`$ etc.
2. All Yukawa coupling constants have nonzero values.
3. An appearance of quark mixing and Cabibbo angle, which is unknown in the scope of standard model.
4. An existence of CP-violating phase in unitary matrix of quark mixing. We shall discuss it in the next section.
Below, we attempt to give an explanation to quark mixing and Cabibbo angle. Here for simplicity, we consider this problem on the example of four quarks $`u,d,s,c`$. The further implication of all quarks would complicate the problem only in algebraic sense. Let consider four left-handed quarks forming a $`SU(2)_L`$ doublets mixed with Cabibbo angle $`\left(\begin{array}{c}u^{}\\ d\end{array}\right)_L,\text{and}\left(\begin{array}{c}c^{}\\ s\end{array}\right)_L,`$ where $`u^{}=u\mathrm{cos}\theta +c\mathrm{sin}\theta ,c^{}=u\mathrm{sin}\theta +c\mathrm{cos}\theta `$. One must distinguish two kind of fermion states: an eigenstate of gauge interactions, i.e. the fields of $`u^{}`$ and $`c^{}`$; an eigenstate of mass-matrices, i.e the fields of $`u`$ and $`c`$. The qualitative properties of Cabibbo flavour mixing could be understood in terms of Yukawa couplings. Unlike the case of leptons, where the Yukawa couplings are characterized by two constants $`f_e`$ and $`f_\mu `$, the interaction of Higgs boson with $`u^{}`$ and $`c^{}`$ is due to following three terms:
$$\begin{array}{c}\frac{1}{\sqrt{2}}f_u^{}(\overline{u}_L^{}u_R^{}+\overline{u}_R^{}u_L^{})(\eta +\chi )=\frac{1}{\sqrt{2}}f_u^{}(\overline{u}^{}u^{})(\eta +\chi ),\hfill \\ \\ \frac{1}{\sqrt{2}}f_c^{}(\overline{c}_L^{}c_R^{}+\overline{c}_R^{}c_L^{})(\eta +\chi )=\frac{1}{\sqrt{2}}f_c^{}(\overline{c}^{}c^{})(\eta +\chi ),\hfill \\ \\ \frac{1}{\sqrt{2}}f_{u^{}c^{}}(\overline{c}_L^{}u_R^{}+\overline{c}_R^{}u_L^{}+\overline{u}_R^{}c_L^{}+\overline{u}_L^{}c_R^{})(\eta +\chi )=\frac{1}{\sqrt{2}}f_{u^{}c^{}}(\overline{c}^{}u^{}+\overline{u}^{}c^{})(\eta +\chi ).\hfill \end{array}$$
Our discussion here throughout a small part will be a standard one \[97-101\], except, instead of mixing of fields $`d^{}`$ and $`s^{}`$ we consider a quite equivalent mixing of $`u^{}`$ and $`c^{}`$. The last expression may be diagonalized by means of rotation right through Cabibbo angle. In the sequel one gets $`m_u\overline{u}u+m_c\overline{c}c,`$ where $`m_u`$ and $`m_c`$ are masses of quarks $`u`$ and $`c`$. Straightforward comparison of two states gives
$$\begin{array}{c}m_u=\frac{1}{\sqrt{2}}(f_u^{}\mathrm{cos}^2\theta +f_c^{}\mathrm{sin}^2\theta 2f_{u^{}c^{}}\mathrm{cos}\theta \mathrm{sin}\theta )\eta ,\hfill \\ \\ m_c=\frac{1}{\sqrt{2}}(f_u^{}\mathrm{sin}^2\theta +f_c^{}\mathrm{cos}^2\theta +2f_{u^{}c^{}}\mathrm{cos}\theta \mathrm{sin}\theta )\eta ,\hfill \\ \\ \mathrm{tan}2\theta =\frac{2f_{u^{}c^{}}}{f_c^{}f_u^{}}0.\hfill \end{array}$$
Similar formulas can be worked out for the other mixing. Hence, the nonzero value of Cabibbo angle arises due to nonzero coupling constant $`f_{u^{}c^{}}`$. The problem is to calculate all coupling constants $`f_{u^{}c^{}}`$,$`f_{c^{}t^{}}`$, and $`f_{t^{}u^{}}`$ generating three Cabibbo angles
$$\mathrm{tan}2\theta _3=\frac{2f_{u^{}c^{}}}{f_c^{}f_u^{}},\mathrm{tan}2\theta _1=\frac{2f_{c^{}t^{}}}{f_t^{}f_c^{}},\mathrm{tan}2\theta _2=\frac{2f_{t^{}u^{}}}{f_u^{}f_t^{}}.$$
Taking into account the explicit form of Q-components of quark fields eq.(6.4.1)
$$\underset{Q}{\Psi }_u^{}=(q_2q_3)^Q,\underset{Q}{\Psi }_c^{}=(q_3q_1)^Q,\underset{Q}{\Psi }_t^{}=(q_1q_2)^Q,$$
also eq.(6.7.2) and eq.(6.10.1), we may write down
$$f_u^{}\frac{1}{2}\{\underset{u}{\overline{\mathrm{\Psi }}}_u^{}\widehat{p}_u\underset{u}{\Psi }_u^{}\left(\underset{u}{\overline{\mathrm{\Psi }}}_u^{}\stackrel{}{\widehat{p}}_u\right)\underset{u}{\Psi }_u^{}\}=(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3})\underset{u}{\overline{\mathrm{\Psi }}}_u^{}\underset{u}{\Psi }_u^{}(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3}),$$
(10.2)
where according to the eq.(6.10.7) $`\overline{\mathrm{\Sigma }}{}_{Qu}{}^{i}=tr\left(\rho _u\mathrm{\Sigma }_Q^i\right)`$, $`\rho _u`$ is given in eq.(6.10.8) and $`\widehat{p}_Qq_i^Q=\mathrm{\Sigma }_Q^iq_i^Q.`$ In analogy, the $`f_c^{}`$ and $`f_{u^{}c^{}}`$ imply
$$f_c^{}\frac{1}{2}\{\underset{u}{\overline{\mathrm{\Psi }}}_c^{}\widehat{p}_u\underset{u}{\Psi }_c^{}\left(\underset{u}{\overline{\mathrm{\Psi }}}_c^{}\stackrel{}{\widehat{p}}_u\right)\underset{u}{\Psi }_c^{}\}(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1}),$$
(10.3)
and
$$\begin{array}{c}f_{u^{}c^{}}\frac{1}{4}\left\{\underset{Q}{\overline{\mathrm{\Psi }}}_u^{}\widehat{p}_Q\underset{Q}{\Psi }_c^{}+\underset{Q}{\overline{\mathrm{\Psi }}}_c^{}\widehat{p}_Q\underset{Q}{\Psi }_u^{}\left(\underset{Q}{\overline{\mathrm{\Psi }}}_u^{}\stackrel{}{\widehat{p}}_Q\right)\underset{Q}{\Psi }_c^{}\left(\underset{Q}{\overline{\mathrm{\Psi }}}_c^{}\stackrel{}{\widehat{p}}_Q\right)\underset{Q}{\Psi }_u^{}\right\}=\hfill \\ \\ =\frac{1}{2}\{(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3})\underset{Q}{\overline{\mathrm{\Psi }}}_u^{}\underset{Q}{\Psi }_c^{}+(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})\underset{Q}{\overline{\mathrm{\Psi }}}_c^{}\underset{Q}{\Psi }_u^{}\}.\hfill \end{array}$$
(10.4)
In accordance with eq.(10.1), one has
$$\begin{array}{c}\underset{Q}{\overline{\mathrm{\Psi }}}_u^{}\underset{Q}{\Psi }_c^{}=(\overline{q_2q_3})^Q(q_3q_1)^Q=f_{Q3}^{u^{}c^{}}(\overline{q}_3q_3)^Q,\hfill \\ \\ \underset{Q}{\overline{\mathrm{\Psi }}}_c^{}\underset{Q}{\Psi }_u^{}=(\overline{q_3q_1})^Q(q_2q_3)^Q=f_{Q3}^{c^{}u^{}}(\overline{q}_3q_3)^Q,\hfill \end{array}$$
where
$$f_{Q3}^{u^{}c^{}}(\overline{q}_3q_3)^Q=\left(f_{Q3}^{u^{}c^{}}(\overline{q}_3q_3)^Q\right)^{}=f_{Q3}^{c^{}u^{}}(\overline{q}_3q_3)^Q\overline{f}_3.$$
Hence $`f_{u^{}c^{}}={\displaystyle \frac{\overline{f}_2}{3}}(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})`$ and
$$\begin{array}{c}\mathrm{tan}2\theta _3=\frac{\overline{f}_3(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})}{(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1}\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3})},\mathrm{tan}2\theta _1=\frac{\overline{f}_1(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2})}{(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2}\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})},\hfill \\ \\ \mathrm{tan}2\theta _2=\frac{\overline{f}_2(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3})}{(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3}\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2})},\hfill \end{array}$$
(10.5)
where the rest of $`\overline{f}_i`$ reads $`\overline{f}_1f_{Q1}^{c^{}t^{}}=f_{Q1}^{t^{}c^{}}`$ and $`\overline{f}_2f_{Q2}^{t^{}u^{}}=f_{Q2}^{u^{}t^{}}`$. Thus, the Q-components of the quark fields $`u^{},c^{}`$ and $`t^{}`$ contain at least one identical subquark, due to which in eq.(6.10.1) the partial formfactors $`\overline{f}_i`$ have nonzero values causing a quark mixing with the Cabibbo angles eq.(10.5). Therefore, the unimodular orthogonal group of global rotations arises, and the quarks $`u^{},c^{}`$ and $`t^{}`$ come up in doublets $`(u^{},c^{})`$,$`(c^{},t^{})`$, and $`(t^{},u^{})`$. For the leptons these formfactors equal zero $`\overline{f}_i^{lept}0`$, because of eq.(6.1.1), namely the lepton mixing is absent. In conventional notation $`\left(\begin{array}{c}u^{}\\ d\end{array}\right)_L,\left(\begin{array}{c}c^{}\\ s\end{array}\right)_L,\left(\begin{array}{c}t^{}\\ b\end{array}\right)_L\left(\begin{array}{c}u\\ d^{}\end{array}\right)_L,\left(\begin{array}{c}c\\ s^{}\end{array}\right)_L,\left(\begin{array}{c}t\\ b^{}\end{array}\right)_L,`$ which gives rise to $`f_{u^{}c^{}}f_{d^{}s^{}},f_{c^{}t^{}}f_{s^{}b^{}},f_{t^{}u^{}}f_{b^{}d^{}},f_u^{}f_d^{},f_c^{}f_s^{},f_t^{}f_b^{},f_df_u,f_sf_c,f_bf_t.`$
## 11 The Appearance of the CP-Violating Phase
The required magnitude of the CP-violating complex parameter $`\epsilon `$ depends upon the specific choice of theoretical model for explaining the $`K_2^02\pi `$ decay . From the experimental data it is somewhere $`|\epsilon |2.3\times 10^3.`$ In the framework of Kobayashi-Maskawa (KM) parametrization of unitary matrix of quark mixing , this parameter may be expressed in terms of three Eulerian angles of global rotations in the three dimensional quark space and one phase parameter. Below we attempt to derive the KM-matrix with an explanation given to an appearance of the CP-violating phase. We recall that during the realization of multiworld structure the P-violation compulsory occurred in the W-world provided by the spanning eq.(6.8.1). The three dimensional effective space $`W_v^{loc}(3)`$ arises as follows:
$$\begin{array}{c}W_v^{loc}(3)q_v^{(3)}=\left(\begin{array}{c}q_R^w(\stackrel{}{T}=0)\\ \\ q_L^w(\stackrel{}{T}=\frac{1}{2})\end{array}\right)\hfill \\ \\ \left(\begin{array}{c}u_R,d_R\\ \\ \left(\begin{array}{c}u^{}\\ d\end{array}\right)_L\end{array}\right),\left(\begin{array}{c}c_R,s_R\\ \\ \left(\begin{array}{c}c^{}\\ s\end{array}\right)_L\end{array}\right),\left(\begin{array}{c}t_R,b_R\\ \\ \left(\begin{array}{c}t^{}\\ b\end{array}\right)_L\end{array}\right)\left(\begin{array}{c}q_3^w\\ \\ \left(\begin{array}{c}q_1^w\\ q_2^w\end{array}\right)\end{array}\right),\left(\begin{array}{c}q_1^w\\ \\ \left(\begin{array}{c}q_2^w\\ q_3^w\end{array}\right)\end{array}\right),\left(\begin{array}{c}q_2^w\\ \\ \left(\begin{array}{c}q_3^w\\ q_1^w\end{array}\right)\end{array}\right),\hfill \end{array}$$
(11.1)
where the subscript $`(v)`$ formally specifies a vertical direction of multiplet, the subquarks $`q_\alpha ^w(\alpha =1,2,3)`$ associate with the local rotations around corresponding axes of three dimensional effective space $`W_v^{loc}(3)`$. The local gauge transformations $`f_{exp}^v`$ are implemented upon the multiplet $`q_{}^{}{}_{v}{}^{(3)}=f_{exp}^vq_v^{(3)}`$, where $`f_{exp}^vSU^{loc}(2)U^{loc}(1)`$. If for the moment we leave it intact and make a closer examination of the content of the middle row in eq.(11.1), then we distinguish the other symmetry arisen along the horizontal line $`(h)`$. Hence, we may expect a situation similar to those of subsec.6.8 will be held in present case. The procedure just explained therein can be followed again. We have to realize that due to the specific structure of W-world implying the condition of realization of the MW connections eq.(6.1.5) with $`\stackrel{}{T}0,Y^w0`$, the subquarks $`q_\alpha ^w`$ tend to be compulsory involved into triplet. They form one “doublet” $`\stackrel{}{T}0`$ and one singlet $`Y^w0`$. Then the quarks $`u_L^{},c_L^{}`$ and $`t_L^{}`$ form a $`SO^{gl}(2)`$ “doublet” and a $`U^{gl}(1)`$ singlet
$$\begin{array}{c}\left((u_L^{},c_L^{})t_L^{}\right)\left((q_1^w,q_2^w)q_3^w\right)q_h^{(3)}W_h^{gl}(3),\hfill \\ \\ (u_L^{},(c_L^{},t_L^{}))(q_1^w,(q_2^w,q_3^w)),\left((t_L^{},u_L^{})c_L^{}\right)((q_3^w,q_1^w),q_2^w).\hfill \end{array}$$
(11.2)
Here $`W_h^{gl}(3)`$ is the three dimensional effective space in which the global rotations occur. They are implemented upon the triplets through the transformation matrix $`f_{exp}^h`$:
$`q_{}^{}{}_{h}{}^{(3)}=f_{exp}^hq_h^{(3)},`$ which reads (eq.(11.2))
$$f_{exp}^h=\left(\begin{array}{ccc}f_{33}& 0& 0\\ 0& c& s\\ 0& s& c\end{array}\right)$$
in the notation $`c=\mathrm{cos}\theta ,s=\mathrm{sin}\theta `$. This implies the incompatibility relation eq.(4.4.5), namely
$$f_{exp}^h=f_{33}(f_{11}f_{22}f_{12}f_{21})=f_{33}\epsilon _{123}\epsilon _{123}f_{exp}^hf_{33}^{}.$$
(11.3)
That is $`f_{33}f_{33}^{}=1,`$ or $`f_{33}=e^{i\delta }`$ and $`f_{exp}^h=1`$. The general rotation in $`W_h^{gl}(3)`$ is described by Eulerian three angles $`\theta _1,\theta _2,\theta _3`$. If we put the arisen phase only in the physical sector then a final KM-matrix of quark flavour mixing would result
$$\begin{array}{c}(\overline{u}_L,\overline{c}_L,\overline{t}_L)V_{KM}\left(\begin{array}{c}d\\ s\\ b\end{array}\right)(\overline{u}_L^{},\overline{c}_L^{},\overline{t}_L^{})\left(\begin{array}{c}d\\ s\\ b\end{array}\right)(\overline{u}_L,\overline{c}_L,\overline{t}_L)\left(\begin{array}{c}d^{}\\ s^{}\\ b^{}\end{array}\right)=\hfill \\ \\ =(\overline{u}_L,\overline{c}_L,\overline{t}_L)\left(\begin{array}{ccc}1& 0& 0\\ 0& c_2& s_2\\ 0& s_2& c_2\end{array}\right)\left(\begin{array}{ccc}c_1& s_1& 0\\ s_1& c_1& 0\\ 0& 0& e^{i\delta }\end{array}\right)\left(\begin{array}{ccc}1& 0& 0\\ 0& c_3& s_3\\ 0& s_3& c_3\end{array}\right)\left(\begin{array}{c}d\\ s\\ b\end{array}\right),\hfill \end{array}$$
(11.4)
where
$$(\overline{u}_L^{},\overline{c}_L^{},\overline{t}_L^{})(\overline{u}_L,\overline{c}_L,\overline{t}_L)V_{KM},\left(\begin{array}{c}d^{}\\ s^{}\\ b^{}\end{array}\right)V_{KM}\left(\begin{array}{c}d\\ s\\ b\end{array}\right).$$
The CP-violating parameter $`\epsilon `$ approximately is written $`\epsilon s_1s_2s_3\mathrm{sin}\delta 0.`$ While the spanning $`W_v^{loc}(2)W_v^{loc}(3)`$ eq.(11.1) underlies the P-violation and the expanded symmetry $`G_v^{loc}(3)=SU^{loc}(2)U^{loc}(1)`$, the CP-violation stems from the similar spanning $`W_h^{gl}(2)W_h^{gl}(3)`$ eq.(11.2) with the expanded global symmetry group.
## 12 The Mass-Spectrum of Leptons and Quarks
The mass-spectrum of leptons and quarks stems from their internal MW-structure eq.(6.3.1) and eq.(6.4.1) incorporated with the quark mixing eq.(10.5). We start a discussion with the leptons. It might be worthwhile to adopt a simple viewpoint on Higgs sector. Following the sec.8, the explicit expressions of the lepton masses read $`m_i={\displaystyle \frac{\eta }{\sqrt{2}}}f_i`$ and $`m_i^\nu ={\displaystyle \frac{\eta }{\sqrt{2}}}f_i^\nu `$, that $`m_e:m_\mu :m_\tau =f_e:f_\mu :f_\tau =L_1^2:L_2^2:L_3^2`$ provided by $`L_1^2={\displaystyle \frac{m_i}{M}}`$ and $`\sqrt{M}={\displaystyle \underset{i}{}}\sqrt{m_i}`$. Thus, $`L_1=(8.9;7.8)\times 10^3,L_2=0.13;0.11,L_3=0.9;0.88`$.
Taking into account the eq.(6.1.6) and eq.(6.10.7) the coupling constants of the quarks $`d,s`$ and $`b`$ can be written
$$\begin{array}{c}f_d=L_1tr(\rho _d\mathrm{\Sigma }_Q)L_1\stackrel{~}{f}_d,f_s=L_2tr(\rho _s\mathrm{\Sigma }_Q)L_2\stackrel{~}{f}_s,f_b=L_3tr(\rho _b\mathrm{\Sigma }_Q)L_3\stackrel{~}{f}_b,\hfill \\ \\ \rho _d=\rho ^Q\rho _d^B,\rho _s=\rho ^Q\rho _s^B\rho ^s,\rho _b=\rho ^Q\rho _b^B\rho ^b.\hfill \end{array}$$
(12.1)
Hence $`m_d={\displaystyle \frac{\eta }{\sqrt{2}}}f_d,m_s={\displaystyle \frac{\eta }{\sqrt{2}}}f_s,m_b={\displaystyle \frac{\eta }{\sqrt{2}}}f_b,`$ and $`m_d:m_s:m_b=(L_1\stackrel{~}{f}_d):(L_2\stackrel{~}{f}_s):(L_3\stackrel{~}{f}_b).`$ According to the subsec.6.10, we derive the masses of the $`u,c`$ and $`t`$ quarks
$$\begin{array}{c}m_u=\frac{\eta }{\sqrt{2}}\{(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3})\mathrm{cos}^2\theta _3+(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})\mathrm{sin}^2\theta _3\frac{\overline{f}_3}{2}(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3}+\hfill \\ \\ \overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})\mathrm{sin}2\theta _3\}=\frac{\eta }{\sqrt{2}}\{(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3})\mathrm{cos}^2\theta _2+(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2})\mathrm{sin}^2\theta _2+\hfill \\ \\ \frac{\overline{f}_2}{2}\left(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3}+\right)\mathrm{sin}2\theta _2\},\hfill \end{array}.$$
(12.2)
$$\begin{array}{c}m_c=\frac{\eta }{\sqrt{2}}\{(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3})\mathrm{sin}^2\theta _3+(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})\mathrm{cos}^2\theta _3+\frac{\overline{f}_3}{2}(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3}+\hfill \\ \\ \overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})\mathrm{sin}2\theta _3\}=\frac{\eta }{\sqrt{2}}\{(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})\mathrm{cos}^2\theta _1+(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2})\mathrm{sin}^2\theta _1\hfill \\ \\ \frac{\overline{f}_1}{2}\left(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2}+\right)\mathrm{sin}2\theta _1\},\hfill \end{array}.$$
(12.3)
$$\begin{array}{c}m_t=\frac{\eta }{\sqrt{2}}\{(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2})\mathrm{cos}^2\theta _1+(\overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})\mathrm{sin}^2\theta _1+\frac{\overline{f}_1}{2}(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2}+\hfill \\ \\ \overline{\mathrm{\Sigma }}{}_{Qc}{}^{3}+\overline{\mathrm{\Sigma }}{}_{Qc}{}^{1})\mathrm{sin}2\theta _1\}=\frac{\eta }{\sqrt{2}}\{(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2})\mathrm{cos}^2\theta _2+(\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3})\mathrm{sin}^2\theta _2\hfill \\ \\ \frac{\overline{f}_2}{2}\left(\overline{\mathrm{\Sigma }}{}_{Qt}{}^{1}+\overline{\mathrm{\Sigma }}{}_{Qt}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{2}+\overline{\mathrm{\Sigma }}{}_{Qu}{}^{3}+\right)\mathrm{sin}2\theta _2\}.\hfill \end{array}.$$
(12.4)
## 13 The Physical Outlook and Concluding Remarks
In this section we briefly expose the main features of our physical outlook and draw a number of conclusions.
Our purpose in the part I is to develop the OM formalism, which is the mathematical framework for our physical outlook embodied in the idea that the geometry and fields, with the internal symmetries and all interactions, as well the four major principles of relativity (special and general), quantum, gauge and colour confinement, are derivative. They come into being simultaneously in the stable system of the underlying “primordial structures” involved in the “linkage” establishing processes (sec.4). The OM formalism is the generalization of secondary quantization of the field theory with appropriate expansion over the geometric objects leading to the quantization of geometry different from all existing schemes. Below, once again, we resume its relevant steps and major points:
$``$ 1. An extension of the four-dimensional Minkowski space $`M_4`$ by the additional one sample of the internal world (a simplified case of the one $`u`$-channel) in order to introduce the mass operator of the fields defined on the internal world (sec.2). 2. A two-step passage for each sample of the $`M_4G_6`$ (sec.2), which restores the complete equivalence between the three spacial and the three time components with the subsequent rotation of the basis vectors on the $`45^0`$ angle providing an adequate algebra for geometry quantization.
3. We deal in terms of first degree of the line element, which entails an additional phase multiplier for vectors (the origin of the fields). All this leads to the simplified scheme of smooth 12-dimensional manifold $`G`$ (sec.2-4). The passage back to the $`M_4`$ may be performed whenever it will be necessary (subsec.2.1).
$``$ Building up the OM formalism we proceed at once to the quantization of geometry by substituting the basis vectors for the creation and annihilation operators acting in the configuration space of occupation numbers (subsec.2.1). They include also the Pauli’s matrices, due to which all the states defined on the OM are degenerate at the very outset with degeneracy degree equal 2. It implies the half-spin quantum number, which subsequently gives rise to the spins of the particles. This rule for spin quantum number is not without an important reason. The argument for this conclusion is compulsory suggested by the properties of Pauli’s operators (subsec.2.1).
$``$ In this framework we derive the matrix element of the line element eq.(2.2.2), which gives the most important relation for the realization of the $`G`$.
$``$ Based on configuration space mechanics with antisymmetric state functions, we discuss in detail the quantum field and differential geometric aspects of the OM (subsec.2.1, App.A).
$``$ We have chosen a simple setting and considered the primordial structures, which are designed to posses certain physical properties satisfying the stated in subsec.4.1 general rules and have involved in the linkage establishing processes. The processes of their creation and annihilation in the lowest state (the regular structures) just are described by the OM formalism. In all the higher states the primordial structures are distorted ones, namely they have undergone the distortion transformations (subsec.4.2). These transformations yield the “quark” and “antiquark” fields defined on the simplified geometry (one $`u`$-channel) given in the subsec.4.3, and skeletonized for illustrative purposes. Due to geometry realization conditions held in the stable systems of primordial structures they emerge in confined phase (subsec.4.2). This scheme still should be considered as the preliminary one, which is further elaborated in the subsec.5.3 to get the physically more realistic picture.
$``$ The distortion transformation functions are the operators acting in the space of the internal degrees of freedom (colours) and imply the incompatibility relations eq.(4.4.5), which hold for both the local and the global distortion rotations. They underly the most important symmetries such as the internal symmetries $`U(1),SU(2),SU(3)`$, the $`SU(2)U(1)`$ symmetry of electroweak interactions, etc., (see part II).
$``$ We generalize the OM formalism via the concept of the OMM yielding the MW geometry involving the spacetime continuum and the internal worlds of the given number (sec.5). In an enlarged framework of the OMM we define and clarify the conceptual basis of subquarks and their characteristics stemming from the various symmetries of the internal worlds (subsec.5.3). They imply subcolour confinement and gauge principle.
By this we have arrived at an entirely satisfactory answer to the question of the physical origin of the geometry and fields, the internal symmetries and interactions, as well the principles of relativity, quantum, gauge and subcolour confinement.
The value of the present version of hypothesis of existence of the MW structures resides in solving in part II some key problems of the SM, wherein we attempt to suggest a microscopic approach to the properties of particles and interactions.
$``$ Within this approach the fields have composite nontrivial internal structure (sec.6). The condition of realization of the MW connections is arisen due to the symmetry of Q-world of electric charge and embodied in the Gell-Mann-Nishijima relation (subsec.6.1). During the realization of the MW-structure the symmetries of corresponding internal worlds are unified into more higher symmetry including also the operators of isospin and hypercharge. Such approach enables to conclude that only possible at low-energy the three lepton generations consist of six lepton fields with integer electric and leptonic charges and being free of confinement (subsec.6.3). Also the three quark generations exist composed of six possible quark fields. They carry fractional electric and baryonic charges and obey confinement condition (subsec.6.4). The global group unifying all global symmetries of the internal worlds of quarks is the flavour group $`SU_f(6)`$ (subsec.6.5). The whole complexity of leptons, quarks and other composite particles, and their interactions arises from the primary field, which has nontrivial MW internal structure and involves nonlinear fermion self-interaction of the components (subsec.6.6). This Lagrangian contains only two free parameters, which are the coupling constants of nonlinear fermion and gauge interactions. Due to specific structure of the W-world of weak interactions implying the condition of realization of the MW connections, the spanning eq.(6.8.1) takes place, which underlies the P-violation in W-world. It is expressed in the reduction of initial symmetry of the right-handed subquarks. Such reduction is characterized by the Weinberg mixing angle with the value fixed at $`30^0`$ (subsec.6.9). It gives rise to the expanded local symmetry $`SU(2)U(1)`$, under which the left-handed fermions transform as six independent $`SU(2)`$ doublets, while the right-handed fermions transform as twelve independent singlets (subsec.6.8). Due to vacuum rearrangement in Q-world the Yukawa couplings arise between the fermion fields and corresponding isospinor-scalar $`\phi `$-meson in conventional form (subsec.6.10).
$``$ We suggest the microscopic approach to Higgs bosons with self-interaction and Yukawa couplings (sec.7). It involves the Higgs bosons as the collective excitations of bound quasi-particle iso-pairs. In the framework of local gauge invariance of the theory incorporated with the P-violation in weak interactions we propose a mechanism providing the Bose-condensation of iso-pairs, which is due to effective attraction between the relativistic fermions caused by the exchange of the mediating induced gauge quanta in the W-world. We consider the four-component Bose-condensate, where due to self-interaction its spin part is vanished. Based on it we show that the field of symmetry-breaking Higgs boson always must be counted off from the gap symmetry restoring value as the point of origin. Then the Higgs boson describes the excitations in the neighbourhood of stable vacuum of the W-world.
$``$ In contrast to the SM, the suggested approach predicts the electroweak symmetry breakdown in the $`W`$-world by the VEV of spin zero Higgs bosons and the transmission of electroweak symmetry breaking from the $`W`$world to the $`M_4`$ spacetime continuum (sec.8). The resulting Lagrangian of unified electroweak interactions of leptons and quarks ensues, which in lowest order approximation leads to the Lagrangian of phenomenological SM. In general, the self-energy operator underlies the Yukawa coupling constant, which takes into account a mass-spectrum of all expected collective excitations of bound quasi-particle pairs.
$``$ If the MSM proves viable it becomes an crucial issue to hold in experiments the two testable solid implications given in sec.9, which are drastically different from those of conventional models.
$``$ The implication of quarks into this scheme is carried out in the same manner except that of appearance of quark mixing with Cabibbo angles and the existence of CP-violating complex phase in unitary matrix of quark mixing. The Q-components of the quarks $`u^{},c^{}`$ and $`t^{}`$ contain at least one identical subquark, due to which the partial formfactors gain nonzero values. This underlies the quark mixing with Cabibbo angles (sec.10). In the case of the leptons these formfactors are vanished and the mixing is absent. The CP-violation stems from the spanning eq.(11.2) incorporated with the expanded group of global rotations. With a simple viewpoint on Higgs sector the masses of leptons and quarks are given in sec.12.
Acknowledgements
I am pleased to mention the most valuable discussions with the late V.H.Ambartsumian on the various issues treated in this paper. I wish to thank S.P.Novikov for useful discussions of some points of mathematical framework, especially, of the reflection formalism (App. B) I express my gratitude to G.Jona-Lasinio for fruitful comments and suggestions. This work was supported in part by the program of Jumelage. It is pleasure to thank for their hospitality the “Centre de Physique Theorique” (CNRS Division 7061, Marseille) and its director P.Chiappetta, and ’D.A.R.C.(Observatoir de Paris-CNRS). I am grateful to R.Triay for comments and suggestions, and acknowledge useful discussions with H.Fliche, G.Sigle and M.Lemoine as well as all the participants of the seminars. I’m indebted to A.M.Vardanian and K.L.Yerknapetian for support.
## Appendix A
## 1 Mathematical Background
$``$Field Aspect of the OM
The free state of $`i`$-type fermion with definite values of momentum $`p_i`$ and spin projection $`s`$ is described by means of plane waves ($`\mathrm{}=1,c=1`$): $`\underset{\eta }{\Psi }{}_{p_\eta }{}^{}(\eta )=\left({\displaystyle \frac{m}{E_\eta }}\right)^{1/2}\underset{\eta }{𝑈}(p_\eta ,s)e^{ip_\eta \eta }`$, etc, where $`E_i\underset{i}{𝑝}{}_{0}{}^{}=\stackrel{}{\underset{i}{𝑝}}{}_{0}{}^{},\underset{i}{𝑝}{}_{0\alpha }{}^{}={\displaystyle \frac{1}{\sqrt{2}}}(\underset{i}{𝑝}{}_{(+\alpha )}{}^{}+\underset{i}{𝑝}{}_{(\alpha )}{}^{}),\stackrel{}{\underset{i}{𝑝}}={\displaystyle \frac{1}{\sqrt{2}}}(\stackrel{}{\underset{i}{𝑝}}{}_{+}{}^{}\stackrel{}{\underset{i}{𝑝}}{}_{}{}^{}),p_\eta ^2=E_\eta ^2\stackrel{}{p}_\eta ^2=p_u^2=E_u^2\stackrel{}{p}_u^2=m^2.`$
Suppose that the $`i`$-th fermion is found in the state $`r_i`$ with the vector function $`\mathrm{\Phi }_{r_i}^{(\lambda _i,\mu _i,\alpha _i)}=\zeta _{r_i}^{(\lambda _i,\mu _i,\alpha _i)}\mathrm{\Phi }_{r_i}^{\lambda _i,\mu _i}(\zeta _{r_i})`$ and
$`\zeta _{r_i^{\lambda _i,\mu _i}}={\displaystyle \underset{\alpha _i=1}{\overset{3}{}}}e_{(\lambda _i,\mu _i,\alpha _i)}^{r_i}\zeta _{r_i}^{(\lambda _i,\mu _i,\alpha _i)}`$, $`\zeta _{r_i}={\displaystyle \underset{\lambda _i,\mu _i=1}{\overset{2}{}}}\zeta _{r_i^{\lambda _i,\mu _i}}\stackrel{~}{𝒰}_{r_i}`$,
the $`\stackrel{~}{𝒰}_{r_i}`$ is the open neighbourhood of the point $`\zeta _{r_i};`$ the $`r_i`$ implies a set $`(r_i^{11},r_i^{12},r_i^{21},r_i^{22})`$. Let the $`^{(1)}`$ is a Hilbert space used for quantum mechanical description of one particle, namely $`^{(1)}`$ is a finite or infinite dimensional complex space provided with scalar product $`(\mathrm{\Phi },\mathrm{\Psi }),`$ which is linear with respect to $`\mathrm{\Psi }`$ and antilinear to $`\mathrm{\Phi }`$. The $`^{(1)}`$ is complete in norm $`|\mathrm{\Phi }|=(\mathrm{\Phi },\mathrm{\Phi })^{1/2}`$, i.e. each fundamental sequence $`\{\mathrm{\Phi }_n\}`$ of vectors of the $`^{(1)}`$ has converged by norm on $`^{(1)}`$. One particle state function is written $`\mathrm{\Phi }_{r_i}^{(1)}={\displaystyle \underset{\lambda _i,\mu _i=1}{\overset{2}{}}}\mathrm{\Phi }_{r_i^{\lambda _i,\mu _i}}^{(1)}_{r_i}^{(1)}`$, where $`_{r_i}^{(1)}={\displaystyle \underset{\lambda _i,\mu _i=1}{\overset{2}{}}}_{r_i^{\lambda _i,\mu _i}}^{(1)}`$. Define
$$\stackrel{~}{\mathrm{\Phi }}^{(1)}=\zeta _i\mathrm{\Phi }_{r_i}^{(1)}\stackrel{~}{G}_{r_i}^{(1)}=\stackrel{~}{𝒰}_{r_i}^{(1)}_{r_i}^{(1)}.$$
(A.1.1)
For description of n particle system we introduce Hilbert space
$$\overline{}_{(r_1,\mathrm{},r_n)}^{(n)}=_{r_1}^{(1)}\mathrm{}_{r_n}^{(1)}$$
(A.1.2)
by considering all sequences
$$\mathrm{\Phi }_{(r_1,\mathrm{},r_n)}^{(n)}=\{\mathrm{\Phi }_{r_1}^{(1)},\mathrm{},\mathrm{\Phi }_{r_n}^{(1)}\}=\mathrm{\Phi }_{r_1}^{(1)}\mathrm{}\mathrm{\Phi }_{r_n}^{(1)},$$
(A.1.3)
where $`\mathrm{\Phi }_{r_i}^{(1)}_{r_i}^{(1)}`$ provided, as usual, with the scalar product
$$(\mathrm{\Phi }_{(r_1,\mathrm{},r_n)}^{(n)},\mathrm{\Psi }_{(r_1,\mathrm{},r_n)}^{(n)})=\underset{i=1}{\overset{n}{}}(\mathrm{\Phi }_{r_i}^{(1)},\mathrm{\Psi }_{r_i}^{(1)}).$$
(A.1.4)
We consider the space $`_{(r_1,\mathrm{},r_n)}^{(n)}`$ of all the limited linear combinations of eq.(A.1.2) and continue by linearity the scalar product eq.(A.1..4) on the $`_{(r_1,\mathrm{},r_n)}^{(n)}`$. The wave function $`\mathrm{\Phi }_{(r_1,\mathrm{},r_n)}^{(n)}_{(r_1,\mathrm{},r_n)}^{(n)}`$ must be antisymmetrized over its arguments. We distinguish the antisymmetric part $`{}_{}{}^{A}\overline{}_{}^{(n)}`$ of Hilbert space $`\overline{}^{(n)}`$ by considering the functions
$${}_{}{}^{A}\mathrm{\Phi }_{(r_1,\mathrm{},r_n)}^{(n)}=\frac{1}{\sqrt{n!}}\underset{\sigma S(n)}{}sgn(\sigma )\mathrm{\Phi }_{\sigma (r_1,\mathrm{},r_n)}^{(n)}.$$
(A.1.5)
The summation is extended over all permutations of indices $`(r_1^{\lambda \mu },\mathrm{},r_n^{\lambda \mu })`$ of the integers $`1,2,\mathrm{},n,`$ where the antisymmetrical eigenfunctions are sums of the same terms with alternating signs in dependence of a parity $`sgn(\sigma )`$ of transposition. One continues the reflection $`\mathrm{\Phi }^{(n)}{}_{}{}^{A}\mathrm{\Phi }_{}^{(n)}`$ by linearity on the $`^{(n)}`$, which is limited and enables the expansion by linearity on the $`{}_{}{}^{A}\overline{}_{}^{(n)}`$. The region of values of this reflection is the $`{}_{}{}^{A}\overline{}_{}^{(n)}`$, namely an antisymmetrized tensor product of $`n`$ identical samples of $`^{(1)}`$. We introduce
$$\begin{array}{c}{}_{}{}^{A}\stackrel{~}{\mathrm{\Phi }}_{(r_1,\mathrm{},r_n)}^{(n)}=\frac{1}{\sqrt{n!}}\underset{\sigma S(n)}{}sgn(\sigma )\stackrel{~}{\mathrm{\Phi }}_{\sigma (r_1,\mathrm{},r_n)}^{(n)}=\hfill \\ =\frac{1}{\sqrt{n!}}\underset{\sigma S(n)}{}sgn(\sigma )\stackrel{~}{\mathrm{\Phi }}_{r_1}^{(1)}\mathrm{}\stackrel{~}{\mathrm{\Phi }}_{r_n}^{(1)}{}_{}{}^{A}\stackrel{~}{G}_{(r_1,\mathrm{},r_n)}^{(n)}=\stackrel{~}{𝒰}_{(r_1,\mathrm{},r_n)}^{(n)}{}_{}{}^{A}\widehat{}_{(r_1,\mathrm{},r_n)}^{(n)}.\hfill \end{array}$$
(A.1.6)
and consider a set $`{}_{}{}^{A}\stackrel{~}{}`$ of all sequences $`{}_{}{}^{A}\stackrel{~}{\mathrm{\Phi }}=\{{}_{}{}^{A}\stackrel{~}{\mathrm{\Phi }}_{}^{(0)},{}_{}{}^{A}\stackrel{~}{\mathrm{\Phi }}_{}^{(1)}\mathrm{},{}_{}{}^{A}\stackrel{~}{\mathrm{\Phi }}_{}^{(n)}\mathrm{}\},`$ with a finite number of nonzero elements. Therewith, the set $`{}_{}{}^{A}q:{}_{}{}^{A}\mathrm{\Phi }=\{{}_{}{}^{A}\mathrm{\Phi }_{}^{(0)},{}_{}{}^{A}\mathrm{\Phi }_{}^{(1)}\mathrm{},{}_{}{}^{A}\mathrm{\Phi }_{}^{(n)}\mathrm{}\}`$ is provided by the structure of the Hilbert subspace implying the composition rules
$$\begin{array}{c}{}_{}{}^{A}(\lambda \mathrm{\Phi }+\mu \mathrm{\Psi })_{}^{(n)}=\lambda {}_{}{}^{A}\mathrm{\Phi }_{}^{(n)}+\mu {}_{}{}^{A}\mathrm{\Psi }_{}^{(n)},\lambda ,\mu C,\hfill \\ ({}_{}{}^{A}\mathrm{\Phi },{}_{}{}^{A}\mathrm{\Psi })=\underset{n=0}{\overset{\mathrm{}}{}}({}_{}{}^{A}\mathrm{\Phi }_{}^{(n)},{}_{}{}^{A}\mathrm{\Psi }_{}^{(n)}).\hfill \end{array}$$
(A.1.7)
The wave manifold $`𝒢`$ stems from the $`{}_{}{}^{A}\stackrel{~}{}`$ because of the expansion by metric induced as a scalar product on $`{}_{}{}^{A}`$
$$𝒢=\underset{n=0}{\overset{\mathrm{}}{}}𝒢^{(n)}=\underset{n=0}{\overset{\mathrm{}}{}}\left(\stackrel{~}{𝒰}^{(n)}{}_{}{}^{A}\overline{}_{}^{(n)}\right).$$
(A.1.8)
The creation $`\widehat{\gamma }_r`$ and annihilation $`\widehat{\gamma }^r`$ operators for each $`^{(1)}`$ can be defined as follows: one must modify the basis operators in order to provide an anticommutation in arbitrary states
$$\widehat{\gamma }{}_{(\lambda ,\mu ,\alpha )}{}^{r}\widehat{\gamma }{}_{(\lambda ,\mu ,\alpha )}{}^{r}\eta _{r}^{\lambda \mu },(\eta _r^{\lambda \mu })^+=\eta _r^{\lambda \mu },$$
(A.1.9)
for given $`\lambda ,\mu ,\alpha `$, where $`\eta _r`$ is a diagonal operator in the space of occupation numbers, while, at $`r_i<r_j`$ one gets $`\widehat{\gamma }{}_{}{}^{r_i}\eta _{r_j}^{}=\eta _{r_j}\widehat{\gamma }{}_{}{}^{r_i},\widehat{\gamma }{}_{}{}^{r_j}\eta _{r_i}^{}=\eta _{r_i}\widehat{\gamma }{}_{}{}^{r_j}.`$ The operators of corresponding occupation numbers (for given $`\lambda ,\mu ,\alpha )`$ are $`\widehat{N}_r=\widehat{\gamma }^r\widehat{\gamma }_r`$. Since the diagonal operators $`(12\widehat{N}_r)`$ anticommute with the $`\widehat{\gamma }^r`$, then $`\eta _{r_i}=_{r=1}^{r_i1}(12\widehat{N}_r),`$ where
$$\eta _{r_i}^{11}\mathrm{\Phi }(n_1,\mathrm{},n_N;0;0;0)=\underset{r=1}{\overset{r_i1}{}}(1)^{n_r}\mathrm{\Phi }(n_1,\mathrm{},n_N;0;0;0),$$
(A.1.10)
etc. Here the occupation numbers $`n_r(m_r,q_r,t_r)`$ are introduced, which refer to the $`r`$-th states corresponding to operators $`\widehat{\gamma }_{(1,1,\alpha )}^r`$, etc., either empty ($`n_r,\mathrm{},t_r=0`$) or occupied ($`n_r,\mathrm{},t_r=1`$). To save writing we abbreviate the modified operators by the same symbols. For example, acting on free state $`0>_{r_i}`$ the creation operator $`\widehat{\gamma }_{r_i}`$ yields the one occupied state $`1>_{r_i}`$ with the phase $`+`$ or $``$ depending of parity of the number of quanta in the states $`r<r_i`$. Modified operators satisfy the same anticommutation relations of the basis operators (subsec.2.2). It is convenient to make use of notation $`\widehat{\gamma }{}_{r}{}^{(\lambda ,\mu ,\alpha )}e_r^{(\lambda ,\mu ,\alpha )}\widehat{b}{}_{(r\alpha )}{}^{\lambda \mu },`$ and abbreviate the pair of indices $`(r\alpha )`$ by the single symbol $`r`$. Then for each $`\mathrm{\Phi }{}_{}{}^{A}_{}^{(n)}`$ and any vector $`f^{(1)}`$ the operators $`\widehat{b}(f)`$ and $`\widehat{b}^{}(f)`$ imply
$$\begin{array}{c}\widehat{b}(f)\mathrm{\Phi }=\frac{1}{\sqrt{(n1)!}}\underset{\sigma S(n)}{}sgn(\sigma )\left(f\mathrm{\Phi }_{\sigma (1)}^{(1)}\right)\mathrm{\Phi }_{\sigma (2)}^{(1)}\mathrm{}\mathrm{\Phi }_{\sigma (n)}^{(1)},\hfill \\ \widehat{b}^{}(f)\mathrm{\Phi }=\frac{1}{\sqrt{(n+1)!}}\underset{\sigma S(n+1)}{}sgn(\sigma )\mathrm{\Phi }_{\sigma (0)}^{(1)}\mathrm{\Phi }_{\sigma (1)}^{(1)}\mathrm{}\mathrm{\Phi }_{\sigma (n)}^{(1)},\hfill \end{array}$$
(A.1.11)
where $`\mathrm{\Phi }_{(0)}^{(1)}f`$. One continues the $`\widehat{b}(f)`$ and $`\widehat{b}^{}(f)`$ by linearity to linear reflections, which are denoted by the same symbols acting respectively from $`{}_{}{}^{A}_{}^{(n)}`$ onto $`{}_{}{}^{A}_{}^{(n1)}`$ or $`{}_{}{}^{A}_{}^{(n+1)}`$. They are limited over the values $`\sqrt{n}|f|`$ and $`\sqrt{(n+1)}|f|`$ and can be expanded by continuation up to the reflections acting from $`{}_{}{}^{A}\overline{}_{}^{(n)}`$ onto $`{}_{}{}^{A}\overline{}_{}^{(n1)}`$ or $`{}_{}{}^{A}\overline{}_{}^{(n+1)}`$. Finally, they must be continued by linearity up to the linear operators acting from $`{}_{}{}^{A}`$ onto $`{}_{}{}^{A}`$ defined on the same closed region in $`{}_{}{}^{A}\overline{}_{}^{(n)}`$, namely in $`{}_{}{}^{A}`$, which is invariant with respect to reflections $`\widehat{b}(f)`$ and $`\widehat{b}^{}(f)`$. Hence, at $`f_i,g_i^{(1)}`$ $`(i=1,\mathrm{},n;j=1,\mathrm{},m)`$ all polynomials over $`\{\widehat{b}^{}(f_i)\}`$ and $`\{\widehat{b}(g_j)\}`$ are completely defined on $`{}_{}{}^{A}`$. While, for given $`\lambda ,\mu `$, one has
$$\begin{array}{c}<\lambda ,\mu \{\widehat{b}_r^{\lambda \mu }(f),\widehat{b}_{\lambda \mu }^r^{}(g)\}\lambda ,\mu >=\delta _r^r^{}.\hfill \end{array}$$
(A.1.12)
The mean values $`<\phi ;\widehat{b}_r^{\lambda \mu }(f)\widehat{b}_{\lambda \mu }^r(f)>`$ calculated at fixed $`\lambda ,\mu `$ for any element $`\mathrm{\Phi }{}_{}{}^{A}`$ equal to mean values of the symmetric operator of occupation number in terms of $`\widehat{N}^r=\widehat{b}_r(f)\widehat{b}^r(f)`$, with a wave function $`f`$ in the state described by $`\mathrm{\Phi }`$. Here, as usual, it is denoted $`<\phi ;A\mathrm{\Phi }>=TrP_\phi A=(\mathrm{\Phi },A\mathrm{\Phi })`$ for each vector $`\mathrm{\Phi }`$ with $`|\mathrm{\Phi }|=1`$, while the $`P_\phi `$ is projecting operator onto one dimensional space $`\{\lambda \mathrm{\Phi }|\lambda C\}`$ generated by $`\mathrm{\Phi }`$. Therewith, the probability of transition $`\phi \psi `$ is given $`Pr\{\phi |\psi \}=\left|(\psi ,\phi )\right|^2`$. The linear operator $`A`$ defined on the elements of linear manifold $`𝒟(A)`$ of $``$ takes the values in $``$. The $`𝒟(A)`$ is an overall closed region of definition of $`A`$, namely the closure of $`𝒟(A)`$ by the norm given in $``$ coincides with $``$. Meanwhile, the $`𝒟(A)`$ included in $`𝒟(A^{})`$ and $`A`$ coincides with the reduction of $`A^{}`$ on $`𝒟(A)`$, because $`𝒟(A)`$ is the symmetric operator such that the linear operator $`A^{}`$ is the maximal conjugated to $`A`$. That is, any operator $`A^{}`$ conjugated to $`A`$ \- $`(\mathrm{\Psi },A^{}\mathrm{\Phi })=(A^{}\mathrm{\Psi },\mathrm{\Phi })`$ for all $`\mathrm{\Phi }𝒟(A)`$ and $`\mathrm{\Psi }𝒟(A^{})`$ coincides with the reduction of $`A^{}`$ on some linear manifold $`𝒟(A^{})`$ included in $`𝒟(A^{})`$. Thus, the operator $`A^{}`$ is closed symmetric expansion of operator $`A`$, namely it is a closure of $`A`$. Self conjugated operator $`A`$, the closure of which is self conjugated as well, allows only the one self conjugated expansion $`A^{}`$. Hence, self conjugated closure $`\widehat{N}`$ of operator $`{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\widehat{b}^{}(f_i)\widehat{b}(f_i)`$, where $`\{f_i|i=1,\mathrm{},n\}`$ is an arbitrary orthogonal basis on $`^{(1)}`$, can be regarded as the operator of occupation number. For the vector $`\chi ^0{}_{}{}^{A}`$ and $`\chi ^{0(n)}=\delta _{0n}`$ one gets $`<\chi ^{0(n)},\widehat{N}(f)>=0`$ for all $`f^{(1)}`$. Thus, $`\chi ^0`$ is the vector of vacuum state: $`\widehat{b}(f)\chi ^0=0`$ for all $`f^{(1)}`$. If $`f=\{f_i|i=1,2,\mathrm{}\}`$ is an arbitrary orthogonal basis on $`^{(1)}`$, then due to irreducibility of operators $`\widehat{b}^{}(f_i)|f_if`$, the $`{}_{}{}^{A}`$ includes the $`0`$ and whole space $`{}_{}{}^{A}`$ as invariant subspaces with respect to all $`\widehat{b}^{}(f)`$. To define the 12 dimensional operator manifold $`\widehat{G}`$ we consider a set $`\widehat{}`$ of all the sequences $`\widehat{\mathrm{\Phi }}=\{\widehat{\mathrm{\Phi }}^{(0)},\widehat{\mathrm{\Phi }}^{(1)},\mathrm{},\widehat{\mathrm{\Phi }}^{(n)},\mathrm{}\}`$ with a finite number of nonzero elements provided by
$$\begin{array}{c}\widehat{\mathrm{\Phi }}_{(r_1,\mathrm{},r_n)}^{(n)}=\widehat{\mathrm{\Phi }}_{r_1}^{(1)}\mathrm{}\widehat{\mathrm{\Phi }}_{r_n}^{(1)}\widehat{G}^{(n)},\widehat{\mathrm{\Phi }}_{r_i}^{(1)}=\widehat{\zeta }_{r_i}\mathrm{\Phi }_{r_i}^{(1)}\widehat{G}_i^{(1)}=\widehat{𝒰}_i^{(1)}_i^{(1)},\hfill \\ \widehat{\zeta }_{r_i}\underset{\alpha _i=1}{\overset{3}{}}\widehat{\gamma }_{(\lambda _i,\mu _i,\alpha _i)}^{r_i}\zeta _{r_i}^{(\lambda _i,\mu _i,\alpha _i)}\widehat{𝒰}_{r_i}^{(1)},\widehat{G}^{(n)}=\widehat{𝒰}^{(n)}\overline{}^{(n)},\widehat{𝒰}_{(r_1,\mathrm{},r_n)}^{(n)}=\widehat{𝒰}_{r_1}^{(1)},\mathrm{}\widehat{𝒰}_{r_n}^{(1)}.\hfill \end{array}$$
(A.1.13)
Then, on the analogy of eq.(A.1.8) the operator manifold $`\widehat{G}`$ ensues
$$\widehat{G}=\underset{n=0}{\overset{\mathrm{}}{}}\widehat{G}^{(n)}=\underset{n=0}{\overset{\mathrm{}}{}}\left(\widehat{𝒰}^{(n)}\overline{}^{(n)}\right).$$
(A.1.14)
To define the secondary quantized form of one particle observable $`A`$ on $``$, following let consider a set of identical samples $`\widehat{}_i`$ of one particle space $`^{(1)}`$ and operators $`A_i`$ acting on them. To each closed linear operator $`A^{(1)}`$ in $`^{(1)}`$ with overall closed region of definition $`𝒟(A^{(1)})`$ following operators are corresponded:
$$\begin{array}{c}A_1^{(n)}=A^{(1)}I\mathrm{}I,\hfill \\ \mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\hfill \\ A_n^{(n)}=II\mathrm{}A^{(1)}.\hfill \end{array}$$
(A.1.15)
The sum $`{\displaystyle \underset{j=1}{\overset{n}{}}}A_j^{(n)}`$ is given on the intersection of regions of definition of operator terms including a linear manifold $`𝒟(A^{(1)})\mathrm{}𝒟(A^{(n)})`$ closed in $`\widehat{}^{(n)}`$. While, the $`A^{(n)}`$ is a minimal closed expansion of this sum with $`𝒟(A^{(n)})`$. One considers a linear manifold $`𝒟(\mathrm{\Omega }(A))`$ in $`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\widehat{}^{(n)}`$ defined as a set of all the vectors $`\mathrm{\Psi }`$ such as $`\mathrm{\Psi }^{(n)}𝒟(A^{(n)})`$ and $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left|A^{(n)}\mathrm{\Psi }^{(n)}\right|^2<\mathrm{}`$ . The manifold $`𝒟(\mathrm{\Omega }(A))`$ is closed in $``$. On this manifold one defines a closed linear operator $`\mathrm{\Omega }(A)`$ acting as $`\mathrm{\Omega }(A)^{(n)}=A^{(n)}\mathrm{\Psi }^{(n)}`$, namely $`\mathrm{\Omega }(A)\mathrm{\Phi }={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}A^{(n)}\mathrm{\Psi }^{(n)}`$, while the $`\mathrm{\Omega }(A)`$ is self conjugated operator with overall closed region of definition. We suppose that the vector $`\mathrm{\Phi }^{(n)}^{(n)}`$ is in the form eq.(A.1.3), where $`\mathrm{\Phi }_i𝒟(A)`$. Then, $`<\phi ^{(n)};A^{(n)}>={\displaystyle \underset{i=1}{\overset{n}{}}}<\phi _i;A>`$, which enables the expansion by continuing onto $`𝒟(A)`$. Thus, $`A^{(n)}`$ is the $`n`$ particle observable corresponding to one particle observable $`A`$. So $`<\phi ;\mathrm{\Omega }(A)>={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}<\phi ^{(n)};A^{(n)}>`$ for any $`\mathrm{\Phi }_i𝒟\left(\mathrm{\Omega }(A)\right)`$. While, the $`\mathrm{\Omega }(A)`$ reflects $`{}_{}{}^{A}𝒟=𝒟\left(\mathrm{\Omega }(A)\right){}_{}{}^{A}`$ into $`{}_{}{}^{A}`$. The reduction of $`\mathrm{\Omega }(A)`$ on $`{}_{}{}^{A}`$ is the self conjugated in the region $`{}_{}{}^{A}𝒟`$, because $`{}_{}{}^{A}`$ is the closed subspace of the $``$. Hence, the $`\mathrm{\Omega }(A)`$ is the secondary quantized form of one particle observable $`A`$ on the $``$.
The vacuum state reads
$$\begin{array}{ccc}\chi ^0(\nu _1,\nu _2,\nu _3,\nu _4)=1,1>^{\nu _1}1,2>^{\nu _2}2,1>^{\nu _3}2,2>^{\nu _4},\hfill & & \\ \nu _i=\{\begin{array}{cc}1\hfill & \text{if }\nu =\nu _i\text{for some }i\text{,}\hfill \\ 0\hfill & \text{otherwise},\hfill \end{array},\hfill & & \end{array}$$
(A.1.16)
where $`\chi _{}(1)>\chi ^0(1,0,0,0),\chi _+(1)>\chi ^0(0,0,0,1),\chi _{}(2)>\chi ^0(0,0,1,0),\chi _+(2)>\chi ^0(0,1,0,0),<\chi _\pm (\lambda )\chi _\pm (\mu )>=\delta _{\lambda \mu },`$ and $`<\chi _\pm (\lambda )\chi _{}(\mu )>=0,`$ provided by $`<\chi _\pm A\chi _\pm >{\displaystyle \underset{\lambda }{}}<\chi _\pm (\lambda )A\chi _\pm (\lambda )>`$ and the normalization condition
$$<\chi ^0(\nu _1^{},\nu _2^{},\nu _3^{},\nu _4^{})\chi ^0(\nu _1,\nu _2,\nu _3,\nu _4)>=\underset{i=1}{\overset{4}{}}\delta _{\nu _i\nu _i^{}}.$$
The state vectors are introduced
$$\begin{array}{c}\chi (\{n_r\}_1^N;\{m_r\}_1^M;\{q_r\}_1^Q;\{t_r\}_1^T;\{\nu _r\}_1^4)=(\widehat{b}_N^{11})^{n_N}\mathrm{}(\widehat{b}_1^{11})^{n_1}\hfill \\ (\widehat{b}_M^{12})^{m_M}\mathrm{}(\widehat{b}_1^{12})^{m_1}(\widehat{b}_Q^{21})^{q_Q}\mathrm{}(\widehat{b}_1^{21})^{q_1}(\widehat{b}_T^{22})^{t_T}\mathrm{}(\widehat{b}_1^{22})^{t_1}\chi ^0(\nu _1,\nu _2,\nu _3,\nu _4),\hfill \end{array}$$
(A.1.17)
where $`\{n_r\}_1^N=n_1,\mathrm{},n_N`$, etc., which are the eigenfunctions of modified operators. They form a whole set of orthogonal vectors
$$\begin{array}{c}<\chi ,\chi ^{}>=\underset{r=1}{\overset{N}{}}\delta _{n_rn_r^{}}\underset{r=1}{\overset{M}{}}\delta _{m_rm_r^{}}\underset{r=1}{\overset{Q}{}}\delta _{q_rq_r^{}}\underset{r=1}{\overset{T}{}}\delta _{t_rt_r^{}}\underset{r=1}{\overset{4}{}}\delta _{\nu _r\nu _r^{}}.\hfill \end{array}$$
(A.1.18)
Considering an arbitrary superposition
$$\chi =\underset{a=\{n_r\}_1^N,\{m_r\}_1^M,\{q_r\}_1^Q,\{t_r\}_1^T=0}{\overset{1}{}}c^{}(a)\chi (a),$$
(A.1.19)
the coefficients $`c^{}`$ of expansion are the corresponding amplitudes of probabilities. Taking into account eq.(A.1.12), the nonvanishing matrix elements of operators $`\widehat{b}_{r_k}^{11}`$ and $`\widehat{b}_{11}^{r_k}`$ read
$$\begin{array}{c}<\chi (\{n_r^{}\}_1^N;0;0;0;1,0,0,0)|\widehat{b}_{r_k}^{11}\chi (\{n_r\}_1^N;0;0;0;1,0,0,0)>=\hfill \\ \\ =<1,1\widehat{b}_{11}^{r_1^{}}\mathrm{}\widehat{b}_{11}^{r_n^{}}\widehat{b}_{r_k}^{11}\widehat{b}_{r_n}^{11}\mathrm{}\widehat{b}_{r_1}^{11}1,1>=\hfill \\ \\ =\{\begin{array}{cc}(1)^{n^{}k^{}}\hfill & \text{if }n_r=n_r^{}\text{ for }rr_k\text{ and }n_{r_k}=0;n_{r_k}^{}=1,\hfill \\ 0\hfill & \text{otherwise},\hfill \end{array}\hfill \\ <\chi (\{n_r^{}\}_1^N;0;0;0;1,0,0,0)|\widehat{b}_{11}^{r_k}\chi (\{n_r\}_1^N;0;0;0;1,0,0,0)>=\hfill \\ \\ =<1,1\widehat{b}_{11}^{r_1^{}}\mathrm{}\widehat{b}_{11}^{r_n^{}}\widehat{b}_{11}^{r_k}\widehat{b}_{r_n}^{11}\mathrm{}\widehat{b}_{r_1}^{11}1,1>=\hfill \\ \\ =\{\begin{array}{cc}(1)^{nk}\hfill & \text{if }n_r=n_r^{}\text{ for }rr_k\text{ and }n_{r_k}^{}=0;n_{r_k}=1,\hfill \\ 0\hfill & \text{otherwise},\hfill \end{array}\hfill \end{array}$$
(A.1.20)
where one denotes $`n={\displaystyle \underset{r=1}{\overset{N}{}}}n_r,n^{}={\displaystyle \underset{r=1}{\overset{N}{}}}n_r^{},`$ the $`r_k`$ and $`r_k^{}`$ are $`k`$-th and $`k^{}`$-th terms of regulated sets of $`\{r_1,\mathrm{},r_n\}(r_1<r_2<\mathrm{}<r_n)`$ and $`\{r_1^{},\mathrm{},r_n^{}\}(r_1^{}<r_2^{}<\mathrm{}<r_n^{})`$, respectively. Continuing along this line we get a whole set of explicit forms of matrix elements of the rest of operators $`\widehat{b}_{r_k}`$ and $`\widehat{b}^{r_k}`$. Hence
$$\underset{\{\nu _r\}=0}{\overset{1}{}}<\chi ^0\widehat{\mathrm{\Phi }}(\zeta )\chi >=\underset{r=1}{\overset{N}{}}c_{n_r}^{}e_{(1,1,\alpha )}^{n_r}\mathrm{\Phi }_{n_r}^{(1,1,\alpha )}+\mathrm{},$$
(A.1.21)
provided by
$$c_{n_r}^{}\delta _{1n_r}c^{}(0,\mathrm{},n_r,\mathrm{},0;0;0;0),\mathrm{}$$
(A.1.22)
Hereinafter we change a notation of the coefficients
$$\begin{array}{c}\overline{c}(r^{11})=c_{n_r}^{},\overline{c}(r^{21})=c_{q_r}^{},N_{11}=N,N_{21}=Q,\hfill \\ \overline{c}(r^{12})=c_{m_r}^{},\overline{c}(r^{22})=c_{t_r}^{},N_{12}=M,N_{22}=T,\hfill \end{array}$$
(A.1.23)
and make use of convention
$$F_{r^{\lambda \mu }}=\underset{\alpha }{}e_{(\lambda ,\mu ,\alpha )}^{r^{\lambda \mu }}\mathrm{\Phi }_{r^{\lambda \mu }}^{(\lambda ,\mu ,\alpha )},\underset{\{\nu _r\}=0}{\overset{1}{}}<\chi ^0\widehat{A}\chi ><\chi ^0\widehat{A}\chi >,$$
(A.1.24)
The matrix elements of operator vector and covector fields take the final forms
$$\begin{array}{c}<\chi ^0\widehat{\mathrm{\Phi }}(\zeta )\chi >=\underset{\lambda \mu =1}{\overset{2}{}}\underset{r^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}\overline{c}(r^{\lambda \mu })F_{r^{\lambda \mu }}(\zeta ),\hfill \\ <\chi \overline{\widehat{\mathrm{\Phi }}}(\zeta )\chi ^0>=\underset{\lambda \mu =1}{\overset{2}{}}\underset{r^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}\overline{c}^{}(r^{\lambda \mu })F^{r^{\lambda \mu }}(\zeta ).\hfill \end{array}$$
(A.1.25)
In the following we shall use a convention: $`\left\{{\displaystyle \underset{\lambda \mu }{\overset{2}{}}}\right\}_1^n{\displaystyle \underset{\lambda _1\mu _1}{\overset{2}{}}}\mathrm{}{\displaystyle \underset{\lambda _n\mu _n}{\overset{2}{}}},`$ $`r_i^{\lambda \mu }r^{\lambda _i\mu _i}`$ and $`\overline{c}(r_1^{11},\mathrm{},r_n^{11})=c^{}(n_1,\mathrm{},n_N;0;0;0),`$ etc. The anticommutation relations ensue
$$<\chi _{}\{\underset{i}{\widehat{b}}{}_{r}{}^{+},\underset{i}{\widehat{b}}{}_{+}{}^{r^{}}\}\chi _{}>=<\chi _+\{\underset{i}{\widehat{b}}{}_{r}{}^{},\underset{i}{\widehat{b}}{}_{}{}^{r^{}}\}\chi _+>=\delta _r^r^{},$$
(A.1.26)
provided by $`\underset{i}{\widehat{\gamma }}{}_{r}{}^{(\lambda \alpha )}=\underset{i}{\widehat{e}}{}_{r}{}^{(\lambda \alpha )}\underset{i}{\widehat{b}}{}_{(r\alpha )}{}^{\lambda },(r\alpha )r.`$ The state functions
$$\chi =(\underset{\eta }{\widehat{b}}_N^+)^{n_N}\mathrm{}(\underset{\eta }{\widehat{b}}_1^+)^{n_1}(\underset{\eta }{\widehat{b}}_M^{})^{m_M}\mathrm{}(\underset{\eta }{\widehat{b}}_1^{})^{m_1}(\underset{u}{\widehat{b}}_Q^+)^{q_Q}\mathrm{}(\underset{u}{\widehat{b}}_1^+)^{q_1}(\underset{u}{\widehat{b}}_T^{})^{t_T}\mathrm{}(\underset{u}{\widehat{b}}_1^{})^{t_1}\chi _{}(\lambda )\chi _+(\mu ),$$
(A.1.27)
form a whole set of orthogonal eigenfunctions of corresponding operators of occupation numbers $`\underset{i}{\widehat{N}}{}_{r}{}^{\lambda }=\underset{i}{\widehat{b}}{}_{r}{}^{\lambda }\underset{i}{\widehat{b}}_\lambda ^r`$ with the expectation values 0,1.
$``$Differential Geometric Aspect of the OM
An explicit form of matrix element of operator tensor reads
$$\begin{array}{c}\frac{1}{\sqrt{n!}}<\chi ^0\widehat{\mathrm{\Phi }}(\zeta _1)\mathrm{}\widehat{\mathrm{\Phi }}(\zeta _n)\chi >=\hfill \\ =\left\{\underset{\lambda \mu =1}{\overset{2}{}}\right\}_1^n\underset{r_1^{\lambda \mu },\mathrm{},r_n^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}\overline{c}(r_1^{\lambda \mu },\mathrm{},r_n^{\lambda \mu })F_{r_1^{\lambda \mu }}(\zeta _1)\mathrm{}F_{r_n^{\lambda \mu }}(\zeta _n),\hfill \end{array}$$
(A.1.28)
where $``$ stands for exterior product.
The linear operator form of $`1`$ degree $`\widehat{\omega }^1`$ is a linear operator valued function on $`\widehat{𝐓}_{\mathrm{\Phi }_p}`$, namely $`\widehat{\omega }^1(\widehat{𝐀}_p):\widehat{𝐓}_{\mathrm{\Phi }_p}\widehat{R}`$, where $`\widehat{𝐀}_p\widehat{𝐓}_{\mathrm{\Phi }_p}`$, and the operator $`\widehat{\omega }^1(\widehat{𝐀})=<\widehat{\omega }^1,𝐀>\widehat{R}`$ corresponds to $`\widehat{𝐀}_p`$ at the point $`𝚽_p`$, provided, according to eq.(A.1.25), with
$$<\chi \widehat{\omega }^1\chi ^0>=\underset{\lambda ,\mu =1}{\overset{2}{}}\underset{r^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}\widehat{c}^{}(r^{\lambda \mu })\omega _{r^{\lambda \mu }}^1,$$
(A.1.29)
where $`\omega _{r^{\lambda \mu }}^1=e_{r^{\lambda \mu }}^{(\lambda ,\mu ,\alpha )}\omega _{(\lambda ,\mu ,\alpha )}^{r^{\lambda \mu }}`$, the $`<\omega _{r^{\lambda \mu }}^1,𝐀>=\omega _{r^{\lambda \mu }}^1(𝐀)`$ is a linear form on $`𝐓_p`$, and
$$\begin{array}{c}\widehat{\omega }^1(\lambda _1\widehat{𝐀}_1+\lambda _2\widehat{𝐀}_2)=\lambda _1\widehat{\omega }^1(\widehat{𝐀}_1)+\lambda _2\widehat{\omega }^1(\widehat{𝐀}_2),\hfill \\ \lambda _1,\lambda _2R,\widehat{𝐀}_1,\widehat{𝐀}_2\widehat{𝐓}_{\mathrm{\Phi }_p}.\hfill \end{array}$$
(A.1.30)
The set of all linear operator forms defined at the point $`𝚽_p`$ fill up the operator vector space $`\widehat{𝐓}_{\mathrm{\Phi }_p}^{}`$ dual to $`\widehat{𝐓}_{\mathrm{\Phi }_p}`$. While, the $`\{\widehat{\gamma }_r\}`$ serves as a basis for them. The operator $`n`$ form is defined as the exterior product of operator 1 forms
$$\widehat{\omega }^n(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n)=\left(\widehat{\omega }_1^1\mathrm{}\widehat{\omega }_n^1\right)(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n)=\begin{array}{ccc}\widehat{\omega }_1^1(\widehat{𝐀}_1)\mathrm{}\mathrm{}\hfill & \widehat{\omega }_n^1(\widehat{𝐀}_1\hfill & \\ \mathrm{}\hfill & \mathrm{}\hfill & \\ \widehat{\omega }_1^1(\widehat{𝐀}_n)\mathrm{}\mathrm{}\hfill & \widehat{\omega }_n^1(\widehat{𝐀}_n)\hfill & \end{array}.$$
(A.1.31)
Here as well as for the rest of this section we abbreviate the set of indices $`(\lambda _i,\mu _i,\alpha _i)`$ by the single symbol $`i`$. If $`\{\widehat{\gamma }{}_{i}{}^{r_i}\}`$ and $`\{\widehat{\gamma }{}_{r_i}{}^{i}\}`$ are dual basis respectively in $`\widehat{𝐓}_{\mathrm{\Phi }_p}`$ and $`\widehat{𝐓}_{\mathrm{\Phi }_p}^{}`$, then the $`\{\widehat{\gamma }{}_{1}{}^{r_1}\mathrm{}\widehat{\gamma }{}_{p}{}^{r_p}\widehat{\gamma }{}_{s_1}{}^{1}\mathrm{}\widehat{\gamma }{}_{s_q}{}^{q}\}`$ will be the basis in operator space $`\widehat{𝐓}_q^p=\underset{p}{\underset{}{\widehat{𝐓}_{\mathrm{\Phi }_p}\mathrm{}\widehat{𝐓}_{\mathrm{\Phi }_p}}}\underset{q}{\underset{}{\widehat{𝐓}_{\mathrm{\Phi }_p}^{}\mathrm{}\widehat{𝐓}_{\mathrm{\Phi }_p}^{}}}.`$ Any operator tensor $`\widehat{𝐓}\widehat{𝐓}_q^p(𝚽_p)`$ can be written
$$\widehat{𝐓}=T_{j_1\mathrm{}j_q}^{i_1\mathrm{}i_p}(r_1,\mathrm{},r_p,s_1,\mathrm{},s_q)\widehat{\gamma }{}_{i_1}{}^{r_1}\mathrm{}\widehat{\gamma }{}_{i_p}{}^{r_p}\widehat{\gamma }{}_{s_1}{}^{j_1}\mathrm{}\widehat{\gamma }{}_{s_q}{}^{j_q},$$
where $`T_{j_1\mathrm{}j_q}^{i_1\mathrm{}i_p}(r_1,\mathrm{},r_p,s_1,\mathrm{},s_q)=T(\widehat{\gamma }{}_{r_1}{}^{i_1}\mathrm{}\widehat{\gamma }{}_{r_p}{}^{i_p}\widehat{\gamma }{}_{j_1}{}^{s_1}\mathrm{}\widehat{\gamma }{}_{j_q}{}^{s_q})`$ are the components of $`\widehat{𝐓}`$ in $`\{\widehat{\gamma }{}_{i}{}^{r_i}\}`$ and $`\{\widehat{\gamma }{}_{r_i}{}^{i}\}`$. Any antisymmetric operator tensor of $`\widehat{(0,n)}`$ type reads
$$\widehat{𝐓}^{}=T_{i_1\mathrm{}i_n}\gamma ^{\widehat{ı}_1}\mathrm{}\gamma ^{\widehat{ı}_n}=\underset{i_1<\mathrm{}<i_n}{}T_{i_1\mathrm{}i_n}d\mathrm{\Phi }^{\widehat{ı_1}}\mathrm{}d\mathrm{\Phi }^{\widehat{ı_n}}.$$
(A.1.32)
Let the $`\widehat{𝒟}_1`$ and $`\widehat{𝒟}_2`$ are two compact convex parallelepipeds in oriented $`n`$ dimensional operator space $`\widehat{𝐑}^n`$ and the $`f:\widehat{𝒟}_1\widehat{𝒟}_2`$ is differentiable reflection of interior of $`\widehat{𝒟}_1`$ into $`\widehat{𝒟}_2`$ retaining an orientation, namely for any function $`\phi C^{\mathrm{}}`$ defined on $`\widehat{𝒟}_2`$ it holds $`\phi fC^{\mathrm{}}`$ and $`f^{}\phi \left(𝚽_p\right)=\phi \left(f\left(𝚽_p\right)\right)`$, where $`f^{}`$ is an image of function $`\phi \left(f\left(𝚽_p\right)\right)`$ on $`\widehat{𝒟}_1`$ at the point $`𝚽_p`$. Hence, the function $`f`$ induces a linear reflection $`\widehat{d}f:\widehat{𝐓}\left(\widehat{𝒟}_1\right)\widehat{𝐓}\left(\widehat{𝒟}_2\right)`$ as an operator differential of $`f`$ implying $`\widehat{d}f\left(\widehat{𝐀}_p\right)\phi =\widehat{𝐀}_p(\phi f)`$ for any operator vector $`\widehat{𝐀}_p\widehat{𝐓}_{\mathrm{\Phi }_p}`$ and for any function $`\phi C^{\mathrm{}}`$ defined in the neighbourhood of $`𝚽_{}^{}{}_{p}{}^{}=f\left(𝚽_p\right)`$. If the function $`f`$ is given in the form $`\mathrm{\Phi }_{}^{}{}_{}{}^{i}=\mathrm{\Phi }_{}^{}{}_{}{}^{i}\left(\mathrm{\Phi }_p\right)`$ and $`\widehat{𝐀}_p=\left(A^i\widehat{}/\mathrm{\Phi }^i\right)_p`$, then in terms of local coordinates one gets $`\left(\widehat{d}f\right)\widehat{𝐀}_p=A^i\left({\displaystyle \frac{\mathrm{\Phi }_{}^{}{}_{}{}^{j}}{\mathrm{\Phi }^i}}\right)_p\left({\displaystyle \frac{\widehat{}}{\mathrm{\Phi }_{}^{}{}_{}{}^{j}}}\right)_p^{}.`$ So, if $`f_1:\widehat{𝒟}_1\widehat{𝒟}_2`$ and $`f_2:\widehat{𝒟}_2\widehat{𝒟}_3`$ then $`\widehat{d}\left(f_2f_1\right)=\widehat{d}f_2\widehat{d}f_1`$. The differentiable reflection $`f:\widehat{𝒟}_1\widehat{𝒟}_2`$ induces the reflection $`\widehat{f}^{}:\widehat{𝐓}^{}\left(\widehat{𝒟}_2\right)\widehat{𝐓}^{}\left(\widehat{𝒟}_1\right)`$ conjugated to $`\widehat{f}_{}`$. The latter is the operator differential of $`f`$, while
$$<\widehat{f}^{}\widehat{\omega ^{}}^1,\widehat{𝐀}>_{\mathrm{\Phi }_p}=<\widehat{\omega ^{}}^1,\widehat{f}_{}\widehat{𝐀}>|_{f\left(\mathrm{\Phi }_p\right)},$$
(A.1.33)
where $`\widehat{𝐀}|_{f\left(\mathrm{\Phi }_p\right)}=\left(\widehat{d}f\right)\widehat{𝐀}_p`$ and $`\widehat{\omega ^{}}^1\widehat{𝐓}^{}|_{f\left(\mathrm{\Phi }_p\right)}`$. Hence $`\widehat{f}^{}\left(\widehat{d}\phi \right)=\widehat{d}\left(\widehat{f}^{}\phi \right)`$ and
$$\begin{array}{c}\widehat{f}^{}T(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n)|_{\mathrm{\Phi }_p}=T(\widehat{f}_{}\widehat{𝐀}_1,\mathrm{},\widehat{f}_{}\widehat{𝐀}_n)|_{f\left(\mathrm{\Phi }_p\right)},\hfill \\ T(\widehat{f}^{}\widehat{\omega }_1^1,\mathrm{},\widehat{f}^{}\widehat{\omega }_n^1)|_{\mathrm{\Phi }_p}=\widehat{f}_{}T(\widehat{\omega }_1^1,\mathrm{},\widehat{f}^{}\widehat{\omega }_n^1)|_{f\left(\mathrm{\Phi }_p\right)}.\hfill \end{array}$$
(A.1.34)
For any differential operator $`n`$ form $`\widehat{\omega }^n`$ on $`\widehat{𝒟}_2`$ the reflection $`f`$ induces the operator $`n`$ form $`\widehat{f}^{}\widehat{\omega }^n`$ on $`\widehat{𝒟}_1`$
$$\left(\widehat{f}^{}\widehat{\omega }^n\right)(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n)|_{\mathrm{\Phi }_p}=\widehat{f}_{}\widehat{\omega }^n(\widehat{f}_{}\widehat{𝐀}_1,\mathrm{},\widehat{f}_{}\widehat{𝐀}_n)|_{f\left(\mathrm{\Phi }_p\right)}.$$
(A.1.35)
If $`\widehat{\omega ^{}}^1=\alpha _i^{}d\mathrm{\Phi }_{}^{}{}_{}{}^{\widehat{ı}}`$ then $`\widehat{f}^{}\left(\alpha _i^{}d\mathrm{\Phi }_{}^{}{}_{}{}^{\widehat{ı}}\right)=\alpha _i^{}{\displaystyle \frac{\mathrm{\Phi }_{}^{}{}_{}{}^{i}}{\mathrm{\Phi }^j}}d\mathrm{\Phi }^{\widehat{ȷ}}.`$ This can be extended up to $`\widehat{\omega ^{}}^n\widehat{\omega }^n`$
$$\widehat{f}^{}\left(\underset{i_1<\mathrm{}<i_n}{}T_{}^{}{}_{i_1\mathrm{}i_n}{}^{}d\mathrm{\Phi }_{}^{}{}_{}{}^{\widehat{ı_1}}\mathrm{}d\mathrm{\Phi }_{}^{}{}_{}{}^{\widehat{ı_n}}\right)=\underset{\begin{array}{c}i_1<\mathrm{}<i_n\hfill \\ j_1<\mathrm{}<j_n\hfill \end{array}}{}T_{}^{}{}_{i_1\mathrm{}i_n}{}^{}\frac{\mathrm{\Phi }_{}^{}{}_{}{}^{i_1}}{\mathrm{\Phi }^{j^1}}\mathrm{}\frac{\mathrm{\Phi }_{}^{}{}_{}{}^{i_n}}{\mathrm{\Phi }^{j^n}}d\mathrm{\Phi }^{\widehat{ı_1}}\mathrm{}d\mathrm{\Phi }^{\widehat{ı_n}},$$
(A.1.36)
namely $`\widehat{f}^{}\widehat{\omega ^{}}^n=J_\mathrm{\Phi }\widehat{\omega }^n=\left(detdf\right)\widehat{\omega }^n,`$ where $`J_\mathrm{\Phi }`$ is the Jacobian of reflection $`J_\mathrm{\Phi }={\displaystyle \frac{\mathrm{\Phi }_{}^{}{}_{}{}^{i}}{\mathrm{\Phi }^j}}`$. While
$$\left(\widehat{f}_1\widehat{f}_2\right)^{}=\widehat{f}_1^{}\widehat{f}_2^{},\widehat{f}^{}\left(\widehat{\omega }_1\widehat{\omega }_2\right)=\widehat{f}^{}\left(\widehat{\omega }_1\right)\widehat{f}^{}\left(\widehat{\omega }_2\right).$$
We may consider the integration of operator $`n`$ form implying $`{\displaystyle _{\widehat{𝒟}_1}}\widehat{f}^{}\widehat{\omega }^n={\displaystyle _{\widehat{𝒟}_2}}\widehat{\omega }^n.`$ In general, let $`\widehat{𝒟}_1`$ is the limited convex $`n`$ dimensional parallelepiped in the $`n`$ dimensional operator space $`\widehat{𝐑}^n`$. One defines the $`n`$ dimensional $`i`$-th piece of integration path $`\widehat{\sigma }^i`$ in $`\widehat{G}`$ as $`\widehat{\sigma }^i=(\widehat{𝒟}_i,f_i,Or_i)`$, where $`\widehat{𝒟}_i\widehat{𝐑}^n,f_i:\widehat{𝒟}_i\widehat{G}`$ and the $`Or_i`$ is an orientation of $`\widehat{𝐑}^n`$. Then, the integral over the operator $`n`$ form $`\widehat{\omega }^n`$ along the operator $`n`$ dimensional chain $`\widehat{c}_n={\displaystyle m_i\widehat{\sigma }^i}`$ may be written
$$_{\widehat{c}_n}\widehat{\omega }^n=m_i_{\widehat{\sigma }^i}\widehat{\omega }^n=m_i_{\widehat{𝒟}_i}\widehat{f}^{}\widehat{\omega }^n,$$
where the $`m_i`$ is a multiple number. Taking into account the eq.(A.1.29), the matrix element yields
$$<\chi _{\widehat{c}_n}\widehat{\omega }^n\chi ^0>\left\{\underset{\lambda \mu =1}{\overset{2}{}}\right\}_1^n\underset{r_1^{\lambda \mu },\mathrm{},r_n^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}m_i\overline{c}(r_1^{\lambda \mu },\mathrm{},r_n^{\lambda \mu })_{\widehat{𝒟}_i}\widehat{f}^{}\widehat{\omega }^n(r_1^{\lambda \mu },\mathrm{},r_n^{\lambda \mu }).$$
(A.1.37)
Next, we may apply the analog of exterior differentiation. We define the operator $`(n+1)`$ form $`\widehat{d}\widehat{\omega }^n`$ on $`(n+1)`$ operator vectors $`\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_{n+1}\widehat{𝐓}_{\mathrm{\Phi }_p}`$ by considering diffeomorphic reflection $`f`$ of the neighbourhood of the point $`0`$ in $`\widehat{𝐑}^n`$ into neighbourhood of the point $`𝚽_p`$ in $`\widehat{G}`$. The prototypes of operator vectors $`\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_{n+1}\widehat{𝐓}_{\mathrm{\Phi }_p}\left(\widehat{G}\right)`$ at the operator differential of $`f`$ belong to tangent operator space $`\widehat{𝐑}^n`$ in $`0`$. Namely, the prototypes are the operator vectors $`\widehat{\xi }_1,\mathrm{},\widehat{\xi }_{n+1}\widehat{𝐑}^n`$. Let $`f`$ reflects the parallelepiped $`\widehat{𝚷}^{},`$ stretched over the $`\widehat{\xi }_1,\mathrm{},\widehat{\xi }_{n+1}`$, onto the $`(n+1)`$ dimensional piece $`\widehat{𝚷}`$ on the $`\widehat{G}`$. While the border of the $`n`$ dimensional chain $`\widehat{𝚷}`$ in $`\widehat{𝐑}^{n+1}`$ defined as follows: the pieces $`\widehat{\sigma }^i`$ of the chain $`\widehat{𝚷}`$ are $`n`$ dimensional facets $`\widehat{𝚷}_i`$ of parallelepiped $`\widehat{𝚷}`$ with the reflections embedding the facets into $`\widehat{𝐑}^{n+1}`$: $`f_i:\widehat{𝚷}_i\widehat{𝐑}^{n+1}`$, and the orientations $`Or_i`$ has defined as $`\widehat{𝚷}={\displaystyle \widehat{\sigma }^i},\widehat{\sigma }^i=(\widehat{𝚷}_i,f_i,Or_i)`$ Considering the curvilinear parallelepiped
$$F(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n)=_{\widehat{𝚷}}\widehat{\omega }^n,$$
one may state that the unique operator of the $`(n+1)`$-form $`\widehat{\mathrm{\Omega }}`$ exists on $`\widehat{𝐓}_{\mathrm{\Phi }_p}`$, which is the principle $`(n+1)`$ linear part in $`0`$ of integral over the border of $`F(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_n)`$, namely
$$F(\epsilon \widehat{𝐀}_1,\mathrm{},\epsilon \widehat{𝐀}_n)=\epsilon ^{n+1}\widehat{\mathrm{\Omega }}(\widehat{𝐀}_1,\mathrm{},\widehat{𝐀}_{n+1})+O\left(\epsilon ^{n+1}\right),$$
(A.1.38)
where $`\widehat{\mathrm{\Omega }}`$ is independent of choice of the coordinates used in definition of $`F`$. The prove of it is the same to those of similar one given in the differential geometry . If in local coordinates $`\widehat{\omega }^n={\displaystyle \underset{i_1<\mathrm{}<i_n}{}}T_{i_1\mathrm{}i_n}d\mathrm{\Phi }^{\widehat{ı_1}}\mathrm{}d\mathrm{\Phi }^{\widehat{ı_n}},`$ then
$$\widehat{\mathrm{\Omega }}=\widehat{d}\widehat{\omega }^n=\underset{i_1<\mathrm{}<i_n}{}\widehat{d}T_{i_1\mathrm{}i_n}d\mathrm{\Phi }^{\widehat{ı_1}}\mathrm{}d\mathrm{\Phi }^{\widehat{ı_n}}.$$
(A.1.39)
The operator of exterior differential $`\widehat{d}`$ commutes with the reflection $`f:\widehat{G}\widehat{G}`$
$$\widehat{d}\left(\widehat{f}^{}\widehat{\omega }^n\right)=\widehat{f}^{}\left(\widehat{d}\widehat{\omega }^n\right).$$
Define the exterior differential by operator (n+1) form
$$\widehat{d}\widehat{\omega }^n=\underset{\begin{array}{c}i_0\hfill \\ i_1<\mathrm{}<i_n\hfill \end{array}}{}\frac{T_{i_1\mathrm{}i_n}}{\mathrm{\Phi }^{i_0}}d\mathrm{\Phi }^{\widehat{i_0}}d\mathrm{\Phi }^{\widehat{i_1}}\mathrm{}d\mathrm{\Phi }^{\widehat{i_n}}=\underset{i_1<\mathrm{}<i_n}{}(\widehat{d}T_{i_1\mathrm{}i_n})d\mathrm{\Phi }^{\widehat{i_1}}\mathrm{}d\mathrm{\Phi }^{\widehat{i_n}},$$
(A.1.40)
one gets
$$\begin{array}{c}<\chi \widehat{d}\widehat{\omega }^n\chi ^0>\hfill \\ \underset{i_1<\mathrm{}<i_n}{}\left\{\underset{\lambda ,\mu =1}{\overset{2}{}}\right\}_1^n\underset{r_1^{\lambda \mu },\mathrm{},r_n^{\lambda \mu }=1}{\overset{N_{\lambda \mu }}{}}\overline{c}(r_1^{\lambda \mu },\mathrm{}r_n^{\lambda \mu })(dT(r_1^{\lambda \mu },\mathrm{},r_n^{\lambda \mu })_{i_1\mathrm{}i_n})d\mathrm{\Phi }_{r_1^{\lambda \mu }}^{i_1}\mathrm{}d\mathrm{\Phi }_{r_n^{\lambda \mu }}^{i_n}.\hfill \end{array}$$
(A.1.41)
Reflected upon the results far obtained within this section we may draw a conclusion that the matrix element of any geometric object of operator manifold $`\widehat{G}`$ yields corresponding geometric object of wave manifold $`𝒢`$.
## Appendix B
## 1 Reflection of the Fermi Fields
The rotation angles (subsec.4.2) are determined from the constraint imposed upon distortion transformations that a sum of distorted parts of corresponding basis vectors $`O_\lambda `$ and $`\sigma _\beta `$ should be zero for given $`\lambda `$
$$<O_{(\lambda \alpha )},O_\tau >_{\tau \lambda }+\frac{1}{2}\epsilon _{\alpha \beta \gamma }\frac{<\sigma _{(\lambda \beta )},\sigma _\gamma >}{<\sigma _{(\lambda \beta )},\sigma _\beta >}=0,$$
(B.1.1)
where $`\epsilon _{\alpha \beta \gamma }`$ is an antisymmetric unit tensor. Thereupon, $`\mathrm{tan}\theta _{(\lambda \alpha )}=\kappa a_{(\lambda \alpha )}`$, where $`\theta _{(\lambda \alpha )}`$ is the particular rotation angle around the axis $`\sigma _\alpha `$. Since the $`R`$ should be independent of the sequence of rotation axes, then it implies the mean value $`R={\displaystyle \frac{1}{6}}{\displaystyle \underset{ijk}{}}R^{(ijk)}`$, where $`R^{(ijk)}`$ the matrix of rotations occurred in the given sequence $`(ijk)`$ $`(i,j,k=1,2,3)`$. The field $`a_f`$ is due to the distortion of basis pseudovector $`O_\lambda `$, while the distortion of $`\sigma _\alpha `$ follows from eq.(B.1.1). We consider the reflection of the Fermi fields and their dynamics from the flat manifold $`\underset{u}{𝐺}`$ to distorted manifold $`\stackrel{~}{\underset{u}{𝐺}}`$, and vice versa. While we construct a diffeomorphism $`u(u_f):\underset{u}{𝐺}\stackrel{~}{\underset{u}{𝐺}}`$, where the holonomic functions $`u(u_f)`$ satisfy defining relation
$$e\psi =e^f+\chi ^f(𝐁_f).$$
(B.1.2)
Here $`e^f`$ and $`e`$ are the basis vectors on $`\underset{u}{𝐺}`$ and $`\stackrel{~}{\underset{u}{𝐺}}`$. The $`\psi `$ is taken to denote $`\psi {\displaystyle \frac{u}{u_f}}.`$ The covector
$$\chi _{(\tau \beta )}^f(𝐁_f)=e_{(\lambda \alpha )}\chi _{(\tau \beta )}^{(\lambda \alpha )}=\frac{1}{2}e_{(\lambda \alpha )}_0^{u_f}(\underset{u}{}{}_{(\rho \gamma )}{}^{f}D_{(\tau \beta )}^{(\lambda \alpha )}\underset{u}{}{}_{(\tau \beta )}{}^{f}D_{(\rho \gamma )}^{(\lambda \alpha )})𝑑u_f^{(\rho \gamma )}$$
(B.1.3)
realizes the coordinates $`u`$ by providing a criteria of integration and undegeneration <sup>2</sup><sup>2</sup>2I wish to thank S.P.Novikov for valuable discussion of this question. The Lagrangian $`L(x)`$ of fields $`\mathrm{\Psi }(u)`$ may be obtained under the reflection from the Lagrangian $`L_f(u_f)`$ of corresponding shadow fields $`\mathrm{\Psi }_f(u_f)`$ and vice versa. The $`\mathrm{\Psi }_f(u_f)`$ is defined as the section of vector bundle associated with the primary gauge group G by reflection $`\mathrm{\Psi }_f:\underset{u}{𝐺}E`$ such that $`p\mathrm{\Psi }_f(u_f)=u_f`$, where $`u_f\underset{u}{𝐺}`$ is the point of flat manifold $`\underset{u}{𝐺}`$ (specified by index $`({}_{f}{}^{})`$). The $`\mathrm{\Psi }_f`$ takes value in standard fiber $`F_{u_f}`$ upon $`u_f:p^1(U^{(f)})=U^{(f)}\times F_{u_f}`$, where $`U^{(f)}`$ is the region of base of principle bundle upon which an expansion into direct product $`p^1(U^{(f)})=U^{(f)}\times G`$ is defined. The fiber is Hilbert vector space on which a linear representation $`U_f`$ of the group G is given. Respectively $`\mathrm{\Psi }(u)F_u`$, where $`F_u`$ is the fiber upon $`u:p^1(U)=U\times F_u`$, $`U`$ is the region of base $`\stackrel{~}{\underset{u}{𝐺}}`$. Thus, the reflection of bispinor fields may be written down
$$\begin{array}{c}\mathrm{\Psi }(u)=R(𝐁_f)\mathrm{\Psi }_f(u_f),\overline{\mathrm{\Psi }}(u)=\overline{\mathrm{\Psi }}_f(u_f)\stackrel{~}{R}^+(𝐁_f),g(u)\mathrm{\Psi }(u)=\hfill \\ S(B_f)R(𝐁_f)\gamma _fD\mathrm{\Psi }_f(u_f),\left(\overline{\mathrm{\Psi }}(u)\right)g(u)=S(B_f)\left(D\overline{\mathrm{\Psi }}_f(u_f)\right)\gamma _f\stackrel{~}{R}^+(𝐁_f).\hfill \end{array}$$
(B.1.4)
Reviewing the notation $`𝐁(u_f)=T^a𝐁^a(u_f)`$ is the gauge field of distortion with the values in Lie algebra of group G, $`R(𝐁_f)`$ is the reflection matrix (see eq.(B.1.7)), $`\stackrel{~}{R}=\gamma ^0R\gamma ^0,`$ $`D=^fig𝐁,g^{(\lambda \alpha )}(\theta )=V_{(i,l)}^{(\lambda \alpha )}(\theta )\gamma _f^{(i,l)}`$, $`V_{(i,l)}^{(\lambda \alpha )}(\theta )`$ are congruence parameters of curves (Latin indices refer to tetrad components). The matrices $`\gamma _f^{(\pm \alpha )}={\displaystyle \frac{1}{\sqrt{2}}}(\gamma ^0\sigma ^\alpha \pm \gamma ^\alpha )`$, $`\gamma ^0,\gamma ^\alpha `$ are Dirac matrices. $``$ is covariant derivative defined on $`\stackrel{~}{\underset{u}{𝐺}}`$: $`=\underset{u}{}+\mathrm{\Gamma }`$, where the connection $`\mathrm{\Gamma }(\theta )`$ in terms of Ricci rotation coefficients reads $`\mathrm{\Gamma }_{(\lambda \alpha )}(\theta )={\displaystyle \frac{1}{4}}\mathrm{\Delta }_{(\lambda \alpha )(i,l)(m,p)}\gamma _f^{(i,l)}\gamma _f^{(m,p)},\overline{\mathrm{\Gamma }}_{(\lambda \alpha )}(\theta )={\displaystyle \frac{1}{4}}\mathrm{\Delta }_{(\lambda \alpha )(i,l)(m,p)}\gamma _f^{(m,p)}\gamma _f^{(i,l)}.`$
According to the general gauge principle , the physical system of the fields $`\mathrm{\Psi }(u)`$ is required to be invariant under the finite local gauge transformations
$$\begin{array}{c}\mathrm{\Psi }^{}(u)=U_R\mathrm{\Psi }(u),\left(g(u)\mathrm{\Psi }(u)\right)^{}=U_R\left(g(u)\mathrm{\Psi }(u)\right),U_R=R(𝐁_f^{})U_fR^1(𝐁_f),\hfill \end{array}$$
(B.1.5)
of the Lie group of gravitation $`G_R(U_R)`$ generated by G, where the gauge field $`𝐁_f(u_f)`$ is transformed under G in standard form. The physical meaning of the general principle is as follows: one has conventional G-gauge theory on flat manifold in terms of curviliniear coordinates if curvature tensor is zero, to which the zero vector eq.(B.1.3) is corresponded. Otherwise it yields the gravitation interaction.
Out of a set of arbitrary curvilinear coordinates in $`\stackrel{~}{\underset{u}{𝐺}}`$ the real curvilinear coordinates may be distinguished, which satisfy eq.(B.1.2) under all the possible Lorentz and gauge transformations. There is a single -valued conformity between corresponding tensors with various suffixes on $`\stackrel{~}{\underset{u}{𝐺}}`$ and $`\underset{u}{𝐺}`$. While, each index transformation is incorporated with the function $`\psi `$. The transformation of real curvilinear coordinates $`uu^{}`$ is due to some Lorentz $`(\mathrm{\Lambda })`$ and gauge $`(B_fB_f^{})`$ transformations
$$\frac{u^{}}{u}=\psi (B_f^{})\psi (B_f)\mathrm{\Lambda }.$$
(B.1.6)
There would then exist preferred systems and group of transformations of real curvilinear coordinates in the $`\stackrel{~}{\underset{u}{𝐺}}`$. The wider group of transformations of arbitrary curvilinear coordinates in the $`\stackrel{~}{\underset{u}{𝐺}}`$ would then be of no consequence for the field dynamics. A straightforward calculation gives the reflection matrix
$$R(u,u_f)=R_f(u_f)R_g(u)=\mathrm{exp}\left[i\mathrm{\Theta }_f(u_f)\mathrm{\Theta }_g(u)\right],$$
(B.1.7)
where
$$\mathrm{\Theta }_f(u_f)=g_0^{u_f}𝐁(u_f)𝑑u_f,\mathrm{\Theta }_g(u)=\frac{1}{2}_0^uR_f^+\{g\mathrm{\Gamma }R_f,gdu\}.$$
(B.1.8)
Then
$$S(B_f)=\frac{1}{8K}\psi \{\stackrel{~}{R}^+gR,\gamma ^0\}=inv,$$
(B.1.9)
where $`K=\stackrel{~}{R}^+R=\stackrel{~}{R}_g^+R_g=1.`$ and $`\stackrel{~}{U}_R^+U_R=\gamma ^0U_R^+\gamma ^0U_R=1.`$ The Lagrangian of shadow Fermi field may be written
$$\begin{array}{c}L_f(u_f)=J_\psi L(u)=\hfill \\ J_\psi \left\{S(B_f)\frac{i}{2}\left[\overline{\mathrm{\Psi }}_f(u_f)\gamma _fD\mathrm{\Psi }_f(u_f)\left(D\overline{\mathrm{\Psi }}_f(u_f)\right)\gamma _f\mathrm{\Psi }_f(u_f)\right]m\overline{\mathrm{\Psi }}_f(u_f)\mathrm{\Psi }_f(u_f)\right\}.\hfill \end{array}$$
(B.1.10)
provided by $`J_\psi =\psi \sqrt{g}(1+2<e^f,\chi ^f>+<\chi ^f,\chi ^f>)^{1/2}.`$ The Lagrangian $`L(u)`$ of the field $`\mathrm{\Psi }(u)`$ reads
$$\sqrt{g}L(u)=\frac{\sqrt{g}}{2}\left\{i\overline{\mathrm{\Psi }}(u)g(\underset{u}{}\mathrm{\Gamma })\mathrm{\Psi }(u)+i\overline{\mathrm{\Psi }}(u)(\underset{u}{\stackrel{}{}}\overline{\mathrm{\Gamma }})g\overline{\mathrm{\Psi }}(u)+2m\overline{\mathrm{\Psi }}(u)\mathrm{\Psi }(u)\right\},$$
(B.1.11)
yielding the field equations
$$\left[ig(\underset{u}{}\mathrm{\Gamma })m\right]\mathrm{\Psi }(u)=0,\overline{\mathrm{\Psi }}(u)\left[i(\underset{u}{\stackrel{}{}}\overline{\mathrm{\Gamma }})gm\right]=0.$$
(B.1.12)
## Appendix C
## 1 The Solution of Wave Equation of Distorted Structure
To solve the equation (B.1.12)
$$\left[ig(\underset{u}{}\mathrm{\Gamma })m\right]\mathrm{\Psi }(u)=0,$$
(C.1.1)
we transform it into
$$\{^2m^2(g\mathrm{\Gamma })^2+2(\mathrm{\Gamma }g)+(g)(g\mathrm{\Gamma })\}\mathrm{\Psi }=0,$$
(C.1.2)
where we abbreviate the indices $`(\lambda \alpha )`$ by the single symbol $`\mu `$, and Latin indices $`(im)(i=\pm ,m=1,2,3)`$ by $`i`$, also denote $`\widehat{𝑝}_u\widehat{p}`$ and
$$\begin{array}{c}\frac{1}{2}\sigma ^{\mu \nu }F_{\mu \nu }=(g)(g\mathrm{\Gamma })(\mathrm{\Gamma }),(g)(g\mathrm{\Gamma })=g^\mu g^\nu _\mu \mathrm{\Gamma }_\nu ,\hfill \\ \frac{1}{2}\sigma ^{\mu \nu }[\mathrm{\Gamma }_\mu ,\mathrm{\Gamma }_\nu ]=(g\mathrm{\Gamma })^2\mathrm{\Gamma }^2,^2=^\mu _\mu ,\mathrm{\Gamma }^2=\mathrm{\Gamma }^\mu \mathrm{\Gamma }_\mu ,\hfill \\ 2g^{\mu \nu }=\{g^\mu ,g^\nu \},2\sigma ^{\mu \nu }=[g^\mu ,g^\nu ],\underset{\mu \nu }{𝐹}=_\mu \mathrm{\Gamma }_\nu _\nu \mathrm{\Gamma }_\mu .\hfill \end{array}$$
(C.1.3)
We are looking for a solution given in the form $`\mathrm{\Psi }=e^{ipu}F(\phi )`$, where $`p_\mu `$ is a constant sixvector $`pu=p_\mu u_\mu `$, and admit that the field of distortion may be switched on at $`u_0=\mathrm{}`$ smoothly. Then the function $`\mathrm{\Psi }`$ must match onto the wave function of ordinary regular structure. Smoothness requires that the numbers $`p_\mu `$ become the components of link momentum of regular structure and satisfy the boundary condition $`p_\mu p_\mu =m^2=p_\eta ^2`$. Due to it we cancel unwanted solutions and clarify the normalization of wave functions
$$\mathrm{\Psi }_p^{}^{}\mathrm{\Psi }_pd^3u=\overline{\mathrm{\Psi }}_p^{}\gamma ^0\mathrm{\Psi }_pd^3u=(2\pi )^3\delta (\stackrel{}{p^{}}\stackrel{}{p}).$$
(C.1.4)
We suppose that at $`\sqrt{g}1`$ the gradient of the function $`\phi `$ reads
$$_\mu \phi =V_\mu ^ik_i,^\mu \phi =V_i^\mu k^i,$$
where $`k_i`$ are arbitrary constant numbers satisfying the condition $`k_ik_i=0`$. Thus $`^\mu \phi _\mu \phi =(V_i^\mu V_\mu ^j)k^ik_j=0`$. The eq.(C.1.2) gives rise to $`F^{}=A(\theta )F`$, where $`(\mathrm{})^{}`$ stands for the derivative with respect to $`\phi `$, and
$$A(\theta )=\frac{2i(p\mathrm{\Gamma })+m^2p^2+(g\mathrm{\Gamma })^2(g)(g\mathrm{\Gamma })}{2i(kVp)(kDV)};$$
(C.1.5)
where
$$\begin{array}{cc}(kVp)=k^iV_i^\mu p_\mu ,(kDV)=k^iD_\mu V_i^\mu ,D_\mu =_\mu 2\mathrm{\Gamma }_\mu ,\hfill & \\ p^2=p^\mu p_\mu =g^{\mu \nu }(\theta )p_\mu p_\nu (kVdu)=k_iV_\mu ^idu^\mu .\hfill & \end{array}$$
(C.1.6)
We are interested in the right-handed eigenvectors $`F_r(r=1,2,3,4)`$ corresponding to eigenvalues $`\mu _r`$ of matrix $`A:AF_r=\mu _rF_r`$, which are the roots of polynomial characteristic equation
$$c(\mu )=(\mu IA)=0.$$
Thus, one gets $`F_r^{}=\mu _rF_r`$ and $`F={\displaystyle \underset{r=1}{\overset{4}{}}}F_r`$. Hence $`(\mathrm{ln}F)^{}={\displaystyle \underset{r=1}{\overset{4}{}}}\mu _r=trA`$ and $`(\mathrm{ln}F)^{}=X_R(\theta )iX_J(\theta )`$, provided
$$\begin{array}{c}X_R(\theta )=trA_R(\theta )=tr\left\{\frac{(kDV)\left[m^2p^2+(g\mathrm{\Gamma })^2(g)(g\mathrm{\Gamma })\right]+4(kVp)(p\mathrm{\Gamma })}{(kDV)^2+4(kVp)^2}\right\},\hfill \\ X_J(\theta )=trA_J(\theta )=2tr\left\{\frac{(kVp)\left[m^2p^2+(g\mathrm{\Gamma })^2(g)(g\mathrm{\Gamma })\right]+(kDV)(p\mathrm{\Gamma })}{(kDV)^2+4(kVp)^2}\right\}.\hfill \end{array}$$
(C.1.7)
The solution of eq.(C.1.1) reads
$$F(\theta )=C\left(\frac{m}{E_u}\right)^{1/2}U\mathrm{exp}\{\chi _R(\theta )i\chi _J(\theta )\},$$
(C.1.8)
where $`C=1`$ is the normalization constant, $`U`$ is the constant bispinor, and
$$\chi _R(\theta )=_0^{u^\mu }(kVdu)X_R(\theta ),\chi _J(\theta )=_0^{u^\mu }(kVdu)X_J(\theta ).$$
(C.1.9)
## Appendix D
## 1 Field Aspect of the OMM
The quantum field and differential geometric aspects of the $`\widehat{G}_N`$ may be discussed on the analogy of $`\widehat{G}_{N=1}`$. Here we turn only to some points of the field aspect. A Lagrangian of free field reads
$$\stackrel{~}{L}_0(D)=\frac{i}{2}\{\overline{\mathrm{\Psi }}_e(\zeta ){}_{}{}^{i}\gamma _{}^{(\lambda ,\mu ,\alpha )}\underset{i}{}{}_{(\lambda ,\mu ,\alpha )}{}^{}\mathrm{\Psi }_{e}^{}(\zeta )\underset{i}{}{}_{(\lambda ,\mu ,\alpha )}{}^{}\overline{\mathrm{\Psi }}_{e}^{}(\zeta ){}_{}{}^{i}\gamma _{}^{(\lambda ,\mu ,\alpha )}\mathrm{\Psi }_e(\zeta )\},$$
(D.1.1)
where we have adopted the following convention:
$$\begin{array}{c}\mathrm{\Psi }_e(\zeta )=e\mathrm{\Psi }(\zeta )=\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)\mathrm{\Psi }(\zeta ),\overline{\mathrm{\Psi }}_e(\zeta )=e\overline{\mathrm{\Psi }}(\zeta ),\overline{\mathrm{\Psi }}(\zeta )=\mathrm{\Psi }^+(\zeta )\gamma ^0,\hfill \\ {}_{}{}^{i}\gamma _{}^{(\lambda ,\mu ,\alpha )}={}_{}{}^{i}\stackrel{~}{O}_{}^{\lambda ,\mu }\stackrel{~}{\sigma }^\alpha ,{}_{}{}^{i}\stackrel{~}{O}_{}^{\lambda ,\mu }=\frac{1}{\sqrt{2}}\left(\nu _i\xi _0\stackrel{~}{O}^\mu +\epsilon _\lambda \xi {}_{}{}^{i}\stackrel{~}{O}_{}^{\mu }\right),\hfill \\ \epsilon _\lambda =\{\begin{array}{cc}1& \lambda =1\\ 1& \lambda =2\end{array},<\nu _i,\nu _j>=\delta _{ij},\{{}_{}{}^{i}\stackrel{~}{O}_{}^{\lambda },{}_{}{}^{j}\stackrel{~}{O}_{}^{\mu }\}=\delta _{ij}{}_{}{}^{}\delta _{}^{\lambda \mu },\hfill \\ \stackrel{~}{O}^\mu =\frac{1}{\sqrt{2}}\left(\xi _0+\epsilon _\mu \xi \right),\stackrel{~}{O}^\lambda ={}_{}{}^{}\delta _{}^{\lambda \mu }\stackrel{~}{O}_\mu =(\stackrel{~}{O}_\lambda )^+,{}_{}{}^{i}\stackrel{~}{O}_{}^{\mu }=\frac{1}{\sqrt{2}}\left(\xi _{0i}+\epsilon _\mu \xi _i\right),\hfill \\ \underset{i}{}{}_{(\lambda ,\mu ,\alpha )}{}^{}=/{}_{}{}^{i}\zeta _{}^{(\lambda ,\mu ,\alpha )},\xi _0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\xi =\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\hfill \\ \xi _{0}^{}{}_{}{}^{2}=\xi ^2=\xi _{0i}^{}{}_{}{}^{2}=\xi _i^2=1,\{\xi _0,\xi \}=\{\xi _0,\xi _{0i}\}=\{\xi _0,\xi _i\}=\hfill \\ =\{\xi ,\xi _{0i}\}=\{\xi ,\xi _i\}=\{\xi _{0i},\xi _j\}_{ij}=\{\xi _{0i},\xi _{0j}\}_{ij}=\{\xi _i,\xi _j\}_{ij}=0.\hfill \end{array}$$
(D.1.2)
Field equations are written
$$\begin{array}{c}(\underset{\eta }{\widehat{p}}m)\underset{\eta }{\Psi }(\eta )=0,\overline{\underset{\eta }{\Psi }}(\eta )(\underset{\eta }{\widehat{p}}m)=0,\hfill \\ (\underset{u}{\widehat{p}}m)\underset{u}{\Psi }(u)=0,\overline{\underset{u}{\Psi }}(u)(\underset{u}{\widehat{p}}m)=0,\hfill \end{array}$$
(D.1.3)
where
$$\begin{array}{cc}\underset{\eta }{\widehat{p}}=i\underset{\eta }{\widehat{}},\underset{u}{\widehat{p}}=i\underset{u}{\widehat{}},\underset{u}{\widehat{}}={}_{}{}^{i}\gamma _{}^{(\lambda \alpha )}\underset{u_i}{}_{(\lambda \alpha )},\underset{\eta }{}{}_{(\lambda \alpha )}{}^{}=/\eta ^{(\lambda \alpha )},\underset{u_i}{}{}_{(\lambda \alpha )}{}^{}=/u_i^{(\lambda \alpha )},\hfill & \\ {}_{}{}^{i}\underset{\eta }{𝛾}_{}^{(\lambda \alpha )}={}_{}{}^{i}\stackrel{~}{\underset{\eta }{𝑂}}_{}^{\lambda }\stackrel{~}{\sigma }^\alpha =\nu _i\xi _0\gamma ^{(\lambda \alpha )}=\nu _i\xi _0\stackrel{~}{O}^\lambda \stackrel{~}{\sigma }^\alpha ,\hfill & \\ {}_{}{}^{i}\underset{u}{𝛾}_{}^{(\lambda \alpha )}={}_{}{}^{i}\stackrel{~}{\underset{u}{𝑂}}_{}^{\lambda }\stackrel{~}{\sigma }^\alpha =\xi {}_{}{}^{i}\gamma _{}^{(\lambda \alpha )}=\xi {}_{}{}^{i}\stackrel{~}{O}_{}^{\lambda }\stackrel{~}{\sigma }^\alpha ,\hfill & \\ \left(\gamma ^{(\lambda \alpha )}\right)^+={}_{}{}^{}\delta _{}^{\lambda \tau }\delta ^{\alpha \beta }\gamma ^{(\tau \beta )}=\gamma _{(\lambda \alpha )},\left({}_{}{}^{i}\underset{u}{𝛾}_{}^{(\lambda \alpha )}\right)^+={}_{}{}^{i}\underset{u}{𝛾}_{(\lambda \alpha )}^{}.\hfill & \end{array}$$
(D.1.4)
The state of free ordinary structure of $`{}_{}{}^{i}u`$-type with the given values of link momentum $`\underset{u_i}{𝑝}`$ and spin projection $`s_i`$ is described by means of plane wave. It is necessary to consider also the solution of negative frequencies with the normalized bispinor amplitude.
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# 1 Introduction
## 1 Introduction
There is a great surge of interest recently on noncommutativity along space-time directions in string/M theory . This noncommutativity implies a stringy uncertainty principle which appears as a special case of a fundamental one in string/M theory as advocated in . A field theory defined on such a noncommutative geometry cannot be unitary and therefore if there exists a decoupled theory on such geometry in string/M theory, it cannot be a field theory. In the context of D-branes with pure electric flux, it was shown in that such a decoupled theory indeed exists and is a noncommutative open string (NCOS) living on the brane. In obtaining such a decoupled theory, the electric field must be close to its critical value such that it almost balances the original string tension. We therefore end up effectively with an almost tensionless string. Such a tiny string tension defines a new energy scale for the decoupled NCOS as well as the scale for the noncommutativity<sup>1</sup><sup>1</sup>1We can scale the original tension to infinity while keeping the new effective tension fixed.. In this decoupling limit, the closed string coupling blows up while the coupling for the decoupled NCOS remains fixed and therefore is small in comparison with the infinitely large closed string coupling. This also implies that the transition of NCOS into closed strings cannot occur easily, indicating that the NCOS decouples from the bulk closed strings.
The same conclusion for the existence of NCOS was also obtained in but from a rather different viewpoint. Consider a 4-dimensional noncommutative Yang-Mills (NCYM) which is the decoupled field theory of D3 branes in a purely magnetic field background. The decoupling limit for this theory requires the closed string coupling to scale to zero. When the gauge coupling for the NCYM is large, the natural way to deal with this theory is to go to its S-dual description. The authors in asked: What is the S-dual of this NCYM? A worldvolume magnetic field is mapped to an electric field under S-duality<sup>2</sup><sup>2</sup>2If one instead uses a $`B`$-field with a nonvanishing spatial component, one cannot reach the conclusion directly from the worldvolume point of view since a spatial $`B`$-field under S-duality turns into a spatial RR 2-form field. One often says that D3-branes (in general Dp-branes) with spatial nonvanishing $`B`$-field give rise to NCYM in the corresponding decoupling limit. From the field theory (or open string point of view) side, this is correct since all we need for the NCYM is the closed string metric, closed string coupling and the $`B`$-field as demonstrated in . However, from the gravity description of D3-branes with spatial nontrivial $`B`$-field, we must also have D-strings present such that the resulting configuration is BPS. These D-strings require a nonvanishing space-time component of RR 2-form. Under S-duality, this nonvanishing component becomes a $`B`$-field which is needed for the space-time noncommutativity as demonstrated in . So this seems to indicate that the gravity description of NCYM contains more information than the field theory description. For example, in order to have the correct decoupling limit for NCOS, we have to perform the non-linear S-duality on the worldvolume while on the gravity side the S-duality is linear. This is the precise reason that we chose to work on a concrete gravity system of ((F, D1), D3) bound state for the decoupled NCOS with noncommutativity in both space-time and space-space directions in . But in the prsent paper, we choose to work in the hard way, i.e. from the worldvolume side.. Moreover, it was found that the decoupling limit in the S-dual theory exists only when the electric field attains a critical value. Also, in this case the original vanishing closed string coupling is transformed to its inverse in the S-dual, therefore becoming infinitely large. So, according to what has been obtained in , the S-dual of NCYM is a NCOS.
The NCYM used in arises as a decoupled theory of D3 branes with a pure magnetic flux. We can have decoupled NCYM from D3 branes with both electric and magnetic fields as discussed in and further with nonvanishing axion which will be discussed in this paper<sup>3</sup><sup>3</sup>3While we were writing up this paper, we became aware a paper in which the $`SL(2,Z)`$ duality of the gravity dual description of ((F, D1), D3) bound state, corresponding to $`𝐄||𝐁`$ case considered here, appeared on the net.. We can also have decoupled NCOS from D3 branes in the presence of both electric and magnetic fields as discussed in and further with a nonvanishing axion which, for the $`𝐄||𝐁`$ case, has been discussed in . One may wonder if the S-dual or in general a $`SL(2,Z)`$ dual of NCYM (or NCOS) always give a NCOS (or NCYM). In , the present authors showed that the S-dual of NCYM does not in general give a NCOS. In for the case E$`||`$B, it was shown that a NCOS is mapped to another NCOS under a $`SL(2,Z)`$ tranformation if the value of axion is irrational. Only for rational value of axion, a NCOS is physically equivalent to a NCYM.
We here intend to give a systematic study of the $`SL(2,Z)`$ dual of NCYM/NCOS which arises from D3 branes with nonvanishing axion and in the presence of both electric and magnetic fields in the respective decoupling limits. We will consider both of E$`||`$B and E$``$B cases. Our plan is as follows: In the following section, we give explicit $`SL(2,Z)`$ transformation rules for the worldvolume gauge fields and the generalized worldvolume gauge fields in the most general background. In section 3, we will calculate quantities relevant for decoupling limits for NCYM/NCOS in two versions related by $`SL(2,Z)`$ duality using Seiberg-Witten relations. We consider both E$``$B and E$`||`$B cases. In section 4, we will discuss various decoupling limits for NCOS/NCYM and the $`SL(2,Z)`$ dual of NCYM/NCOS. We conclude this paper in section 5.
## 2 D3 Branes and Nonperturbative $`SL(2,Z)`$ Symmetry
Type IIB string theory is conjectured to have a nonperturbative quantum $`SL(2,Z)`$ symmetry. This symmetry implies that two Type IIB string theories related by an $`SL(2,Z)`$ transformation are physically equivalent. Among various dynamical objects in this theory, D3 branes play a special role in the sense that this object, like type IIB string itself, is invariant under the $`SL(2,Z)`$. For the purpose of this paper, we need to consider only the bosonic sector of the low energy effective action for a single D3 brane coupled to a most general background. This is described by a Born-Infeld type action in Einstein-frame as
$$S_4=\frac{1}{(2\pi )^3\alpha ^2}d^4x\sqrt{det(g^E+e^{\varphi /2})}+\frac{1}{(2\pi )^3\alpha ^2}d^4x(Ce^{})_4,$$
(1)
where $`_2=2\pi \alpha ^{}F_2+B_2`$ and
$$(Ce^{})_4=C_4+C_2+\frac{1}{2}C_0_2_2.$$
(2)
In the above, the 2-form $`F_2`$ is the worldvolume $`U(1)`$ field strength, the worldvolume metric is the pullback of spacetime metric, $`B_2`$ is the pullback of spacetime NSNS 2-form potential and $`C_n`$ is the pullback of spacetime RR $`n`$-form potential. In the following, we take $`\mu ,\nu =0,1,2,3`$ as worldvolume indices and we will denote $`C_0=\chi `$ for notational convenience.
For vanishing $`C_2`$ and $`B_2`$, it was shown in that the equations of motion from the above action are just special forms of a more general action which possess a classical $`SL(2,R)`$ symmetry provided the dilaton and the axion parametrize the coset $`SL(2,R)/SO(2)`$. It was further shown in that with the inclusion of both $`B_2`$ and $`C_2`$, the equations of motion still have such $`SL(2,R)`$ symmetry provided the metric, dilaton, axion, $`B_2`$, $`C_2`$ and $`C_4`$ transform according to the rules determined from the type IIB supergravity, i.e.,
$`g_{\mu \nu }^Eg_{\mu \nu }^E,C_4C_4,\lambda {\displaystyle \frac{a\lambda +b}{c\lambda +d}},`$
$`\left(\begin{array}{c}B_2\\ C_2\end{array}\right)\left(\mathrm{\Lambda }^1\right)^T\left(\begin{array}{c}B_2\\ C_2\end{array}\right),`$ (7)
where the complex scalar<sup>4</sup><sup>4</sup>4Alternatively, the dilaton and axion can be used to parameterize the coset $`SL(2,R)/SO(2)`$ as
$$=\left(\begin{array}{cc}\chi ^2+e^{2\varphi }& \chi \\ \chi & 1\end{array}\right)e^\varphi .$$
Then under $`SL(2,R)`$, $`\mathrm{\Lambda }\mathrm{\Lambda }^T`$ with $`SL(2,R)`$ matrix $`\mathrm{\Lambda }`$ defined in (8). $`\lambda =\chi +ie^\varphi `$, $`\mathrm{\Lambda }`$ is a $`2\times 2`$ $`SL(2,R)`$ matrix defined as
$$\mathrm{\Lambda }=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),adbc=1,$$
(8)
and ‘T’ denotes the transpose of the matrix.
Our goal in this section is to express the transformed $`F_2`$ (or $`_2`$) under $`SL(2,R)`$ in terms of the original fields in a simple way. As is understood, the $`SL(2,R)`$ symmetry is manifest only on equations of motion and it rotates between equations of motion and Bianchi identities for $`F_2`$. Let us define a quantity $`K^{\mu \nu }`$ as
$$\sqrt{detg^E}\frac{K^{\mu \nu }}{2\pi }=\frac{\delta S_4}{\delta F_{\mu \nu }}.$$
(9)
Note that
$$det\left(g^E+e^{\varphi /2}\right)=\left(detg^E\right)\left(1+\frac{1}{2}e^\varphi ^2\frac{1}{16}e^{2\varphi }\left(\right)^2\right),$$
(10)
where $``$ denotes the Hodge-dual on the brane. With the above, we have
$$2\pi \alpha ^{}K_{\mu \nu }=\frac{e^\varphi _{\mu \nu }\frac{1}{4}e^{2\varphi }()()_{\mu \nu }}{\sqrt{1+\frac{1}{2}e^\varphi ^2\frac{1}{16}e^{2\varphi }\left(\right)^2}}+(C)_{\mu \nu }+\chi ()_{\mu \nu }.$$
(11)
With the above expression, the equation of motion for gauge potential $`A`$ ($`F_2=dA`$) is
$$dK_2=0.$$
(12)
So combining with Bianchi identity $`dF_2=0`$, we have
$$d\left(\begin{array}{c}F_2\\ K_2\end{array}\right)=0.$$
(13)
Given any solution of the above equation for $`F_2`$ ($`K_2`$ is given through (11)), it appears that we could obtain another solution from this through a global $`GL(2,R)`$ rotation. But, since D3-branes appear as sources to the bulk gravity, the energy-momentum tensor due to this source must be kept invariant under this symmetry since the Einstein-frame metric is inert to this symmetry. Further the equations of motion for various potentials in the bulk spacetime should be transformed covariantly under this transformation when the D3 brane source is considered. The global symmetry for the bulk gravity therefore restricts us to have only a global classical $`SL(2,R)`$ rather than $`GL(2,R)`$ for the D3 brane. With the D3 brane as source, we can deduce from equations of motion for $`B_2`$ and $`C_2`$, that $`(F_2,K_2)`$ transform in the same way as $`(B_2,C_2)`$ under $`SL(2,R)`$, i.e.,
$$\left(\begin{array}{c}F_2\\ K_2\end{array}\right)\left(\mathrm{\Lambda }^1\right)^T\left(\begin{array}{c}F_2\\ K_2\end{array}\right).$$
(14)
We can define a generalized 2-form $`𝒦_2`$ as
$$𝒦_2=2\pi \alpha ^{}K_2+C_2,$$
(15)
analogous to $`=2\pi \alpha ^{}F_2+B_2`$. Given the transformations for $`(F_2,K_2)`$ and $`(B_2,C_2)`$ under $`SL(2,R)`$, we have
$$\left(\begin{array}{c}_2\\ 𝒦_2\end{array}\right)\left(\mathrm{\Lambda }^1\right)^T\left(\begin{array}{c}_2\\ 𝒦_2\end{array}\right).$$
(16)
This is the key equation which we will use in the following section. With $`K_2`$ given in (11), (15) can be re-expressed as
$$𝒦_2=\frac{e^\varphi ()_2+\frac{1}{4}e^{2\varphi }()_2}{\sqrt{1+\frac{1}{2}e^\varphi ^2\frac{1}{16}e^{2\varphi }\left(\right)^2}}+\chi _2.$$
(17)
The above equation implies the following constraint which generalizes the one given in as
$$(_2,𝒦_2)\left(\begin{array}{c}_2\\ 𝒦_2\end{array}\right)=0,$$
(18)
which is useful for proving the invariance of the energy-momentum tensor.
In the following section, we will use the above equations (16) and (17) for calculating the open string metric and noncommutativity parameters from the Seiberg-Witten relations for both $`𝐄||𝐁`$ and $`𝐄𝐁`$ cases.
## 3 Seiberg-Witten Setup
In this section, we calculate the effective open string metric and noncommutativity parameters for both E$``$B and E$`||`$B cases using Seiberg-Witten formulae. The effective open string metric is
$$G_{\mu \nu }=g_{\mu \nu }(g^1)_{\mu \nu },$$
(19)
and the anti-symmetric noncommutativity parameter is
$$\mathrm{\Theta }^{\mu \nu }=2\pi \alpha ^{}\left(\frac{1}{g+}\right)_A^{\mu \nu },$$
(20)
where ‘A’ denotes the antisymmetric part. The effective open string coupling $`G_s`$ is related to the closed string coupling $`g_s=e^\varphi `$ through the following relation:
$$G_s=g_s\left(\frac{detG}{det(g+)}\right)^{1/2}$$
(21)
As usual, we assume in the following that $`_2`$ is entirely given by the worldvolume field $`F_2`$ and set the NSNS $`B_2`$ to zero. In other words, we trade NSNS $`B_2`$ for the worldvolume $`F_2`$ through a gauge transformation. For either E$``$B or E $`||`$B case, we calculate the above open string quantities from the relevant closed quantities and the worldvolume $``$ in two versions related by $`SL(2,Z)`$-duality. From now on, we limit ourselves to the non-perturbative quantum $`SL(2,Z)`$ symmetry rather than the classical $`SL(2,R)`$. In other words, we consider physically equivalent theories related by $`SL(2,Z)`$. For convenience, let us write down the transformed $`e^{\widehat{\varphi }}`$, $`\widehat{\chi }`$, $`\widehat{}_2`$ and $`\widehat{𝒦}_2`$ in terms of the corresponding original fields and integral $`SL(2,Z)`$ elements $`a,b,c,d`$ which satisfy $`adbc=1`$ as
$`e^{\widehat{\varphi }}=e^\varphi |c\lambda +d|^2,\widehat{\chi }={\displaystyle \frac{ac(\chi ^2+e^{2\varphi })+(ad+bc)\chi +bd}{|c\lambda +d|^2}},`$
$`\widehat{}_2=d_2c𝒦_2,\widehat{𝒦}_2=b_2+a𝒦_2.`$ (22)
where $`𝒦_2`$ is given by (17). The string metric is defined as $`g_{\mu \nu }=e^{\varphi /2}g_{\mu \nu }^E`$ and so we have
$$\widehat{g}_{\mu \nu }=g_{\mu \nu }|c\lambda +d|.$$
(23)
We always denote the corresponding quantities in the $`SL(2,Z)`$ dual with ‘hat’ over the letters as indicated above. Let us begin with the E$``$B case first.
### 3.1 E$``$B Case
Our starting point is to choose constant $`F_{\mu \nu }`$
$$_2=2\pi \alpha ^{}\left(\begin{array}{cccc}0& E& 0& 0\\ E& 0& B& 0\\ 0& B& 0& 0\\ 0& 0& 0& 0\end{array}\right),$$
(24)
and the constant closed string metric in string-frame as $`g_{\mu \nu }=\mathrm{diag}(g_0,g_1,g_2,g_3)`$ with $`g_0,g_1,g_2,g_3`$ all positive fixed parameters.
Using (19), we have the open string metric as
$$G_{\mu \nu }=\left(\begin{array}{cccc}g_0(1\stackrel{~}{E}^2)& 0& \sqrt{g_0g_2}\stackrel{~}{E}\stackrel{~}{B}& 0\\ 0& g_1(1\stackrel{~}{E}^2+\stackrel{~}{B}^2)& 0& 0\\ \sqrt{g_0g_2}\stackrel{~}{E}\stackrel{~}{B}& 0& g_2(1+\stackrel{~}{B}^2)& 0\\ 0& 0& 0& g_3\end{array}\right),$$
(25)
and using (20) the noncommutativity parameter as
$$\mathrm{\Theta }^{\mu \nu }=\frac{1}{1\stackrel{~}{E}^2+\stackrel{~}{B}^2}\left(\begin{array}{cccc}0& \stackrel{~}{E}/E_0& 0& 0\\ \stackrel{~}{E}/E_0& 0& \stackrel{~}{B}/B_0& 0\\ 0& \stackrel{~}{B}/B_0& 0& 0\\ 0& 0& 0& 0\end{array}\right).$$
(26)
The open string coupling can be obtained from eq.(21) as,
$$G_s=g_s\left(1\stackrel{~}{E}^2+\stackrel{~}{B}^2\right)^{1/2},$$
(27)
which implies that the critical field in this case is $`(1+\stackrel{~}{B}^2)^{1/2}`$. In the above, we have defined
$$\stackrel{~}{E}=\frac{E}{E_0},\stackrel{~}{B}=\frac{B}{B_0},$$
(28)
where the parameters $`E_0=\sqrt{g_0g_1}/(2\pi \alpha ^{})`$ and $`B_0=\sqrt{g_1g_2}/(2\pi \alpha ^{})`$.
One might think that the E$``$B case is simpler than the E$`||`$B case. On the contrary, it is a bit more complicated in both the decoupling limits (NCYM and NCOS) which will be studied in the following section and the $`SL(2,Z)`$ dual formulation. Let us derive the open string metric, the noncommutativity parameter and the open string coupling in the $`SL(2,Z)`$ dual. In order to calculate these quantities, we have to express $`\widehat{}_2`$ in terms of relevant quanties in the original version. For constant $`F_{\mu \nu }`$, the equation of motion is satisfied and so we can use the duality relation to calculate $`\widehat{}_2`$. In doing so, we first need to calculate $`𝒦_2`$ from (17). Thus we find,
$$𝒦_{\mu \nu }=\left(\begin{array}{cccc}0& \sqrt{g_0g_1}\stackrel{~}{E}\chi & 0& \sqrt{g_0g_3}\stackrel{~}{B}/G_s\\ \sqrt{g_0g_1}\stackrel{~}{E}\chi & 0& \sqrt{g_1g_2}\stackrel{~}{B}\chi & 0\\ 0& \sqrt{g_1g_2}\stackrel{~}{B}\chi & 0& \sqrt{g_2g_3}\stackrel{~}{E}/G_s\\ \sqrt{g_0g_3}\stackrel{~}{B}/G_s& 0& \sqrt{g_2g_3}\stackrel{~}{E}/G_s& 0\end{array}\right),$$
(29)
where we have used $`g_{\mu \nu }=g_s^{1/2}g_{\mu \nu }^E`$ and the relation between $`G_s`$ and $`g_s`$ given in (27) as well as the definitions for $`\stackrel{~}{E}`$ and $`\stackrel{~}{B}`$ given above. With the above and $`\widehat{}_2=d_2c𝒦_2`$, we have
$$\widehat{}_{\mu \nu }=\left(\begin{array}{cccc}0& \sqrt{g_0g_1}\stackrel{~}{E}(c\chi +d)& 0& \sqrt{g_0g_3}c\stackrel{~}{B}/G_s\\ \sqrt{g_0g_1}\stackrel{~}{E}(c\chi +d)& 0& \sqrt{g_1g_2}\stackrel{~}{B}(c\chi +d)& 0\\ 0& \sqrt{g_1g_2}\stackrel{~}{B}(c\chi +d)& 0& \sqrt{g_2g_3}c\stackrel{~}{E}/G_s\\ \sqrt{g_0g_3}c\stackrel{~}{B}/G_s& 0& \sqrt{g_2g_3}c\stackrel{~}{E}/G_s& 0\end{array}\right)$$
(30)
We notice from the expression of $`\widehat{}_{\mu \nu }`$ in (30) that we now have additional electric and magnetic fields pointing along negative $`x^3`$-direction and $`x^1`$-direction, respectively, in the $`SL(2,Z)`$ dual even though we originally had only electric field pointing along $`x^1`$-direction and magnetic field pointing along $`x^3`$-direction. This implies that we may have more noncommutative directions in the $`SL(2,Z)`$ dual. This differs from E$`||`$B case as we will see.
Using again Seiberg-Witten relations and after some tedious calculations, we find the new open string metric to have the form
$$\widehat{G}_{\mu \nu }=\frac{|cS+d|^2}{|c\lambda +d|}G_{\mu \nu },$$
(31)
and the noncommutativity parameters take the form,
$`\widehat{\mathrm{\Theta }}^{01}={\displaystyle \frac{(c\chi +d)}{|cS+d|^2}}\mathrm{\Theta }^{01},\widehat{\mathrm{\Theta }}^{03}={\displaystyle \frac{2\pi \alpha ^{}}{\sqrt{g_0g_3}}}{\displaystyle \frac{c\stackrel{~}{B}/G_s}{|cS+d|^2}},`$
$`\widehat{\mathrm{\Theta }}^{12}={\displaystyle \frac{(c\chi +d)}{|cS+d|^2}}\mathrm{\Theta }^{12},\widehat{\mathrm{\Theta }}^{23}={\displaystyle \frac{2\pi \alpha ^{}}{\sqrt{g_2g_3}}}{\displaystyle \frac{c\stackrel{~}{E}/G_s}{|cS+d|^2}},`$ (32)
where $`\lambda =\chi +i/g_s`$ defined before, $`S=\chi +i/G_S`$ (since $`𝐄𝐁=0`$ here), and $`G_{\mu \nu }`$, $`\mathrm{\Theta }^{01}`$ and $`\mathrm{\Theta }^{12}`$ are the original open string metric, noncommutative parameters given in (25) and (26), respectively. The open string coupling here is related to the original open string coupling as
$$\widehat{G}_s=|cS+d|^2G_s.$$
(33)
### 3.2 E $`||`$ B Case
We now take
$$=2\pi \alpha ^{}\left(\begin{array}{cccc}0& E& 0& 0\\ E& 0& 0& 0\\ 0& 0& 0& B\\ 0& 0& B& 0\end{array}\right),$$
(34)
and the string frame metric $`g_{\mu \nu }=\mathrm{diag}(g_1,g_1,g_2,g_2)`$ where $`g_1,g_2`$ are positive parameters. Using Seiberg-Witten relations, we have the open string metric as
$$G_{\mu \nu }=\left(\begin{array}{cccc}g_1(1\stackrel{~}{E}^2)& 0& 0& 0\\ 0& g_1(1\stackrel{~}{E}^2)& 0& 0\\ 0& 0& g_2(1+\stackrel{~}{B}^2)& 0\\ 0& 0& 0& g_2(1+\stackrel{~}{B}^2)\end{array}\right),$$
(35)
the noncommutativity parameters as,
$$\mathrm{\Theta }^{01}=\frac{\stackrel{~}{E}}{E_c(1\stackrel{~}{E}^2)},\mathrm{\Theta }^{23}=\frac{\stackrel{~}{B}}{B_0(1+\stackrel{~}{B}^2)},$$
(36)
and the open string coupling as
$$G_s=g_s(1\stackrel{~}{E}^2)^{1/2}(1+\stackrel{~}{B}^2)^{1/2}.$$
(37)
In the above, we have defined $`\stackrel{~}{E}=E/E_c`$ and $`\stackrel{~}{B}=B/B_0`$ with the critical field $`E_c=g_1/(2\pi \alpha ^{})`$ and $`B_0=g_2/(2\pi \alpha ^{})`$.
We now calculate the relevant open string quantities in the $`SL(2,Z)`$ dual. To do so, we need to have $`\widehat{}_2`$ and as before we calculate $`𝒦_{\mu \nu }`$ first. Using (17), we have
$$𝒦_{01}=g_1[\stackrel{~}{E}\chi \stackrel{~}{B}(1\stackrel{~}{E}^2)/G_s],𝒦_{23}=g_2[\stackrel{~}{B}\chi +\stackrel{~}{E}(1+\stackrel{~}{B}^2)/G_s)],$$
(38)
where we have used the relation between $`G_s`$ and $`g_s`$ given in (37), the definitions for $`\stackrel{~}{E}`$ and $`\stackrel{~}{B}`$ given earlier and $`g_{\mu \nu }=g_s^{1/2}g_{\mu \nu }^E`$. With the above, we have using eq.(22)
$$\widehat{}_{01}=g_1[\stackrel{~}{E}(c\chi +d)c\stackrel{~}{B}(1\stackrel{~}{E}^2)/G_s],\widehat{}_{23}=g_2[\stackrel{~}{B}(c\chi +d)+c\stackrel{~}{E}(1+\stackrel{~}{B}^2)/G_s)].$$
(39)
Using Seiberg-Witten relations (19)-(21), we now have the open string metric
$$\widehat{G}_{\mu \nu }=\frac{|cS+d|^2}{|c\lambda +d|}G_{\mu \nu },$$
(40)
the noncommutativity parameters
$$\widehat{\mathrm{\Theta }}^{01}=\frac{(c\chi +d)\mathrm{\Theta }^{01}\frac{c\stackrel{~}{B}}{E_{cr}G_s}}{|cS+d|^2},\widehat{\mathrm{\Theta }}^{23}=\frac{(c\chi +d)\mathrm{\Theta }^{23}\frac{c\stackrel{~}{E}}{B_0G_s}}{|cS+d|^2},$$
(41)
and the open string coupling
$$\widehat{G}_s=|cS+d|^2G_s.$$
(42)
In the above, $`G_s`$, $`\mathrm{\Theta }^{\mu \nu }`$ and $`G_{\mu \nu }`$ are the original open string coupling, noncommutativity parameters and open string metric, respectively. We have now $`S=\chi +\stackrel{~}{E}\stackrel{~}{B}/G_s+i/G_s`$ for which $`𝐄`$ and $`𝐁`$ contribute since $`𝐄𝐁0`$.
## 4 Decoupling Limits and $`SL(2,Z)`$ Duality
We are now ready to discuss the decouping limits for NCYM/NCOS and the $`SL(2,Z)`$ duality for the underlying decoupled theory. Before we discuss the noncommutative theory, we would like to address one question: Can an ordinary theory become a noncommutative theory through $`SL(2,Z)`$ duality? Our examination gives negative answer. Let us point out that there is a general rule regarding whether we can map a strongly coupled theory to a weakly coupled one through a $`SL(2,Z)`$ transformation or not for any theory either ordianry or noncommutative. The rule is: for rational $`\chi `$, we can have two physically equivalent theories which are strong-weak dual to each other while for irrational $`\chi `$, we do not have this. We first discuss $`𝐄𝐁`$ case and then $`𝐄||𝐁`$ case.
### 4.1 $`𝐄𝐁`$ Case
Let us begin with the decoupling limit for NCYM. To have a NCYM, we need to decouple not only the open string ending on the brane from the closed strings in the bulk but also the open string massive modes from the massless ones. So we need to send $`\alpha ^{}0`$. To have a sensible quantum theory, we need to fix the open string coupling and the open string metric in this limit. We also need to fix at least one nonvanishing spatial component of the noncommutative matrix. With these requirements and examining (25) and (26), we can naively have the following three limits:
* 1)
$$\alpha ^{}0,g_1=\left(\frac{\stackrel{~}{b}}{\alpha ^{}}\right)^2,\stackrel{~}{E}^2=1+\stackrel{~}{B}^2\left(\frac{\alpha ^{}}{\stackrel{~}{b}}\right)^2,g_s=G_s\frac{\stackrel{~}{b}}{\alpha ^{}}$$
(43)
with $`g_0,g_2,g_3`$ and $`\stackrel{~}{B}`$ fixed. For simplicity, we choose $`g_0\stackrel{~}{B}^2=1`$, $`g_2(1+\stackrel{~}{B}^2)=1`$ and $`g_3=1`$. So we have the metric
$$G_{\mu \nu }=\left(\begin{array}{cccc}1g_0(\alpha ^{}/\stackrel{~}{b})^2& 0& 1g_2(\alpha ^{}/\stackrel{~}{b})^2/2& 0\\ 0& 1& 0& 0\\ 1g_2(\alpha ^{}/\stackrel{~}{b})^2/2& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),$$
(44)
and the noncommutativity parameters<sup>5</sup><sup>5</sup>5We choose $`\stackrel{~}{B}`$ to be negative for definiteness. For positive $`\stackrel{~}{B}`$, the discussion and the conclusion are basically the same.
$$\mathrm{\Theta }^{01}=\mathrm{\Theta }^{12}=2\pi \stackrel{~}{b}|\stackrel{~}{B}|\sqrt{1+\stackrel{~}{B}^2}.$$
(45)
* 2)
$$\alpha ^{}0,\stackrel{~}{E}=\stackrel{~}{B}=\left(\frac{\stackrel{~}{b}}{\alpha ^{}}\right)^{1/2},g_0=g_2=\frac{\alpha ^{}}{\stackrel{~}{b}},$$
(46)
with $`g_1,g_3`$ and $`g_s`$ fixed. We here choose $`g_0/g_2=1`$ just for simplicity but in general we need only $`g_0/g_2`$ fixed. For simplicity, we also choose $`g_1=g_3=1`$. Now the open string metric is
$$G_{\mu \nu }=\left(\begin{array}{cccc}1\frac{\alpha ^{}}{\stackrel{~}{b}}& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 1+\frac{\alpha ^{}}{\stackrel{~}{b}}& 0\\ 0& 0& 0& 1\end{array}\right),$$
(47)
and the noncommutativity parameters are similar to those in 1) as
$$\mathrm{\Theta }^{01}=\mathrm{\Theta }^{12}=2\pi \stackrel{~}{b}.$$
(48)
So 1) and 2) are quite similar except for the closed string coupling $`g_s`$. $`g_s`$ blows up in the decoupling limit in 1) while it remains fixed in 2). This case corresponds to the light-like NCYM discussed in .
* 3)
$$\alpha ^{}0,g_1=g_2=\left(\frac{\alpha ^{}}{\stackrel{~}{b}}\right)^2,\stackrel{~}{B}=\frac{\stackrel{~}{b}}{\alpha ^{}},g_s=G_s\frac{\alpha ^{}}{\stackrel{~}{b}},$$
(49)
with $`g_0,g_3`$ and $`\stackrel{~}{E}`$ fixed. Here we choose $`g_1/g_2=1`$ for simplicity but in general we only need this ratio to be fixed. For the present case, we do not have the $`\stackrel{~}{E}1`$ requirement. It can be any fixed number or approaching zero. For simplicity we set $`g_3=1`$. This same decoupling limit for $`\stackrel{~}{E}1`$ was discussed in . With the above, we have the open string metric
$$G_{\mu \nu }=\left(\begin{array}{cccc}g_0(1\stackrel{~}{E}^2)& 0& \sqrt{g_0}\stackrel{~}{E}& 0\\ 0& 1& 0& 0\\ \sqrt{g_0}\stackrel{~}{E}& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),$$
(50)
and the nonvanishing noncommutativity parameter
$$\mathrm{\Theta }^{12}=2\pi \stackrel{~}{b}.$$
(51)
We point out that there are both space-time and space-space noncommutativities in 1) and 2) while there is only space-space noncommutativity in 3). The space-time noncommutativity arises because the electric field approaches the critical value in both cases. In general, one expects that the underlying decoupled theories are NCOS rather than a NCYM. At least for 2), the unitarity discussion given in seems to indicate that the resulting theory is a light-like NCYM. Actually, 1) has the same structure as in 2). So we expect that we might also have a light-like NCYM in the limit of 1). The arguments given in are: even though we have both $`\mathrm{\Theta }^{01}`$ and $`\mathrm{\Theta }^{12}`$ in appearance, if we choose light-like coordinates $`x^\pm =(x^0\pm x^2)/\sqrt{2}`$, work in the $`(x^+,x^1,x^{},x^3)`$-system and take $`x^+`$ as the light-cone time, we then have only nonvanishing noncommutativity parameter $`\mathrm{\Theta }^1`$, a space-space noncommutativity. In addition, the underlying theory is unitary, i.e., the inner product $`pp=p_\mu \mathrm{\Theta }^{\mu \rho }G_{\rho \sigma }\mathrm{\Theta }^{\sigma \nu }p_\nu `$ is never negative. Therefore, it appears that the underlying theory is a well-defined NCYM. Let us demonstrate this for both 1) and 2). We now express everything in the light-like coordinate $`(x^+,x^1,x^{},x^3)`$-system. Let us denote the corresponding cases as $`1^{(lc)})`$ and $`2^{(lc)})`$, respectively.
* $`1^{(lc)}`$): We now have the metric
$$G_{\mu \nu }^{(lc)}=\left(\begin{array}{cccc}2& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& \frac{(\alpha ^{}/\stackrel{~}{b})^2}{2g_0g_2}& 0\\ 0& 0& 0& 1\end{array}\right)$$
(52)
and the only nonvanishing noncommutativity parameter
$$\mathrm{\Theta }^{(lc)1}=\mathrm{\Theta }^{(lc)1}=2\sqrt{2}\pi \stackrel{~}{b}|\stackrel{~}{B}|\sqrt{1+\stackrel{~}{B}^2}.$$
(53)
One can check that indeed $`pp`$ is non-negative as $`pp=p_\mu \mathrm{\Theta }^{\mu \rho }G_{\rho \sigma }\mathrm{\Theta }^{\sigma \nu }p_\nu =(p_{})^2(\mathrm{\Theta }^1)^20`$ when $`\alpha ^{}0`$ is taken.
* $`2^{(lc)})`$: We have the open string metric
$$G_{\mu \nu }^{(lc)}=\left(\begin{array}{cccc}2& 0& \frac{\alpha ^{}}{\stackrel{~}{b}}& 0\\ 0& 1& 0& 0\\ \frac{\alpha ^{}}{\stackrel{~}{b}}& 0& 0& 0\\ 0& 0& 0& 1\end{array}\right),$$
(54)
and the nonvanishing noncommutativity parameter
$$\mathrm{\Theta }^{(lc)1}=\mathrm{\Theta }^{(lc)1}=2\sqrt{2}\pi \stackrel{~}{b}.$$
(55)
One can check again $`pp=(p_{})^2(\mathrm{\Theta }^{(lc)1})^20`$.
The above discussion seems to indicate that in terms of the light-like coordinates, $`1^{(lc)})`$ and $`2^{(lc)})`$ cases look no different from the case 3) above. According to the criterion given in , each of the field theories is unitary and each has a space-space noncommutativity. We therefore should call the decoupled field theories in $`1^{(lc)}`$) and $`2^{(lc)}`$) as light-like NCYM.
Let us discuss the $`SL(2,Z)`$ duality for each of the NCYM. Let us denote the corresponding cases as $`1^{(lc)})^{}`$, $`2^{(lc)})^{}`$ and $`3)^{}`$, respectively.
* $`1^{(lc)})^{}`$ We need to consider irrational $`\chi `$ and rational $`\chi `$ separately. a) If $`\chi `$ is irrational, $`|c\lambda +d|=(c\chi +d)`$ since $`g_s\mathrm{}`$ in the decoupling limit 1). $`|cS+d|^2=(c\chi +d)^2+c^2/G_s^2`$ remains fixed. Using (31), (32), (33) and the decoupling limit in 1), we have
$$\widehat{G}_{\mu \nu }^{(lc)}=\frac{|cS+d|^2}{c\chi +d}G_{\mu \nu }^{(lc)},\widehat{G}_s=|cS+d|^2G_s,\widehat{\mathrm{\Theta }}^{(lc)1}=\frac{c\chi +d}{|cS+d|^2}\mathrm{\Theta }^{(lc)1},$$
(56)
where $`G_{\mu \nu }^{(lc)}`$, $`G_s`$ and $`\mathrm{\Theta }^{(lc)1}`$ are the corresponding quantities in $`1^{(lc)})`$. So this theory looks similar to the original theory but it is always strongly coupled. The $`SL(2,Z)`$ duality here is not useful. b) If $`\chi `$ is rational, we can choose $`c\chi +d=0`$. Now $`|c\lambda +d|=|c|/g_s=|c|\alpha ^{}/(G_s\stackrel{~}{b})`$. $`|cS+d|^2=c^2/G_s^2`$. We have now
$$\widehat{G}_{\mu \nu }^{(lc)}=\frac{\stackrel{~}{b}}{\alpha ^{}}\frac{|c|]}{G_s}G_{\mu \nu }^{(lc)},\widehat{G}_s=c^2/G_s,$$
(57)
and all the noncommutativity parameters vanish. If the original light-like NCYM is strongly coupled, we end up with a physically equivalent weakly coupled theory defined on a commutative geometry. Also we have $`pp=0`$ since $`\mathrm{\Theta }^{(lc)\mu \nu }`$ vanish. So it appears that we end up with a light-like OYM. We will comment on this later in this subsection.
* $`2^{(lc)})^{}`$ We have also two cases: a) irrational $`\chi `$ and b) rational $`\chi `$.
a)Irrational $`\chi `$: Now $`|c\lambda +d|`$ is fixed since we have fixed $`g_s`$. So is $`|cS+d|^2`$. Therefore the metric $`\widehat{G}_{\mu \nu }^{(lc)}`$ is essentially the same as the original metric $`G_{\mu \nu }^{(lc)}`$. We also have
$$\widehat{\mathrm{\Theta }}^{(lc)1}=2\sqrt{2}\pi \stackrel{~}{b}\frac{c\chi +d}{|cS+d|^2},\widehat{\mathrm{\Theta }}^{(lc)3}=\frac{2\sqrt{2}\pi c\stackrel{~}{b}}{G_s|cS+d|^2}.$$
(58)
We again have only the space-space noncommutativity in the light-like coordinate system but we double the number of noncommutative pairs. This theory is also unitary since $`pp0`$. We again end up with a light-like NCYM. However, this NCYM is strongly coupled regardless of the coupling of the original theory. So the $`SL(2,Z)`$ duality is not that useful except for increasing the noncommutative directions.
b)Rational $`\chi `$: We can now choose $`c\chi +d=0`$. Since $`g_s`$ is fixed, we can obtain the corresponding quantities simply by setting $`c\chi +d=0`$ in a). We then have $`\widehat{\mathrm{\Theta }}^{(lc)1}=0`$ and $`\widehat{G}_s=c^2/G_s`$. So we end up with a weakly coupled light-like NCYM if the original light-like NCYM is strongly coupled. The number of noncommutative pairs remain the same but we change from $`\mathrm{\Theta }^{(lc)1}`$ to $`\widehat{\mathrm{\Theta }}^{(lc)3}`$ by the $`SL(2,Z)`$ duality.
* $`3)^{}`$ We have two cases: a) irrational $`\chi `$ and b) rational $`\chi `$.
a)Irrational $`\chi `$: $`|c\lambda +d|=|c|\stackrel{~}{b}/(\alpha ^{}G_s)`$. $`|cS+d|^2`$ remains fixed. We then have from (31) and (32)
$`\widehat{G}_{\mu \nu }={\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}}{\displaystyle \frac{G_s}{|c|}}|cS+d|^2G_{\mu \nu },\widehat{G}_s=|cS+d|^2G_s,`$
$`\widehat{\mathrm{\Theta }}^{03}={\displaystyle \frac{2\pi c\stackrel{~}{b}}{\sqrt{g_0}G_s|cS+d|^2}},\widehat{\mathrm{\Theta }}^{12}={\displaystyle \frac{2\pi \stackrel{~}{b}(c\chi +d)}{|cS+d|^2}},\widehat{\mathrm{\Theta }}^{23}={\displaystyle \frac{2\pi c\stackrel{~}{E}\stackrel{~}{b}}{G_s|cS+d|^2}},`$ (59)
Since $`\alpha ^{}G^1`$ is fixed, we therefore have NCOS. Depending on the value of $`\stackrel{~}{E}`$, we can have as many as three independent noncommutativity parameters. However, this theory is strongly coupled regardless of the original theory. So it is again not useful.
b)Rational $`\chi `$: For this case we can choose $`c\chi +d=0`$. Now $`|c\lambda +d|`$ remains the same as the above since $`g_s0`$ and $`|cS+d|^2=c^2/G_s^2`$ still remains fixed. So we still have $`\widehat{G}_{\mu \nu }\alpha ^{}G_{\mu \nu }`$ which implies that we still end up with a NCOS. But now $`\widehat{\mathrm{\Theta }}^{12}=0`$. We have
$$\widehat{\mathrm{\Theta }}^{03}=\frac{2\pi G_s\stackrel{~}{b}}{c\sqrt{g_0}},\widehat{\mathrm{\Theta }}^{23}=2\pi G_s\stackrel{~}{E}\stackrel{~}{b}/c.$$
(60)
We then have a weakly coupled NCOS if the original NCYM is strongly coupled. So now the $`SL(2,Z)`$ duality is useful.
So far it appears that $`1^{(lc)})`$ and $`2^{(lc)})`$ differ from 3) in that the former are light-like NCYM while the latter is usual NCYM. Also their $`SL(2,Z)`$ dualities are quite different. The former give either light-like NCYM or OYM while the latter gives the usual NCOS. There is also a big difference regarding the closed string coupling $`g_s`$. The former have either $`g_s`$ blowing up or fixed while the latter has vanishing $`g_s`$. In this aspect, these light-like NCYM are similar to the NCYM discussed in whose S-duality gives an OYM which is not well-defined because of the singular metric and infinitely large open string coupling. We will discuss the $`SL(2,Z)`$ duality of this kind of NCYM in the following subsection. There is a difference regarding the open string coupling. The $`SL(2,Z)`$ duality of the light-like NCYM has a finite (maybe large) open string coupling while that of the above mentioned NCYM discussed in has blowing up open string coupling. The reason for this is that we here consider $`𝐄𝐁`$ case and $`𝐄`$ and $`𝐁`$ have no contribution to the quantity $`S`$. But for $`𝐄||𝐁`$ case, we do have $`\stackrel{~}{E}\stackrel{~}{B}`$ contribution to the $`S`$ as indicated before which blows up in the decoupling limit for the above mentioned NCYM in .
Up to now we have avoided pointing out the underlying major difference between case $`1^{(lc)})`$ and $`2^{(lc)})`$ and case 3). The $`detG_{\mu \nu }^{(lc)}\alpha ^2`$ vanishes for the former two cases while $`detG_{\mu \nu }`$ is finite for case 3). In other words, the light-like NCYM, if they indeed exist, for the former two cases are defined on a zero-size spacetime, or singular spacetime, while the latter is a well-defined usual NCYM. Because of this, at least we have one component of $`\alpha ^{}G^{\mu \nu }`$ nonvanshing (the same is true in the $`SL(2,Z)`$ dual). One would say that the underlying theory may not be a field theory. One may wonder that the unitarity condition obtained from a one-loop analysis in is sufficient to show the existence of such light-like NCYM. Further study is needed.
One could have non-singular metric by rescaling the light-cone coordinate $`x^{}`$. For $`1^{(lc)})`$, if we rescale $`x^{}=(\stackrel{~}{b}/(\alpha ^{}\sqrt{g_0g_2}))\stackrel{~}{x}^{}`$, then we get the open string metric $`G_{\mu \nu }^{(lc)}=\mathrm{diag}(1/2,1,2,1)`$, which is non-singular with respect to the coordinates $`(x^+,x^1,\stackrel{~}{x}^{},x^3)`$. For $`2^{(lc)})`$, if we rescale $`x^{}=(\stackrel{~}{b}/\alpha ^{})x^{}`$, we have, with respect to $`(x^+,x^1,\stackrel{~}{x}^{},x^3)`$,
$$G_{\mu \nu }^{(lc)}=\left(\begin{array}{cccc}2& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\end{array}\right),$$
(61)
which is also non-singular. Note that the only noncommutativity parameter with respect to $`(x^+,x^1,x^{},x^3)`$ is $`\mathrm{\Theta }^1`$. Actually, we have the two-point function $`x^{}(\tau )x^1(0)=(i/2)\mathrm{\Theta }^1ϵ(\tau )`$. The scaling $`x^{}(1/\alpha ^{})\stackrel{~}{x}^{}`$ implies $`\stackrel{~}{x}^{}x^1\alpha ^{}\mathrm{\Theta }^1ϵ(\tau )0`$. Since now $`\alpha ^{}G^10`$ with respect to $`(x^+,x^1,\stackrel{~}{x}^{},x^3)`$, we therefore end up with light-like OYM<sup>6</sup><sup>6</sup>6We are not sure whether the resulting light-like OYM is physically equivalent to the original theory defined on a zero size 4-dimensional spacetime because the rescaling of $`x^{}`$ is singular. for both 1) and 2).
Because we end up with OYM for 1) and 2), their $`SL(2,Z)`$ dualities still give OYM as discussed at the outset of this section. We will not give the detail here.
We now move on to discuss possible NCOS limit and its $`SL(2,Z)`$ duality. To have decoupling limit for NCOS, we need to keep $`\alpha ^{}G^{\mu \nu }`$ and at least $`\mathrm{\Theta }^{01}`$ fixed when the limit $`EE_c`$ is taken with $`E_c`$ the critical field limit. In general, we do not need to send $`\alpha ^{}0`$ since the open string massive modes are not decoupled from its massless modes. However, it is convenient to choose the $`\alpha ^{}0`$ limit since we will study the $`SL(2,Z)`$ duality of the resulting NCOS which might be a field theory. Now we have a fixed $`\alpha _{eff}^{}`$ for the NCOS which is determined by the noncommutative scale.
For $`𝐄𝐁`$, the critical electric field limit is $`\stackrel{~}{E}^21+\stackrel{~}{B}^2`$. From the previous discussion for NCYM, we expect that for either fixed $`\stackrel{~}{B}`$ or infinitely large $`\stackrel{~}{B}`$ as $`\alpha ^{}0`$, we have similar complications here. We will discuss these cases elsewhere. We here focus on the limit $`\stackrel{~}{B}0`$ along with the above critical electric field limit as $`\alpha ^{}0`$ for NCOS. The relation between the effective open string coupling and the closed string coupling (27) implies $`\stackrel{~}{E}^2<1+\stackrel{~}{B}^2`$, so we should have in general $`\stackrel{~}{E}^2=1+\stackrel{~}{B}^2(\alpha ^{}/\stackrel{~}{b})^\delta `$ with $`\delta >0`$ and $`\stackrel{~}{b}`$ fixed. With the above discussion, we must also have $`\stackrel{~}{B}^2=(\alpha ^{}/\stackrel{~}{b}^{})^\beta `$ with $`\beta >0`$ and $`\stackrel{~}{b}^{}`$ fixed. For $`\beta >\delta `$, the effect of $`𝐁`$ simply drops out and we have purely electric field effect which has been discussed before and we will not repeat this case here. The only other case which gives $`G_{\mu \nu }\alpha ^{}`$ is $`\beta =\delta `$. We have two cases: a) $`\stackrel{~}{b}>\stackrel{~}{b}^{}`$ and b)$`\stackrel{~}{b}<\stackrel{~}{b}^{}`$. Let us discuss each in order. a) We now have the decoupling limit
$`\alpha ^{}0,g_0=\left[\left({\displaystyle \frac{\stackrel{~}{b}}{\stackrel{~}{b}^{}}}\right)^\delta 1\right]^1\left({\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}}\right)^{1\delta },g_1=\left({\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}}\right)^{1\delta },`$
$`g_2=g_3={\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}},g_s=G_s\left({\displaystyle \frac{\stackrel{~}{b}}{\alpha ^{}}}\right)^{\delta /2},\stackrel{~}{E}^2=1+\left[\left({\displaystyle \frac{\stackrel{~}{b}}{\stackrel{~}{b}^{}}}\right)^\delta 1\right]\left({\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}}\right)^\delta ,`$ (62)
with $`\stackrel{~}{B}`$ given above. We then have
$$G_{\mu \nu }=\frac{\alpha ^{}}{\stackrel{~}{b}}\left(\begin{array}{cccc}1& 0& \left[1(\stackrel{~}{b}^{}/\stackrel{~}{b})^\delta \right]^{1/2}& 0\\ 0& 1& 0& 0\\ \left[1(\stackrel{~}{b}^{}/\stackrel{~}{b})^\delta \right]& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),\mathrm{\Theta }^{01}=2\pi \stackrel{~}{b}\left[\left(\frac{\stackrel{~}{b}}{\stackrel{~}{b}^{}}\right)^\delta 1\right]^{1/2}.$$
(63)
We, therefore, have NCOS with nonvanishing $`\mathrm{\Theta }^{01}`$. We now consider case b)<sup>7</sup><sup>7</sup>7This case may be equivalent to the one studied in .. The decoupling limit for this case remains the same except for the scaling for $`g_0`$ which can be obtained by the following replacement:
$$\left(\frac{\stackrel{~}{b}}{\stackrel{~}{b}^{}}\right)^\delta 11\left(\frac{\stackrel{~}{b}}{\stackrel{~}{b}^{}}\right)^\delta .$$
(64)
Now we have
$$G_{\mu \nu }=\frac{\alpha ^{}}{\stackrel{~}{b}}\left(\begin{array}{cccc}1& 0& \left[(\stackrel{~}{b}^{}/\stackrel{~}{b})^\delta 1\right]^{1/2}& 0\\ 0& 1& 0& 0\\ \left[(\stackrel{~}{b}^{}/\stackrel{~}{b})^\delta 1\right]^{1/2}& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),$$
(65)
and the nonvanishing noncommutativity parameter $`\mathrm{\Theta }^{01}`$ which can be obtained from (63) by the same replacement as above. We again end up with a NCOS.
We denote the $`SL(2,Z)`$ duality of the above two cases as $`a)^{}`$ and $`b)^{}`$.
* $`a)^{}`$ We need to consider: 1) irrational $`\chi `$ and 2) rational $`\chi `$.
1)Irrational $`\chi `$: Now since $`g_s\mathrm{}`$, we have $`|c\lambda +d|=|c\chi +d|0`$ and $`|cS+d|^2`$ remains fixed. So we have
$`\widehat{G}_{\mu \nu }={\displaystyle \frac{|cS+d|^2}{|c\chi +d|}}G_{\mu \nu },\widehat{G}_s=|cS+d|^2G_s`$
$`\widehat{\mathrm{\Theta }}^{01}={\displaystyle \frac{c\chi +d}{|cS+d|^2}}\mathrm{\Theta }^{01},\widehat{\mathrm{\Theta }}^{23}={\displaystyle \frac{2\pi \stackrel{~}{b}c}{G_s|cS+d|^2}},`$ (66)
where $`G_{\mu \nu },G_s`$ and $`\mathrm{\Theta }^{01}`$ are the open string metric, noncommutativity parameter and open string coupling in a) above. Since $`\alpha ^{}\widehat{G}^1`$ is fixed, so we still have NCOS. This theory is strongly coupled and again the $`SL(2,Z)`$ duality is not useful.
2)Rational $`\chi `$: We can now choose $`c\chi +d=0`$. Then we have $`|c\lambda +d|=|c|/g_s=(|c|/G_s)(\alpha ^{}/\stackrel{~}{b})^{\delta /2}`$ and $`|cS+d|^2=c^2/G_s^2`$ still remains fixed. We then have
$`\widehat{G}_{\mu \nu }={\displaystyle \frac{|c|}{G_s}}\left({\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}}\right)^{\delta /2}G_{\mu \nu },\widehat{G}_s={\displaystyle \frac{c^2}{G_s}},`$
$`\widehat{\mathrm{\Theta }}^{23}=2\pi \stackrel{~}{b}G_s/c.`$ (67)
We now have $`\alpha ^{}\widehat{G}^10`$ and nonvanishing noncommutativity parameter $`\widehat{\mathrm{\Theta }}^{23}`$, therefore we end up with a NCYM. This NCYM is weakly coupled if the original NCOS is strongly coupled. Therefore the $`SL(2,Z)`$ duality is useful.
* $`b)^{}`$. The discussion for this case is basically the same as in $`a)^{}`$ above and we do not repeat them here.
### 4.2 E $`||`$ B Case
Unlike the previous one, this case is relatively simple since the open string metric is always diagonal and we do not have the same complications as we encountered there. Let us begin with the decoupling limit for NCYM.
For having sensible quantum NCYM, we need to keep the open string metric, the open string coupling and at least one space-space noncommutativity parameter fixed as $`\alpha ^{}0`$. For simplicity, we choose $`G_{\mu \nu }=\eta _{\mu \nu }=(1,1,1,1)`$. From (37), we have $`\stackrel{~}{E}^21`$. We then have the following decoupling limit:
$$\alpha ^{}0,\stackrel{~}{B}=\frac{\stackrel{~}{b}}{\alpha ^{}},g_2=\left(\frac{\stackrel{~}{b}}{\alpha ^{}}\right)^2,g_1(1\stackrel{~}{E}^2)=1,g_s=\frac{G_s}{\sqrt{1\stackrel{~}{E}^2}}\frac{\alpha ^{}}{\stackrel{~}{b}}.$$
(68)
The only nonvanishing noncommutativity parameter is
$$\mathrm{\Theta }^{23}=2\pi \stackrel{~}{b}.$$
(69)
In the above, we have not specified how $`\stackrel{~}{E}`$ scales. It appears that the resulting NCYM does not require this as long as $`\stackrel{~}{E}^21`$. However, the scaling behavior of this parameter has great impact on its $`SL(2,Z)`$ dual description. This dual description may have a small coupling, therefore a good one, in the case when the open string coupling $`G_s`$ is large. For this purpose, let us consider the following three cases which correspond to those studied in :
* a) $`\stackrel{~}{E}`$ is fixed but it equals neither 0 nor unity<sup>8</sup><sup>8</sup>8 $`\stackrel{~}{E}=0`$ corresponds to zero electric field which is not our interest here. $`\stackrel{~}{E}=1`$ gives a singular open string metric and the NCYM is no longer 1 + 3 dimensional which is not our interest here, either. So we exclude these two cases here..
* b) $`\stackrel{~}{E}=1(\alpha ^{}/\stackrel{~}{b}^{})^\delta /2`$ with $`\delta >0`$.
* c) $`\stackrel{~}{E}=(\alpha ^{}/\stackrel{~}{b}^{})^\beta `$, with $`\beta >0`$.
We would like to point out that the electric field in b) becomes critical but does not have effect on NCYM.
Let us study each of the above in the $`SL(2,Z)`$ dual description.
Case a): Using the decoupling limit in (68), we have $`|c\lambda +d|=c/g_s=c\stackrel{~}{b}(1\stackrel{~}{E}^2)^{1/2}/(\alpha ^{}G_s)`$ and $`|cS+d|=c\stackrel{~}{B}\stackrel{~}{E}/G_s=c\stackrel{~}{b}\stackrel{~}{E}/(\alpha ^{}G_s)`$. Using these we have
$$\widehat{G}_{\mu \nu }\eta _{\mu \nu }/\alpha ^{},\widehat{\mathrm{\Theta }}^{01}\alpha ^2,\widehat{\mathrm{\Theta }}^{23}\alpha ^{}\widehat{G}_s1/\alpha ^2.$$
(70)
Since $`\alpha ^{}\widehat{G}^{\mu \nu }\eta ^{\mu \nu }\alpha ^20`$, we still have a field theory but defined in a commutative geometry. However, this theory is bad since it has an infinitely large open string coupling and a singular metric. Even if we rescale the coordinates to have a finite metric but we cannot change the open string coupling. So we cannot turn this theory to a well-defined one. We here reach the same conclusion as in regardless of the fact that $`\chi `$ is rational or not.
Case b): This case is not much different from case a). Even though the scaling of the open string metric depends on whether $`\chi `$ is rational or not, it always blows up as $`\alpha ^{}0`$. So we still end up with a field theory which is not well-defined since the open string coupling blows up in the same way as in case a). The noncommutativity parameters scale as
$$\widehat{\mathrm{\Theta }}^{01}\alpha ^{2+\delta },\widehat{\mathrm{\Theta }}^{23}\alpha ^{}.$$
(71)
Case c): From our experience in on S-duality, we expect that this is the case for which we expect to have NCOS. We now have $`g_s=\alpha ^{}G_s/\stackrel{~}{b}`$ and $`|c\lambda +d|=c/g_s=c\stackrel{~}{b}/(\alpha ^{}G_s)`$. We have three sub-cases to consider: $`0<\beta <1`$, $`\beta =1`$ and $`\beta >1`$. For $`0<\beta <1`$, we reach the same conclusion as in case a) and b) above, i.e., we end up with a field theory which is not well-defined because of the infinitely large open string coupling. This subcase has also been studied in on S-dual rather than on $`SL(2,Z)`$ dual. The conclusion remains the same and we will skip the details. We now focus on $`\beta =1`$ and $`\beta >1`$ subcases. For $`\beta =1`$, $`|cS+d|^2=[(c\chi +d)+\stackrel{~}{b}/(\stackrel{~}{b}^{}G_s)]^2+c^2/G_s^2`$ is fixed and we have
$`\widehat{G}_{\mu \nu }={\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}}{\displaystyle \frac{G_s}{|c|}}|cS+d|^2\eta _{\mu \nu },\widehat{G}_s=G_s|cS+d|^2,`$
$`\widehat{\mathrm{\Theta }}^{01}={\displaystyle \frac{2\pi \stackrel{~}{b}}{G_s|cS+d|^2}},\widehat{\mathrm{\Theta }}^{23}={\displaystyle \frac{2\pi \stackrel{~}{b}}{|cS+d|^2}}\left[(d+c\chi )+{\displaystyle \frac{c\stackrel{~}{b}}{G_s\stackrel{~}{b}^{}}}\right].`$ (72)
The above implies that $`\alpha ^{}\widehat{G}^{\mu \nu }`$ is fixed. We therefore have NCOS rather than NCYM. In other words, the $`SL(2,Z)`$ dual of NCYM for $`\beta =1`$ gives a NCOS whether $`\chi `$ is rational or not. This is due to $`g_s0`$. However, whether $`\chi `$ is rational or not is important in determining the usefulness of the $`SL(2,Z)`$ duality. Our primary purpose is to find a weakly coupled theory by $`SL(2,Z)`$ duality when the open string coupling for NCYM is large. When $`\chi `$ is irrational, we map a strongly coupled theory (NCYM) to another strongly coupled theory (NCOS) by $`SL(2,Z)`$ duality which can be examined from the relation between two open string couplings given in (72). So $`SL(2,Z)`$ duality is not particularly useful in this case. However, when $`\chi `$ is rational, we can always choose $`c\chi +d=0`$ through $`SL(2,Z)`$ duality. Then we can map a strongly coupled theory (NCYM) to a physically equivalent and weakly coupled theory (NCOS). So only for rational $`\chi `$, the S-duality is useful.
For $`\beta >1`$, we continue to have $`|c\lambda +d|=|c|/g_s=|c|\stackrel{~}{b}/(\alpha ^{}G_s)`$ but now $`|cS+d|^2=(c\chi +d)^2+c^2/G_s^2`$. We then have
$`\widehat{G}_{\mu \nu }={\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}}{\displaystyle \frac{G_s}{|c|}}|cS+d|^2\eta _{\mu \nu },\widehat{G}_s=G_s|cS+d|^2,`$
$`\widehat{\mathrm{\Theta }}^{01}={\displaystyle \frac{2\pi \stackrel{~}{b}}{G_s|cS+d|^2}},\widehat{\mathrm{\Theta }}^{23}={\displaystyle \frac{2\pi (c\chi +d)}{|cS+d|^2}}.`$ (73)
We have again NCOS since $`\alpha ^{}\widehat{G}^{\mu \nu }\eta ^{\mu \nu }`$. Only for rational $`\chi `$, a strongly coupled NCYM can be mapped to a weakly coupled NCOS by $`SL(2,Z)`$ duality since we can choose $`c\chi +d=0`$. Once such a choice is made, we have $`\widehat{\mathrm{\Theta }}^{23}=0`$. This case is not different from the one with $`\stackrel{~}{E}=0`$. However, when $`\chi `$ is irrational, we end up not only with a strongly coupled theory but also with nonvanishing $`\widehat{\mathrm{\Theta }}^{23}`$ even if we start with $`\stackrel{~}{E}=0`$.
Let us now discuss the $`SL(2,Z)`$ duality of NCOS. To have NCOS, we need $`\alpha ^{}G^{\mu \nu }`$ and at least $`\mathrm{\Theta }^{01}`$ to be fixed when the critical electric field limit $`\stackrel{~}{E}1`$ is taken. Unlike in the field theory limit, we do not need to take $`\alpha ^{}0`$ since we do not require the open string massive modes to decouple from its massless ones. Our purpose here is to study the $`SL(2,Z)`$ dual of NCOS which might be a NCYM. For this reason, it is convenient to set $`\alpha ^{}0`$ for NCOS such that we can easily discuss its $`SL(2,Z)`$ dual which has the possibility of NCYM. In doing so, the effective open string $`\alpha _{\mathrm{eff}}^{}`$ for the NCOS is still fixed and is determined by the noncommutativity parameter or scale. Since we require $`\alpha ^{}G^{\mu \nu }`$ to be fixed as $`\alpha ^{}0`$, so we can set for simplicity
$$G_{\mu \nu }=\frac{\alpha ^{}}{\stackrel{~}{b}}\eta _{\mu \nu }.$$
(74)
From eqs.(34)-(37), we have the following decoupling limit:
$`\alpha ^{}0,\stackrel{~}{E}=1{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}^{}}}\right)^\delta ,g_1={\displaystyle \frac{\stackrel{~}{b}^{}}{\stackrel{~}{b}}}\left({\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}^{}}}\right)^{1\delta },`$
$`g_2(1+\stackrel{~}{B}^2)={\displaystyle \frac{a^{}}{\stackrel{~}{b}}},g_s=\left({\displaystyle \frac{\stackrel{~}{b}^{}}{\alpha ^{}}}\right)^{\delta /2}{\displaystyle \frac{G_s}{(1+\stackrel{~}{B}^2)^{1/2}}},`$ (75)
with $`\delta >0`$. From the above and (36), we have
$$\mathrm{\Theta }^{01}=2\pi \stackrel{~}{b},\mathrm{\Theta }^{23}=2\pi \stackrel{~}{b}\stackrel{~}{B}.$$
(76)
In the above, we have not yet specified how $`\stackrel{~}{B}`$ scales. In general we can set $`\stackrel{~}{B}=h(\alpha ^{}/\stackrel{~}{b}^{})^\beta `$ with $`h`$ fixed and $`\beta 0`$. For $`\beta =0`$, $`\mathrm{\Theta }^{23}`$ is fixed while it vanishes for $`\beta >0`$. Using this decoupling limit, we try to find the underlying theory after $`SL(2,Z)`$ duality. Whether $`\chi `$ is rational or not is crucial for the conclusion. So we discuss them separately in the following.
Irrational $`\chi `$: In this case we have $`c\chi +d0`$. Now $`|c\lambda +d|=|c\chi +d|`$ since $`g_s\mathrm{}`$. For $`\beta =0`$, $`|cS+d|^2=[c(\chi +h/G_s)+d]^2+c^2/G_s^2`$. Whereas for $`\beta >0`$, $`|cS+d|=(c\chi +d)^2+c^2/G_s^2`$. So for $`\beta 0`$, $`|cS+d|`$ is fixed. From (40), (41) and (42), we have
$`\widehat{G}_{\mu \nu }={\displaystyle \frac{\alpha ^{}}{\stackrel{~}{b}}}{\displaystyle \frac{|cS+d|^2}{|c\chi +d|}}\eta _{\mu \nu },\widehat{G}_s=|cS+d|G_s,`$
$`\widehat{\mathrm{\Theta }}^{01}={\displaystyle \frac{(c\chi +d)}{|cS+d|^2}}\mathrm{\Theta }^{01},\widehat{\mathrm{\Theta }}^{23}={\displaystyle \frac{(c\chi +d)\mathrm{\Theta }^{23}2\pi c\stackrel{~}{b}(1+\stackrel{~}{B}^2)/G_s}{|cS+d|^2}}.`$ (77)
The scaling of the metric $`\widehat{G}_{\mu \nu }`$ tells that we end up actually with a NCOS rather than a NCYM for irrational $`\chi `$. This has been given first in . Notice that we now have nonvanishing $`\widehat{\mathrm{\Theta }}^{23}`$ even if we begin with $`\mathrm{\Theta }^{23}=0`$. However, we map a strongly coupled NCOS to another strongly coupled NCOS. Therefore, for irrational $`\chi `$, the $`SL(2,Z)`$ duality is not that useful. The interesting point in this case is that we can use it to reduce or to increase the space-space noncommutative directions (since we can get a vanishing $`\mathrm{\Theta }^{23}`$ from a nonvanishing $`\widehat{\mathrm{\Theta }}^{23}`$ or vice-versa).
Rational $`\chi `$: Now we can always choose $`c\chi +d=0`$. Then $`|c\lambda +d|=c/g_s=(\alpha ^{}/\stackrel{~}{b}^{})^{\delta /2}c(1+\stackrel{~}{B}^2)^{1/2}/G_s`$. Again $`|cS+d|^2`$ remains fixed. So we have
$$\widehat{G}_{\mu \nu }\alpha ^{1\delta /2}\eta _{\mu \nu },\widehat{\mathrm{\Theta }}^{01}=0,\widehat{\mathrm{\Theta }}^{23}=\frac{2\pi c\stackrel{~}{b}(1+\stackrel{~}{B}^2)}{G_s|cS+d|},$$
(78)
and the open string coupling $`\widehat{G}_s=c^2(1+h^2)/G_s`$ for $`\beta =0`$ and $`\widehat{G}_s=c^2/G_s`$ for $`\beta >0`$.
Since $`\alpha ^{}\widehat{G}^1\alpha ^{\delta /2}0`$, we therefore end up with a NCYM with noncommutative space-space directions. This has also been studied in . So now a strongly coupled NCOS is physically equivalent to a weakly coupled NCYM. In this case, the $`SL(2,Z)`$ is really useful and now a $`SL(2,Z)`$ is not much different from a simple S-duality as studied in . A space-time noncommutativity is also transformed to a space-space one.
## 5 Conclusion
To conclude, we have discussed in this paper various decoupling limits for noncommutative open string/Yang-Mills theory in four-dimensions and their $`SL(2,Z)`$ duality for both $`𝐄𝐁`$ and $`𝐄||𝐁`$ cases. Since $`SL(2,Z)`$ is a non-perturbative quantum symmetry of type IIB string theory, we often use this symmetry to find a physically equivalent and yet weakly coupled theory if the original theory is strongly coupled. However, our study indicates that if the RR scalar in one theory is irrational, the $`SL(2,Z)`$ does not help much and the $`SL(2,Z)`$ dual is always strongly coupled. So when we say that we can use S-duality or in general $`SL(2,Z)`$ duality to transform a strongly coupled theory to a weakly coupled one, one must understand that this can be done only for rational $`\chi `$. Since $`\chi `$ is determined by the underlying (most likely non-perturbative) vacuum, whether $`\chi `$ is rational or not is a rather non-trivial question. We cannot answer this until we understand the non-perturbative type IIB theory completely.
We also find that $`SL(2,Z)`$ symmetry can be used to increase or decrease the number of noncommutative directions but it seems that we cannot turn an OYM to a NCYM/NCOS through this symmetry. We also find that the interplay of electric and magnetic fields are important in controlling the number of noncommutative directions. We show that the $`SL(2,Z)`$ duality of NCYM can be an ordinary theory which is not well-defined or a NCOS regardless of whether $`\chi `$ is rational or not. But only for rational $`\chi `$, the NCOS can be weakly coupled if the original NCYM is strongly coupled. Also when the original NCOS is strongly coupled the $`SL(2,Z)`$ duality is either another strongly coupled NCOS if $`\chi `$ is irrational or a weakly coupled NCYM if $`\chi `$ is rational. Some of the critical electric field limit for $`𝐄𝐁`$ are particularly interesting. Whether we have decoupled light-like NCYM or light-like NCOS or other kinds of OS or light-like OYM is still not clear. Further study is needed.
ACKNOWLEDGMENTS
JXLU acknowledges the support of U. S. Department of Energy.
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# A new method for measuring azimuthal distributions in nucleus-nucleus collisions
## I Introduction
In heavy ion collisions, much work is devoted to the study of the azimuthal distributions of outgoing particles, and in particular of distributions with respect to the reaction plane. Since these distributions reflect the interactions between particles, possible anisotropies, the so-called “flow,” reveal information on the hot stages of the collision: thermalization, pressure gradients, time evolution, etc.
Since the orientation of the reaction plane is not known a priori, flow measurements are usually extracted from two-particle azimuthal correlations. This is based on the idea that azimuthal correlations between two particles are generated by the correlation of the azimuth of each particle with the reaction plane. The assumption that this is the only source of two-particle azimuthal correlations, or at least that other sources can be neglected, dates back to the early days of the flow . It still underlies the analyses done at ultrarelativistic energies, both at the Brookhaven AGS and the CERN SPS.
However, we have shown in recent papers that other sources of azimuthal correlations (which we refer to as “nonflow” correlations) are of comparable magnitude and must be taken into account in the flow analysis. We have studied in detail the well-known correlations due to global momentum conservation and those due to quantum correlations between identical particles . We have also discussed other correlations due to resonance decays and final state interactions .
Nonflow correlations scale with the total multiplicity $`N`$ like like $`1/N`$. Thus, they become large for peripheral collisions. It is important to take them into account, in particular, when studying the centrality dependence of the flow, which has been recently proposed as a sensitive probe of the phase transition to the quark-gluon plasma .
Clearly, a reliable flow analysis should eliminate nonflow correlations. Correlations due to momentum conservation can be calculated analytically and subtracted from the measured correlations, so as to isolate the correlations due to flow ; short range correlations can be measured independently and subtracted in the same way . Other well-identified nonflow correlations can be estimated through a Monte-Carlo simulation . This method was used by the WA93 Collaboration to estimate direct correlations from $`\pi ^0\gamma \gamma `$ decays . Alternatively, one can attempt to eliminate nonflow correlations directly at the experimental level: effects of momentum conservation cancel if the detector used in the flow analysis is symmetric with respect to midrapidity ; short range correlations are eliminated if one correlates two subevents separated by a gap in rapidity. This is the method recently used by the STAR Collaboration at RHIC . In the STAR paper, the correlations between pions of the same charge are also compared with correlations between $`\pi ^+`$ and $`\pi ^{}`$: correlations from $`\rho ^0\pi ^+\pi ^{}`$ are thus found to be negligible.
Nevertheless, nonflow correlations remain which cannot be handled so simply. Correlations due to resonance decays, for instance, are hard to estimate (this would require a detailed knowledge of the collision dynamics) and there is no known systematic way to eliminate them at the experimental level; more importantly, the production of minijets will contribute to azimuthal correlations in the experiments at higher energies, at RHIC and LHC. Finally, the existence of other sources of nonflow correlations, so far unknown, cannot be excluded.
The purpose of this paper is to propose a new method for the flow analysis which requires no knowledge of nonflow correlations. The general idea is to eliminate these latter using higher order azimuthal correlations. Higher order correlations were previously used in to show qualitatively the collectivity of flow. The study presented in this paper is more quantitative : by means of a cumulant expansion, we are able to extract the value of the flow from multiparticle correlations. The method we propose is more reliable, and in many respects simpler than traditional methods . In particular, detector defects, which must be considered carefully when measuring anisotropies of a few percent, can be corrected in a compact and elegant way.
In Sec. II, we give the principle of our method as well as orders of magnitude. We show in particular that this method is more sensitive: it allows measurements of azimuthal anisotropies down to values of order $`1/M`$, instead of $`1/\sqrt{N}`$ with the standard analysis, where $`N`$ denotes the total multiplicity of particles emitted in the collision.
Then, we show how the method can be implemented practically. As usual, the measurement of azimuthal distributions is performed in two steps. First, one reconstructs approximately the orientation of the reaction plane from the directions of many emitted particles, and one estimates the statistical uncertainties associated with this reconstruction. In fact, this first step amounts to measuring the value of the flow, integrated over some region of phase space (corresponding typically to a detector). We show in Sec. III how this measurement can be done using moments of the distribution of the $`Q`$-vector, which generalizes the transverse momentum transfer introduced by Danielewicz and Odyniec in order to estimate the azimuth of the reaction plane . We also discuss an improved version of the subevent method introduced by the same authors to estimate the accuracy of the reaction plane reconstruction.
The second step in the flow analysis is to perform more detailed measurements of azimuthal distributions, for various particles, as a function of rapidity and/or transverse momentum. We refer to these detailed measurements as to “differential flow”. They are usually performed by measuring distributions with respect to the reconstructed reaction plane, and then correcting for the statistical errors in this reconstruction, which have been estimated previously. Here, the differential flow will be extracted directly from the correlation between the azimuths of the outgoing particles and the $`Q`$-vector, as explained in Sec. IV. The discussion applies so far to an ideal detector. A general way of implementing acceptance corrections adapted to our method is discussed in Sec. V. Finally, the correct procedure is summarized in Sec. VI. Readers already familiar with flow analysis and willing to apply our method may go directly to this last section.
## II Cumulant expansion of azimuthal correlations
As the standard methods of flow analysis , our method is based on a Fourier expansion of azimuthal distributions which is defined in Sec. II A. Then, in Sec. II B, we discuss two-particle azimuthal correlations, on which the standard flow analysis relies, and show that they decompose into a contribution from flow and an additional term of order $`1/N`$ which corresponds to nonflow correlations; this latter contribution limits the sensitivity of the traditional method. In Sec. II C, the decomposition is generalized to multiparticle correlations. Finally, in Sec. II D, we show that this decomposition of multiparticle correlations allows us to obtain more sensitive measurements of flow.
### A Fourier coefficients
We call “flow” the azimuthal correlations between the outgoing particles and the reaction plane. These are conveniently characterized in terms of the Fourier coefficients $`v_n`$ which we now define. In most of this paper, we shall work with a coordinate system in which the $`x`$ axis is the impact direction, and $`(x,z)`$ the reaction plane, while $`\varphi `$ denotes the azimuthal angle with respect to the reaction plane. In this frame, the momentum of a particle of mass $`m`$ is
$$𝐩=\left(\begin{array}{c}p_x=p_T\mathrm{cos}\varphi \\ p_y=p_T\mathrm{sin}\varphi \\ p_z=\sqrt{p_T^2+m^2}\mathrm{sinh}y\end{array}\right),$$
(1)
where $`p_T`$ is the transverse momentum and $`y`$ the rapidity. Since the orientation of the reaction plane is unknown in experiments, so is the azimuth $`\varphi `$. Therefore $`p_x`$ and $`p_y`$ are not measured directly.
When necessary, we shall denote by $`\overline{\varphi }`$ the azimuthal angle in the laboratory frame. Unlike $`\varphi `$, $`\overline{\varphi }`$ is a measurable quantity, related to $`\varphi `$ by $`\overline{\varphi }=\varphi +\varphi _R`$, where $`\varphi _R`$ is the unknown azimuthal angle of the reaction plane in the laboratory system.
With these definitions, $`v_n`$ can be expressed as a function of the one-particle momentum distribution $`f(𝐩)\mathrm{d}N/\mathrm{d}^3𝐩`$
$$v_n(𝒟)e^{in\varphi }=\frac{{\displaystyle _𝒟}e^{in\varphi }f(𝐩)\mathrm{d}^3𝐩}{{\displaystyle _𝒟}f(𝐩)\mathrm{d}^3𝐩},$$
(2)
where the brackets denote an average value over many events, and $`𝒟`$ represents a phase space window in the $`(p_T,y)`$ plane where flow is measured, typically corresponding to a detector. Since the particle source is symmetric with respect to the reaction plane for spherical nuclei, $`\mathrm{sin}n\varphi `$ vanishes and $`v_n`$ is real.
The purpose of the flow analysis is to extract $`v_n`$ from the data. Only the first two coefficients $`v_1`$ and $`v_2`$ have been published. They are usually called directed and elliptic flow, respectively. There are so far very few measurements of higher order coefficients. The E877 experiment at the Brookhaven AGS reported values compatible with zero for $`v_3`$ and $`v_4`$ . Nonvanishing values of higher harmonics, up to $`v_6`$, were reported from preliminary analyses at the CERN SPS . However, the latter results are likely to be strongly biased by quantum two-particle correlations . At the energies of the CERN SPS, $`v_1`$ and $`v_2`$ are of the order of a few percent , close to the limit of detectability with the standard methods, hence the need for a new, more sensitive method.
### B Two-particle correlations
Since the actual orientation of the reaction plane is not known experimentally, one can only measure relative azimuthal angles between outgoing particles. The standard flow analysis relies on the measurement of two-particle azimuthal correlations, which involve the two-particle distribution $`f(𝐩_1,𝐩_2)=\mathrm{d}N/\mathrm{d}^3𝐩_1\mathrm{d}^3𝐩_2`$:
$$e^{in(\varphi _1\varphi _2)}_{𝒟_1\times 𝒟_2}=\frac{{\displaystyle _{𝒟_1\times 𝒟_2}}e^{in(\varphi _1\varphi _2)}f(𝐩_1,𝐩_2)\mathrm{d}^3𝐩_1\mathrm{d}^3𝐩_2}{{\displaystyle _{𝒟_1\times 𝒟_2}}f(𝐩_1,𝐩_2)\mathrm{d}^3𝐩_1\mathrm{d}^3𝐩_2}.$$
(3)
The standard analysis neglects nonflow correlations. Under that assumption, the two-particle momentum distribution factorizes:
$$f(𝐩_1,𝐩_2)=f(𝐩_1)f(𝐩_2).$$
(4)
Then, Eqs. (2) and (3) give
$$e^{in(\varphi _1\varphi _2)}_{𝒟_1\times 𝒟_2}=v_n(𝒟_1)v_n(𝒟_2).$$
(5)
This equation means that the only azimuthal correlation between two particles results from their correlation with the reaction plane. Measuring the left-hand side of Eq. (5) in various phase space windows, one can then reconstruct $`v_n`$ from this equation, up to a global sign. For instance, the E877 Collaboration uses the correlations between three rapidity windows to extract flow from their data .
However, nonflow correlations do exist. The two-particle distribution can generally be written as
$$f(𝐩_1,𝐩_2)=f(𝐩_1)f(𝐩_2)+f_c(𝐩_1,𝐩_2),$$
(6)
where $`f_c(𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐})`$ denotes the correlated part of the distribution. There are various sources of such correlations, among which global momentum conservation, resonance decays (in which the decay products are correlated), final state Coulomb, strong or quantum interactions .
In the coordinate system we have chosen, where the reaction plane is fixed, $`f_c(𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐})`$ is typically of order $`1/N`$ relative to the uncorrelated part, where $`N`$ is the total number of particles emitted in the collision. This order of magnitude can easily be understood in the case of correlations between decay products, such as $`\rho \pi \pi `$. A significant fraction of the pions produced in a heavy ion collision originate from this decay, and the conservation of energy and momentum in the decay gives rise to a large correlation between the reaction products. Since a large number of $`\rho `$ mesons are produced in a high energy nucleus-nucleus collision, the probability that two arbitrary pions originate from the same $`\rho `$ is of order $`1/N`$. This $`1/N`$ scaling also holds for the correlation due to global momentum conservation .
Inserting Eq. (6) in expression (3), one finds, instead of Eq. (5),
$$e^{in(\varphi _1\varphi _2)}_{𝒟_1\times 𝒟_2}=v_n(𝒟_1)v_n(𝒟_2)+e^{in(\varphi _1\varphi _2)}_c.$$
(7)
The left-hand side represents the measured two-particle azimuthal correlation. The first term in the right-hand side is the contribution of flow to this correlation, while the second term $`e^{in(\varphi _1\varphi _2)}_c`$ denotes the contribution of the correlated part $`f_c`$. The latter term corresponds to azimuthal correlations which do not arise from flow: we call them “direct” correlations, in opposition to the indirect correlations arising from the correlation with the reaction plane, that is, from flow.
Since the correlated two-particle distribution $`f_c(𝐩_1,𝐩_2)`$ is of order $`1/N`$, so is the second term in the right-hand side of Eq. (7), which therefore reads
$$e^{in(\varphi _1\varphi _2)}_{𝒟_1\times 𝒟_2}=v_n(𝒟_1)v_n(𝒟_2)+O\left(\frac{1}{N}\right).$$
(8)
However, one must be careful with this order of magnitude. Strictly speaking, it holds only when momenta are averaged over a large region of phase space. In the case of the short range correlations due to final state interactions (Coulomb, strong, quantum) the correlations vanish as soon as the phase spaces $`𝒟_1`$ and $`𝒟_2`$ of the two particles are widely separated. This is the method used in to get rid of such correlations. If, on the other hand, $`𝒟_1`$ and $`𝒟_2`$ coincide, the short range correlations are larger than expected from Eq. (8): in this equation, the total number of emitted particles $`N`$ should be replaced by the number of particles $`M`$ used in the flow analysis, which is smaller in practice. Furthermore, in the case of correlations due to the quantum (HBT) effect, the nonflow correlation scales like $`1/N`$ only if the source radius $`R`$ scales like $`N^{1/3}`$ . From now on, we shall omit the subscript $`𝒟`$ for sake of brevity. Note, however, that all the averages we shall consider are over a region of phase space which is not necessarily the whole space, but may be restricted to the $`(p_T,y)`$ acceptance of a detector. This will be especially important in Sec. V, when we discuss acceptance corrections.
Equation (8) shows that nonflow correlations can be neglected if $`v_nN^{1/2}`$. At SPS energies, the flow is weak and this condition is not fulfilled. Indeed, we have shown that the values of flow measured by the NA49 Collaboration at CERN are considerably modified once nonflow correlations are taken into account.
### C Multiparticle correlations and the cumulant expansion
The failure of the standard analysis is due to the impossibility to separate the correlated part from the uncorrelated part in Eq. (6) at the level of two-particle correlations. The main idea of this paper is to perform this separation using multiparticle correlations. The decomposition of the particle distribution into correlated and uncorrelated parts in Eq. (6) can be generalized to an arbitrary number of particles. For instance, the three-particle distribution can be decomposed as
$`{\displaystyle \frac{\mathrm{d}N}{\mathrm{d}𝐩_1\mathrm{d}𝐩_2\mathrm{d}𝐩_3}}f(𝐩_1,𝐩_2,𝐩_3)`$ $`=`$ $`f_c(𝐩_1)f_c(𝐩_2)f_c(𝐩_3)`$ (9)
$`+`$ $`f_c(𝐩_1,𝐩_2)f_c(𝐩_3)+f_c(𝐩_1,𝐩_3)f_c(𝐩_2)+f_c(𝐩_2,𝐩_3)f_c(𝐩_1)`$ (10)
$`+`$ $`f_c(𝐩_1,𝐩_2,𝐩_3),`$ (11)
where $`f_c(𝐩_1)f(𝐩_1)`$. The last term $`f_c(𝐩_1,𝐩_2,𝐩_3)`$ corresponds to the genuine three-particle correlation, which is of order $`1/N^2`$.
To understand this order of magnitude, let us take a simple example: the $`\omega `$ meson decays mostly into three pions. First of all, this decay generates direct two-particle correlations: the relative momentum between any two of the outgoing pions is restricted by energy and momentum conservation. The corresponding correlation is of order $`1/N`$ as discussed previously in the case of $`\rho \pi \pi `$ decays. It corresponds to the second, third and fourth term in the right-hand side of Eq. (9). As stated above, the last term in this equation stands for the direct three-particle correlation. The corresponding correlation between the decay products of a given $`\omega `$ is of order unity, while the probability that three arbitrary pions come from the same $`\omega `$ scales with $`N`$ like $`1/N^2`$. Thus the correlation between three random pions is of order $`1/N^2`$.
More generally, the decomposition of the $`k`$-particle distribution yields a correlated part $`f_c(𝐩_1,\mathrm{},𝐩_k)`$ of order $`1/N^{k1}`$. Generalizing the above discussion of $`\omega \pi \pi \pi `$ decay, the decay of a cluster of $`k`$ particles will generate correlations $`f_c(𝐩_1,\mathrm{},𝐩_k^{})`$ with $`k^{}k`$. For instance, momentum conservation, which is an effect involving all $`N`$ particles emitted in a collision, produces direct $`k`$-particle correlations for arbitrary $`k`$.
Such a decomposition is similar to the cluster expansion which is well known in the theory of imperfect gases . In the language of probability theory, this is known as the cumulant expansion . Equations (6) and (9) can be represented diagrammatically by Figs. 1 and 2. In these figures, correlated distributions $`f_c`$ are represented by enclosed sets of points, i.e. they correspond to connected diagrams.
More generally, in order to decompose the $`k`$-point function $`f(𝐩_1,\mathrm{},𝐩_k)`$, one first takes all possible partitions of the set of points $`\{𝐩_1,\mathrm{},𝐩_k\}`$. To each subset of points $`\{𝐩_{i}^{}{}_{1}{}^{},\mathrm{},𝐩_{i}^{}{}_{m}{}^{}\}`$, one associates the corresponding correlated function $`f_c(𝐩_{i}^{}{}_{1}{}^{},\mathrm{},𝐩_{i}^{}{}_{m}{}^{})`$. The contribution of a given partition is the product of the contributions of each subset. Finally, $`f(𝐩_1,\mathrm{},𝐩_k)`$ is the sum of the contributions of all partitions.
The equations expressing the $`k`$-point functions $`f`$ in terms of the correlated functions $`f_c`$ can be inverted order by order, so as to isolate the term of smallest magnitude:
$`f_c(𝐩_1)`$ $`=`$ $`f(𝐩_1)`$ (12)
$`f_c(𝐩_1,𝐩_2)`$ $`=`$ $`f(𝐩_1,𝐩_2)f(𝐩_1)f(𝐩_2)`$ (13)
$`f_c(𝐩_1,𝐩_2,𝐩_3)`$ $`=`$ $`f(𝐩_1,𝐩_2,𝐩_3)f(𝐩_1,𝐩_2)f(𝐩_3)f(𝐩_1,𝐩_3)f(𝐩_2)f(𝐩_2,𝐩_3)f(𝐩_1)+2f(𝐩_1)f(𝐩_2)f(𝐩_3).`$ (14)
The cumulant expansion has been used previously in high energy physics to characterize multiparticle correlations: it has been applied to correlations in rapidity and to Bose-Einstein quantum correlations . In these studies, the interest was mainly in short range correlations. The use of higher order cumulants was therefore limited by statistics: the probability that three or more particles are very close in phase space is small. In this paper, we are interested in collective flow, which by definition produces a long range correlation, so that the limitation due to statistics is not so drastic. It will indeed be shown in Sec. III D that cumulants up to order 6 can be measured, depending on the event multiplicity and available statistics.
We shall deal with multiparticle azimuthal correlations, which generalize the two-particle azimuthal correlations in Eq. (7), and can be decomposed in the same way. Referring to the diagrammatic representation in Figs. 1 and 2, we shall name the contribution of $`f_c(𝐩_1,\mathrm{},𝐩_k)`$ to an azimuthal correlation, i.e. the genuine $`k`$-particle correlation, the “connected part” of the correlation or, equivalently, the “direct” $`k`$-particle correlation.
### D Measuring flow with multiparticle azimuthal correlations
Our method, which we now explain, allows the detection of small deviations from an isotropic distribution. If the source is isotropic, there is no flow, and the orientation of the reaction plane does not influence the particle distribution. We can therefore consider that the reaction plane has a fixed direction in the laboratory coordinate system, so that the cumulant expansion can be performed in that frame: in other terms, we replace $`\varphi `$ by the measured azimuthal angle $`\overline{\varphi }`$. One then measures the $`k^{\mathrm{th}}`$ cumulant of the multiparticle azimuthal correlation, which is of order $`N^{1k}`$ if the distribution is isotropic. Flow will appear as a deviation from this expected behaviour.
Let us be more explicit. We are dealing with azimuthal correlations. When the source is isotropic, that is, if the $`k`$-particle distribution remains unchanged when all azimuthal angles are shifted by the same quantity $`\alpha `$, the flow coefficients (2) obviously vanish. Therefore, the two-particle azimuthal correlation (7) reduces to its connected part, of order $`1/N`$. As a further consequence of isotropy, averages like $`e^{in(\varphi _1+\varphi _2\varphi _3)}`$ vanish: only $`2k`$-particle azimuthal correlations involving $`k`$ powers of $`e^{in\varphi }`$ and $`k`$ powers of $`e^{in\varphi }`$ are nonvanishing. For instance, the four-particle correlation $`e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}`$ is a priori non vanishing. Introducing the cumulant expansion defined in Sec. II C, this correlation can be decomposed into
$`e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}`$ $`=`$ $`e^{in(\varphi _1\varphi _3)}_ce^{in(\varphi _2\varphi _4)}_c+e^{in(\varphi _1\varphi _4)}_ce^{in(\varphi _2\varphi _3)}_c+e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}_c`$ (15)
$`=`$ $`e^{in(\varphi _1\varphi _3)}e^{in(\varphi _2\varphi _4)}+e^{in(\varphi _1\varphi _4)}e^{in(\varphi _2\varphi _3)}+e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}_c.`$ (16)
Note that most terms in the cumulant expansion disappear as a consequence of isotropy. The first two terms in the right-hand side of Eq. (15) are products of direct two-particle correlations, and are therefore of order $`1/N^2`$, while the last term, which corresponds to the direct four-particle correlation, is much smaller, of order $`1/N^3`$. However, in the case of short range correlations, it may rather be of order $`1/M^3`$, where $`M`$ is the number of particles used in the flow analysis, for the same reasons as discussed in Sec. II B.
We name this latter term the “cumulant” to order 4 and denote it by $`e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}`$. Using Eq. (15), it can be expressed as a function of the measured two- and four-particle azimuthal correlations:
$$e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}e^{in(\varphi _1\varphi _3)}e^{in(\varphi _2\varphi _4)}e^{in(\varphi _1\varphi _4)}e^{in(\varphi _2\varphi _3)}.$$
(17)
The reason why we introduce a new notation here is that the cumulant to order 4 will always be defined by (17) in this paper, even when the source is not isotropic. Now, if the source is not isotropic, the decomposition of the four-particle azimuthal correlation involves many terms which have been omitted in Eq. (15) (see Appendix A 1), so that the cumulant $`e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}`$ no longer corresponds to the connected part $`e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}_c`$.
In the isotropic case, the cumulant $`e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}`$ involves only direct four-particle correlations: the two-particle correlations have been eliminated in the subtraction. In order to illustrate this statement, let us consider two decays $`\rho \pi \pi `$, and “turn off” all other sources of azimuthal correlations. We label 1 and 2 the pions emitted by the first resonance, 3 and 4 the pions emitted by the second. There are correlations between $`\pi _1`$ and $`\pi _2`$, or between $`\pi _3`$ and $`\pi _4`$, so that the measured four-particle correlation, i.e. the left-hand side of Eq. (15), does not vanish. However, there is no direct four-particle correlation between the four outgoing pions, so that the cumulant (17) vanishes. More generally, if particles are produced in clusters of $`k`$ particles, there are measured azimuthal correlations to all orders, but the cumulants to order $`k^{}>k`$ vanish.
Let us now consider small deviations from isotropy, i.e. weak flow. The two-particle azimuthal correlation receives a contribution $`v_n^2`$ according to Eq. (7). For similar reasons, the four-particle correlation gets a contribution $`v_n^4`$. The cumulant defined by Eq. (17) thus becomes (see Appendix A 1)
$$e^{in(\varphi _1+\varphi _2\varphi _3\varphi _4)}=v_n^4+O\left(\frac{1}{N^3}+\frac{v_{2n}^2}{N^2}\right)$$
(18)
where the coefficient $`1`$ in front of $`v_n^4`$ is found by replacing each factor $`e^{in\varphi }`$ or $`e^{in\varphi }`$ in the left-hand side with its average value $`v_n`$. The flow $`v_n`$ can thus be obtained, up to a sign, from the measured two- and four-particle azimuthal correlations, with a better accuracy than when using only two-particle correlations, as we shall see shortly.
It should be noticed that the cumulant involves a contribution from the higher order harmonic $`2n`$, of magnitude $`v_{2n}^2/N^2`$. This contribution does not interfere with the measurement of $`v_n`$ provided the following condition is satisfied:
$$|v_{2n}|Nv_n^2.$$
(19)
Since $`v_n`$ is measurable only if $`v_n1/N`$, as we shall see later in this section, the interference with the harmonic $`2n`$ occurs only if $`|v_{2n}||v_n|`$. In practice, the only situation where this might be a problem is when measuring the directed flow $`v_1`$ at ultrarelativistic energies, where elliptic flow $`v_2`$ is expected to be larger than $`v_1`$. On the other hand, this interference will not endanger the measurement of $`v_2`$, since $`v_4`$ should be much smaller.
In the following, we shall always assume that condition (19) is fulfilled. Then, using Eq. (18), it becomes possible to measure the flow $`v_n`$ as soon as it is much larger than $`N^{3/4}`$. The sensitivity is better than with the traditional methods using two-particle correlations which, as we have seen, require $`v_nN^{1/2}`$.
Similarly, using $`2k`$-particle azimuthal correlations and taking the cumulant, i.e. isolating the connected part (which amounts to getting rid of nonflow correlations of orders less than $`2k`$), one obtains a quantity which is of magnitude $`N^{12k}`$ for an isotropic source. Flow gives a contribution of magnitude $`v_n^{2k}`$. The contribution of higher order harmonics $`v_{kn}`$ can be neglected as soon as
$$|v_{kn}|N^{k1}v_n^k.$$
(20)
If $`|v_n|1/N`$, this is not a problem, unless $`|v_{kn}||v_n|`$. This is unlikely to occur, since one expects $`v_n`$ to decrease rapidly with $`n`$. Neglecting higher order harmonics, there remains the contributions of flow, of magnitude $`v_n^{2k}`$, and of direct $`2k`$-particle correlations, of magnitude $`N^{12k}`$. Therefore, $`2k`$-particle azimuthal correlations allow measurements of $`v_n`$ if it is larger than $`N^{1+1/2k}`$. Since $`k`$ is arbitrarily large, one can ideally measure $`v_n`$ down to values of order $`1/N`$, instead of $`1/\sqrt{N}`$ with the standard methods. A necessary condition for the flow analysis is therefore
$$v_n\frac{1}{N}$$
(21)
which will be assumed throughout this paper. As we shall see in Sec. III D, the sensitivity is in fact limited experimentally by statistical errors due to the finite number of events.
In practice, the cumulants of multiparticle azimuthal correlations will be extracted from moments of the distribution of the $`Q_n`$-vector introduced in next section.
## III Integrated flow
In this section, we show how it is possible to measure the value of $`v_n`$ integrated over a phase space region. This measurement will serve as a reference when we perform more detailed measurements of azimuthal anisotropies, in Sec. IV. We first define in Sec. III A a simple version of the $`Q_n`$-vector, or event flow vector, which is used in the standard flow analysis to estimate the orientation of the reaction plane. We then show, in Sec. III B, that the integrated value of the flow can be obtained from the moments of the $`Q_n`$ distribution: eliminating nonflow correlations up to order $`2k`$ by means of a cumulant expansion, we obtain an accuracy on the integrated $`v_n`$ of magnitude $`N^{1+1/2k}`$, better than the accuracy of standard methods if $`k>1`$. Instead of using a single event vector $`Q_n`$, one can do a similar analysis using subevents (Sec. III C). Since the order $`2k`$ of the calculation is arbitrary, we obtain with either method an infinite set of equations to determine $`v_n`$. The order $`2k`$ which should be chosen when analyzing experimental data depends on the number of events available (Sec. III D). More general forms of the $`Q_n`$-vector, which allow an optimal flow analysis, are discussed in Sec. III E. Finally, in Sec. III F, we recover, as a limiting case, the results obtained in the limit of large multiplicity where the distribution of $`Q_n`$ is Gaussian . In the whole section, we assume the analysis is performed using a perfectly isotropic detector; corrections to this ideal case will be dealt with in Sec. V.
### A The $`Q`$-vector
#### 1 Definition
Consider a collision in which $`M`$ particles are detected with azimuthal angles $`\varphi _1,\mathrm{},\varphi _M`$. In order to detect possible anisotropies of the $`\varphi `$ distribution, it is natural to construct an observable which involves all the $`\varphi _j`$, i.e. a global quantity. For the study of the $`n^{\mathrm{th}}`$ harmonic, one uses the $`n^{\mathrm{th}}`$ transverse event flow vector , which we write as a complex number
$$Q_n=\frac{1}{\sqrt{M}}\underset{j=1}{\overset{M}{}}e^{in\varphi _j},$$
(22)
where $`\varphi _j`$ denotes the azimuthal angle of the $`j^{\mathrm{th}}`$ particle with respect to the reaction plane.
For simplicity, we have associated a unit weight with each particle in Eq. (22). The generalization of our results to arbitrary weights is straightforward and will be given in Sec. III E. The $`Q_n`$-vector generalizes to arbitrary harmonics the transverse momentum transfer introduced by Danielewicz and Odyniec , which corresponds to $`n=1`$ and the transverse sphericity tensor introduced in , which corresponds to the case $`n=2`$.
In practice, the number of particles $`M`$ used for the flow analysis is not equal to the total multiplicity $`N`$ of particles produced in the collision, since all particles are not detected. However, $`M`$ should be taken as large as possible. In this paper, we shall assume that $`M`$ and $`N`$ are of the same order of magnitude. The factor $`1/\sqrt{M}`$ in front of Eq. (22), which does not appear in previous definitions of the flow vector , will be explained in Sec. III A 3.
#### 2 Flow versus nonflow contributions
A nonvanishing value for the average value of the flow vector, $`Q_n`$, signals collective flow. Indeed, using Eqs. (2) and (22), it is related to the Fourier coefficient $`v_n=e^{in\varphi }`$ by
$$Q_n=\sqrt{M}v_n.$$
(23)
Note that $`Q_n`$ is real, as is $`v_n`$, due to the symmetry with respect to the reaction plane.
As stated before, the purpose of the flow analysis is to measure $`v_n`$, i.e. $`Q_n`$. This is not a trivial task because the azimuth of the reaction plane is unknown, so that the phase of $`Q_n`$ is unknown. The only measurable quantity is $`|Q_n|`$, the length of $`Q_n`$. Its square $`Q_nQ_n^{}`$, where $`Q_n^{}`$ denotes the complex conjugate, only depends on relative azimuthal angles:
$$Q_nQ_n^{}=\frac{1}{M}\underset{j,k=1}{\overset{M}{}}e^{in(\varphi _j\varphi _k)}.$$
(24)
In Sec. III B, we shall see that the flow can be deduced from the moments of the distribution of $`|Q_n|^2`$, i.e. from the average values $`|Q_n|^{2k}`$, where $`k`$ is a positive integer. To illustrate how flow enters these expressions, we discuss here the second order moment $`|Q_n|^2`$. Averaging Eq. (24) over many events and using Eq. (7), one obtains
$$|Q_n|^2=\frac{1}{M}\left[M+M(M1)\left(v_n^2+e^{in(\varphi _j\varphi _k)}_c\right)\right].$$
(25)
The first term corresponds to the diagonal terms $`j=k`$, i.e. to “autocorrelations”. If there are no azimuthal correlations (neither flow nor nonflow), only this term remains and the average value of $`|Q_n|^2`$ is exactly 1. The second term corresponds to $`jk`$, i.e. to the two-particle azimuthal correlations discussed in Sec. II B. Since $`e^{in(\varphi _1\varphi _2)}_c`$ is at most of order $`1/M`$, direct correlations give a contribution which is a priori of the same order of magnitude as autocorrelations, although it may be smaller in practice. Equation (25) can thus be written
$`|Q_n|^2`$ $`=`$ $`M\left[v_n^2+{\displaystyle \frac{1}{M}}+O\left({\displaystyle \frac{1}{M}}\right)\right]`$ (26)
$`=`$ $`Q_n^2+1+O(1).`$ (27)
As expected from the discussion of Sec. II B, since $`|Q_n|^2`$ involves two-particle correlations, flow measurements based on $`|Q_n|^2`$ are reliable only if $`|v_n|1/\sqrt{M}`$. Smaller values of flow can be obtained using higher moments of the distribution of $`|Q_n|^2`$, as explained in Sec. III B.
If flow is strong enough, the event flow vector can be used to estimate the orientation of the reaction plane. Indeed, if $`|v_n|1/\sqrt{M}`$, Eqs. (23) and (27) show that $`Q_nQ_n=\sqrt{M}v_n`$. Then the phase of $`Q_n`$ is approximately $`0`$ if $`v_n>0`$ and $`\pi `$ if $`v_n<0`$. Experimentally, one defines $`Q_n`$ as in (22), with the azimuthal angles $`\varphi _j`$ measured with respect to a fixed direction in the laboratory (rather than the reaction plane, which is unknown). Then the azimuthal angle of the reaction plane $`\varphi _R`$ can be estimated from the phase of $`Q_n`$, which we write $`n\varphi _Q`$: $`\varphi _R\varphi _Q`$ (resp. $`\varphi _R\varphi _Q+\pi /n`$) modulo $`2\pi /n`$ if $`v_n>0`$ (resp. $`v_n<0`$).
#### 3 Varying the centrality
Let us now explain the factor $`1/\sqrt{M}`$ in the definition (22). This factor was introduced independently by A. Poskanzer and S. Voloshin, and in . It is important when using events with different multiplicities $`M`$ in the flow analysis, i.e. events with different centralities. This is the case in practice: one takes all events in a given centrality interval in order to increase the available statistics.
If there is no flow, Eq. (27) shows that $`|Q_n|^2`$ is independent of $`M`$ since nonflow correlations scale like $`1/M`$. This can be understood simply: the sum in Eq. (22) is a random walk of $`M`$ unit steps, therefore it has a length of order $`\sqrt{M}`$, which cancels out with the factor $`1/\sqrt{M}`$ in front. Flow, on the other hand, depends strongly on centrality (it vanishes for central and very peripheral collisions): according to Eq. (27), it gives a positive contribution to $`|Q_n|^2`$ which strongly depends on $`M`$. This allows to disentangle flow and nonflow effects.
Note that flow can be detected by studying the variation of $`|Q_n|^2`$ with centrality. This is the method used in : one expects $`|Q_n|^2`$ to be minimum for the most peripheral collisions where the density of particles is too small for collective behaviour to set in, and for central collisions where $`v_n`$ also vanishes from azimuthal symmetry. However, such a method does not allow an accurate measurement of flow: it is impossible to select true (i.e. with $`b=0`$) central collisions experimentally, and there may still be some flow up to large impact parameters, as suggested by hydrodynamic calculations in the case of elliptic flow , and by recent measurements .
The method presented in this paper is more powerful in the sense that it allows flow measurements for a given centrality. The error on the centrality selection (due to the fact that one always selects events within a finite range of impact parameters) is compensated by the factor $`1/\sqrt{M}`$ in the definition of $`Q_n`$.
### B Cumulants of the distribution of $`|Q_n|^2`$
For sake of brevity, we now drop the subscript $`n`$ and set $`n=1`$ until the end of this paper, unless otherwide stated. All our results can be easily generalized to the study of higher order $`v_n`$’s by multiplying all azimuthal angles by $`n`$.
The moments of the $`|Q|^2`$ distribution involve the multiparticle azimuthal correlations discussed in Sec. II D. While $`|Q|^2`$ involves two-particle azimuthal correlations, as seen in Eq. (24), the higher moments $`|Q|^{2k}`$ involve $`2k`$-particle correlations. For instance, we have
$$|Q|^4=\frac{1}{M^2}\underset{j,k,l,m}{}e^{i(\varphi _j+\varphi _k\varphi _l\varphi _m)}.$$
(28)
These higher order azimuthal correlations can be used to eliminate nonflow correlations order by order, as explained in Sec. II D. This will be achieved by taking the cumulants of the distribution of $`|Q|^2`$, which we shall soon define.
#### 1 Isotropic source
Following the procedure outlined in Sec. II D, we first consider an isotropic source (no flow). Using Eq. (27), $`|Q|^2`$ is then of order unity, and so are the higher order moments $`|Q|^{2k}`$. However, by analogy with the cumulant decomposition of multiparticle distributions introduced in Sec. II C, we can construct specific combinations of the moments, namely the cumulants of the $`Q`$ distribution, which are much smaller than unity: the cumulant $`|Q|^{2k}`$ to order $`k`$, built with the $`|Q|^{2j}`$ where $`jk`$, is of magnitude $`1/M^{k1}`$.
As an illustration, let us construct the fourth order cumulant $`|Q|^4`$. If the multiplicity $`M`$ is large, most of the terms in Eq. (28) are nondiagonal, i.e. they correspond to values of $`j`$, $`k`$, $`l`$ and $`m`$ all different. Then, using the cumulant of the four-particle azimuthal correlation defined by Eq. (17) and summing over $`(j,k,l,m)`$, it is natural to define $`|Q|^4`$ as
$$|Q|^4=|Q|^42|Q|^2^2.$$
(29)
The order of magnitude of $`|Q|^4`$ is easy to derive: each term of type (17) is of order $`1/M^3`$ as discussed in Sec. II D; there are $`M^4`$ such terms in the sum (28); taking into account the factor $`1/M^2`$ in front of the sum, $`|Q|^4`$ is finally of order $`1/M`$. As intended, two-particle nonflow correlations, which are of order unity, have been eliminated in the subtraction (29).
A more careful analysis must take into account diagonal terms for which two (or more) indices among $`(j,k,l,m)`$ are equal. This analysis is presented in Appendix A 2, where we show that diagonal terms are also of order $`1/M`$: they give a contribution of the same order of magnitude as direct four-particle correlations. In the following, we shall assume that this property, namely that the contribution of diagonal terms is at most of the magnitude of the contribution of nondiagonal terms, also holds for higher order moments.
Among these diagonal terms are the autocorrelations already encountered in the expansion of $`|Q|^2`$ \[see the discussion below Eq. (25)\], which we define as the terms which remain in the absence of flow and direct correlations. A straightforward calculation (see Appendix A 2) shows that their contribution to the cumulant $`|Q|^4`$ is $`1/M`$. As in the case of the second order moment $`|Q|^2`$ discussed previously, autocorrelations are a priori of the same order of magnitude as other nonflow correlations. As we shall see later in this section, they can easily be calculated and removed order by order.
Arbitrary moments $`|Q|^{2k}`$ can be decomposed into cumulants, which can then be isolated in a similar way. This decomposition can be represented in terms of diagrams, like the decomposition of the multiparticle distribution in Sec. II D. This is explained in detail in Appendix B. For example, the decomposition of $`|Q|^4`$ is displayed in Fig. 3. In these diagrams, each dot on the left (resp. on the right) of the dashed line represents a power of $`Q`$ (resp. $`Q^{}`$), and correlated parts, which correspond to direct correlations, are circled: the equation displayed in Fig. 3 stands for
$$|Q|^4=2|Q|^2^2+|Q|^4.$$
(30)
Since $`|Q|^2=|Q|^2`$, one recovers Eq. (29). More generally, to decompose $`|Q|^{2k}`$, one draws $`k`$ dots on each side of the dashed line. The diagrams combine all possible subsets of the dots on the left with subsets of the dots on the right containing the same number of elements. The latter condition is due to the fact that the average value of $`Q^lQ^m`$ vanishes when $`lm`$, as a consequence of isotropy.
In order to invert these relations, and to express the cumulants as a function of the measured moments, the simplest way consists in using the formalism of generating functions, recalled in Appendix B 2. There, it is shown that the cumulant $`|Q|^{2k}`$ is obtained from the expansion in power series of $`x`$ of the following generating equation, and then the identification of the coefficients of $`x^{2k}`$:
$`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x^{2k}}{(k!)^2}}|Q|^{2k}`$ $`=`$ $`\mathrm{ln}\left({\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x^{2k}}{(k!)^2}}|Q|^{2k}\right)`$ (31)
$`=`$ $`\mathrm{ln}I_0(2x|Q|),`$ (32)
where $`I_0`$ is the modified Bessel function of order 0. Expanding this equation to order $`x^4`$, one recovers Eq. (29); to order $`x^6`$, one obtains the sixth order cumulant
$$|Q|^6=|Q|^69|Q|^4|Q|^2+12|Q|^2^3,$$
(33)
which is of the order of $`1/M^2`$ for an isotropic source.
#### 2 Contribution of flow
Let us now consider small deviations from isotropy. As explained in Sec. II D, these deviations will contribute to the cumulants $`|Q|^{2k}`$ defined above.
The contribution of flow to the fourth order cumulant $`|Q|^4`$ is calculated in detail in Appendix A. It is shown in particular that the diagonal terms in Eq. (28) are at most of the same magnitude as nondiagonal terms, as in the case of an isotropic source. As in Sec. II D, higher order harmonics can be neglected as soon as condition (19) is fulfilled. One then obtains
$$|Q|^4=Q^4\frac{1}{M}+O\left(\frac{1}{M}\right)$$
(34)
where the term $`1/M`$ is the contribution of autocorrelations, i.e. the case $`j=k=l=m`$. From Eqs. (23) and (34), one can measure values of the integrated flow $`v`$ down to $`M^{3/4}`$, instead of $`M^{1/2}`$ with traditional methods.
Increased sensitivity can be attained using higher order cumulants. As shown in Appendix B 3, the cumulants defined by Eq. (32) are related to the flow by the following generating equation:
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{x^{2k}}{(k!)^2}|Q|^{2k}=\mathrm{ln}I_0(2xQ)+M\mathrm{ln}I_0\left(\frac{2x}{\sqrt{M}}\right).$$
(35)
Expanding this equation up to order $`x^{2k}`$, and isolating the coefficient of $`x^{2k}/(k!)^2`$, one obtains a relation with on the left-hand side the cumulant $`|Q|^{2k}`$, while the first term on the right-hand side is the contribution of flow, and the second term corresponds to autocorrelations. This identity holds within an error of order $`M^{1k}`$ due to direct $`2k`$-particle correlations. Using Eq. (23), it therefore allows measurements of $`v`$ within $`O(M^{1+1/2k})`$, as expected from the discussion of Sec. II D. Expanding Eq. (35) to order $`x^4`$, one recovers (34). To order $`x^6`$, one obtains
$$|Q|^6=4Q^6+\frac{4}{M^2}+O\left(\frac{1}{M^2}\right)$$
(36)
which extends the limit of detectability down to $`vM^{5/6}`$.
Since Eqs. (32) and (35) can be expanded to any order, one obtains an infinite set of equations to determine the same quantity $`Q`$. The best choice for the order $`k`$ will be discussed below in Sec. III D. Before we come to this point, we shall discuss an alternative method to measure $`Q`$, the so-called “subevent” method.
### C Subevents
The standard flow analysis, instead of studying the autocorrelation of the event flow vector as in Sec. III B, deals with “subevents”: the set of detected particles is divided randomly into two subsets $`\mathrm{I}`$ and $`\mathrm{I}\mathrm{I}`$ of equal multiplicities, and the two corresponding (subevent) flow vectors $`Q_\mathrm{I}`$ and $`Q_{\mathrm{I}\mathrm{I}}`$ are constructed. Then one studies the azimuthal correlation between $`Q_\mathrm{I}`$ and $`Q_{\mathrm{I}\mathrm{I}}`$ . This is usually done under the assumption that the only azimuthal correlation between the subevents is due to flow. Then, from the flow of two equivalent subevents, one can deduce the flow of the whole event by a simple multiplication by a factor of $`\sqrt{2}`$, as will soon be explained.
A nice feature of that method is that, since the subevents have no particle in common, autocorrelations are automatically removed: only correlations due to flow and direct correlations remain. Therefore, one may prefer to work with subevents when direct correlations are small (although they are, generally, of the same order of magnitude as autocorrelations).
In this section, we shall improve the standard subevent method, in the spirit of Sec. III B: we eliminate nonflow azimuthal correlations order by order by means of a cumulant expansion of the distribution of $`Q_\mathrm{I}`$ and $`Q_{\mathrm{I}\mathrm{I}}`$, thereby increasing the sensitivity of the method.
#### 1 Definitions
Consider two separate subevents of multiplicity $`M_\mathrm{I}`$ and $`M_{\mathrm{I}\mathrm{I}}`$ respectively (in practice, one chooses $`M_\mathrm{I}=M_{\mathrm{I}\mathrm{I}}`$). We can construct the subevent flow vectors $`Q_{\mathrm{I}}^{}{}_{n}{}^{}`$ and $`Q_{\mathrm{I}\mathrm{I}}^{}{}_{n}{}^{}`$ as following:
$`Q_{\mathrm{I}}^{}{}_{n}{}^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{M_\mathrm{I}}}}{\displaystyle \underset{j=1}{\overset{M_\mathrm{I}}{}}}e^{in\varphi _j}=|Q_{\mathrm{I}}^{}{}_{n}{}^{}|e^{in\mathrm{\Psi }_\mathrm{I}},`$ (37)
$`Q_{\mathrm{I}\mathrm{I}}^{}{}_{n}{}^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{M_{\mathrm{I}\mathrm{I}}}}}{\displaystyle \underset{k=1}{\overset{M_{\mathrm{I}\mathrm{I}}}{}}}e^{in\stackrel{~}{\varphi }_k}=|Q_{\mathrm{I}\mathrm{I}}^{}{}_{n}{}^{}|e^{in\mathrm{\Psi }_{\mathrm{I}\mathrm{I}}}.`$ (38)
As in Sec. III B, we set $`n=1`$ and drop the subscript $`n`$ in $`Q_{\mathrm{I}}^{}{}_{n}{}^{}`$ and $`Q_{\mathrm{I}\mathrm{I}}^{}{}_{n}{}^{}`$; generalization to higher $`n`$ is straightforward.
By analogy with Eq. (23), we write
$$Q_\mathrm{I}=\sqrt{M_\mathrm{I}}v_\mathrm{I},Q_{\mathrm{I}\mathrm{I}}=\sqrt{M_{\mathrm{I}\mathrm{I}}}v_{\mathrm{I}\mathrm{I}},$$
(39)
where $`v_\mathrm{I}`$ and $`v_{\mathrm{I}\mathrm{I}}`$ denote the values of $`v_1`$ associated with each subevent. Hereafter, we shall assume that the two subevents are equivalent, i.e. $`v_\mathrm{I}=v_{\mathrm{I}\mathrm{I}}v`$, as is the case if they are chosen randomly.
Therefore, if $`M_\mathrm{I}=M_{\mathrm{I}\mathrm{I}}=M/2`$, the value of $`Q`$ for the whole event is related to the value for one subevent, $`Q_\mathrm{I}`$, by
$$Q=\sqrt{M}v=\sqrt{2}Q_\mathrm{I}.$$
(40)
The purpose here is to measure $`Q_\mathrm{I}`$, which is equivalent to measuring $`Q`$.
#### 2 Limitations of the standard method
In order to extract the azimuthal correlation between the subevents, the simplest possibility is to form the product
$$Q_\mathrm{I}Q_{\mathrm{I}\mathrm{I}}^{}=\frac{1}{\sqrt{M_\mathrm{I}M_{\mathrm{I}\mathrm{I}}}}\underset{j,k}{}e^{i(\varphi _j\stackrel{~}{\varphi }_k)}=|Q_\mathrm{I}Q_{\mathrm{I}\mathrm{I}}|e^{i(\mathrm{\Psi }_\mathrm{I}\mathrm{\Psi }_{\mathrm{I}\mathrm{I}})}.$$
(41)
Using Eq. (7), the average over many events gives
$$Q_\mathrm{I}Q_{\mathrm{I}\mathrm{I}}^{}=\sqrt{M_\mathrm{I}M_{\mathrm{I}\mathrm{I}}}\left(v^2+e^{i(\varphi _j\stackrel{~}{\varphi }_k)}_c\right).$$
(42)
This equation is analogous to Eq. (25), with the important difference that autocorrelations \[the first term in the right-hand side of Eq. (25)\] no longer appear. As a consequence, $`Q_\mathrm{I}Q_{\mathrm{I}\mathrm{I}}^{}`$ vanishes if there are no azimuthal correlations between particles.
However, two-particle nonflow correlations do remain. Since they are of order $`1/N`$, Eq. (42) can be written
$$Q_\mathrm{I}Q_{\mathrm{I}\mathrm{I}}^{}=\sqrt{M_\mathrm{I}M_{\mathrm{I}\mathrm{I}}}\left[v^2+O\left(\frac{1}{N}\right)\right].$$
(43)
One recognizes in the right-hand side of this equation the flow and nonflow contributions to two-particle azimuthal correlations, as in Eq. (8). The only difference lies in the global multiplicative factor $`\sqrt{M_\mathrm{I}M_{\mathrm{I}\mathrm{I}}}`$. In particular, summing over many particles does not decrease the relative weight of nonflow correlations, as might be believed: they add up in the same way as the correlations due to flow.
#### 3 Beyond the standard method
Now, following the procedure outlined in Sec. III B, it is possible to eliminate nonflow correlations between $`Q_\mathrm{I}`$ and $`Q_{\mathrm{I}\mathrm{I}}`$ order by order. This is done by means of a cumulant expansion, which is a trivial generalization of the one presented previously. The equation to an arbitrary order $`2k`$ is obtained by replacing, in Eqs. (32) and (35), $`|Q|^2`$ with $`Q_\mathrm{I}Q_{\mathrm{I}\mathrm{I}}^{}`$, and $`Q^2`$ with $`Q_\mathrm{I}Q_{\mathrm{I}\mathrm{I}}`$. For example, Eqs. (29) and (34) become
$$Q_\mathrm{I}^2Q_{\mathrm{I}\mathrm{I}}^22Q_\mathrm{I}Q_{\mathrm{I}\mathrm{I}}^{}^2=M_\mathrm{I}M_{\mathrm{I}\mathrm{I}}\left[v^4+O\left(\frac{1}{N^3}\right)\right],$$
(44)
which allows measurements of the flow when $`v`$ is much larger than $`1/N^{3/4}`$: the sensitivity is better than with Eq. (43), where $`v`$ is to be compared with $`1/N^{1/2}`$. The term $`1/M`$ in Eq. (34), which reflects autocorrelations, is automatically removed in Eq. (44).
To make a long story short, the same techniques apply to subevents as to the whole event. The only interest of subevents is that they remove autocorrelations. However, they do not remove direct nonflow correlations, which may be of the same order of magnitude. Furthermore, autocorrelations can be also be subtracted systematically when working with the whole event, as shown in Sec. III B. Another drawback of the subevent method is that each subevent contains at most half of the total multiplicity, resulting in increased errors. As a conclusion, the subevent method seems to be obsolete when working with cumulants.
### D Statistical errors
The cumulant expansion allows in principle the measurement of $`v`$ down to values of $`1/N`$ by going to large orders $`k`$, as explained in Sec. II D. In practice, however, since the number of events $`N_{\mathrm{evts}}`$ used in the analysis is finite, the sensitivity is limited by statistical errors. In this section, we determine, as a function of $`M`$ and $`N_{\mathrm{evts}}`$, which order of the cumulant expansion should be chosen so as to obtain the most accurate value of the integrated flow.
First, there is a “systematic” error, which is the error due to nonflow correlations. Expanding Eq. (35) to order $`2k`$, we obtain an equation relating the measured cumulant $`|Q|^{2k}`$ and the integrated flow $`Q`$, which is of the type
$$|Q|^{2k}=a_kQ^{2k}+O\left(M^{1k}\right),$$
(45)
where $`a_k`$ is a numerical coefficient of order unity, and the last term is the systematic error. The resulting error on $`Q`$ is therefore
$$\delta Q_{\mathrm{syst}}Q^{12k}M^{1k}$$
(46)
The systematic error thus decreases with increasing $`k`$, since $`Q\sqrt{M}=Mv_n1`$, as assumed in Eq. (21).
Let us now discuss the statistical error. When averaging a quantity over a large number of events $`N_{\mathrm{evts}}`$, the statistical error is generally of relative order $`1/\sqrt{N_{\text{evts}}}`$. Since the moments of the distribution of $`|Q|^2`$ are of order unity, the absolute statistical error on the moments is of order $`1/\sqrt{N_{\text{evts}}}`$. The same error applies to the cumulants, which are constructed from the moments. If there is no flow, a more accurate calculation shows that the statistical error on the cumulant is
$$\delta |Q|^{2k}_{\mathrm{stat}}=\frac{k!}{\sqrt{N_{\mathrm{evts}}}}.$$
(47)
If the flow is weak, that is if $`Q1`$, this formula still holds approximately. Using Eq. (45), one thus derives the statistical error on the integrated flow
$$\delta Q_{\mathrm{stat}}Q^{12k}N_{\mathrm{evts}}^{1/2}.$$
(48)
Since we have assumed that $`Q1`$, the statistical error increases with $`k`$, unlike the systematic error.
It is very likely that this property still holds in the more general case when $`Q`$ is not much smaller than unity. However, we have not been able to derive a general formula for the statistical error for arbitrary $`Q`$ and $`k`$. We only have formulas for the lowest order cumulants. Using the cumulant to order 2 ($`k=1`$),
$$\delta Q_{\mathrm{stat}}=\frac{1}{2Q}\sqrt{\frac{1+2Q^2}{N_{\mathrm{evts}}}},$$
(49)
and with the fourth order cumulant ($`k=2`$),
$$\delta Q_{\mathrm{stat}}=\frac{1}{2Q^3}\sqrt{\frac{1+4Q^2+Q^4+2Q^6}{N_{\mathrm{evts}}}}.$$
(50)
One sees on these two formulas that for very strong flow ($`Q1`$), the statistical error $`\delta Q_{\mathrm{stat}}`$ is of order $`1/\sqrt{N_{\mathrm{evts}}}`$, independent of $`Q`$. This remains true for higher order cumulants. Note, moreover, that both formulas give Eq. (48) in the limit of small $`Q`$.
Since the systematic error decreases with $`k`$ and the statistical error increases with $`k`$, the best accuracy is achieved for the value of $`k`$ such that both are of the same order of magnitude. Using Eqs. (46) and (48), one thus obtains the optimal value of the order $`2k`$:
$$2k2+\frac{\mathrm{ln}N_{\text{evts}}}{\mathrm{ln}M}.$$
(51)
Since, in practice, $`M`$ is at least of the order of a hundred at ultrarelativistic energies, the fourth order cumulant $`2k=4`$ \[i.e. Eq. (34)\] gives the best accuracy if the number of events lies in the range $`10^3<N_{\mathrm{evts}}<10^6`$. Higher order cumulants may be useful if a large statistics is available and/or if the multiplicity $`M`$ is low, as for instance in a peripheral collision.
The flow is detectable only if $`Q`$ is larger than both statistical and systematic errors. Taking for instance $`N_{\mathrm{evts}}=10^5`$ and $`M=300`$, statistical and systematic errors are of the same order. One then obtains, using Eq. (50), that flow can be seen if $`Q>0.3`$. Using Eq. (23), $`v`$ can be measured down to $`1.6\%`$ using the fourth order cumulant. If $`v=3\%`$, a typical value at the CERN SPS, then $`Q0.5`$. Using Eq. (50), the typical error is then $`\delta Q0.02`$, i.e. $`\delta v=0.1\%`$.
### E Weighted $`Q`$-vectors
The vector $`Q_n`$ has been defined in Eq. (22) with unit weights. A more general definition is
$$Q_n=\frac{1}{\sqrt{_{j=1}^Mw_j^2}}\underset{j=1}{\overset{M}{}}w_je^{in\varphi _j},$$
(52)
where the weight $`w_j`$ is an arbitrary function of $`p_T`$, $`y`$, the particle type, and the order of the harmonic under study. As a consequence, we shall restore the index $`n`$ in this subsection.
#### 1 Flow analysis with arbitrary weights
The method exposed in Sec. III B also applies with this more general definition. There are only two slight differences. The first is that the average value of $`Q_n`$, which we have denoted by $`Q_n`$, is no longer related to the average value $`v_n`$ of the flow by Eq. (23). This modification is not important for what follows: we shall see in Sec. IV B that measurements of differential flow depend on the value of $`Q_n`$ rather than $`v_n`$. The second difference is that autocorrelations cannot be removed so simply: the procedure given in Appendix B 4 is no longer valid, so that the subevent method, which avoids autocorrelations, may regain some interest.
Apart from this difference, the procedure is the same as in Sec. III B. In particular, the cumulants of the event flow vector distribution are expressed in the same way in terms of the moments. The generating equation (35) still holds, with the caveat that the last term, corresponding to autocorrelations, is no longer exact.
However, autocorrelations are unchanged at the lowest order: a calculation analogous to the one leading to Eq. (25) shows that $`|Q_n|^2=1`$ if there are no azimuthal correlations between particles, up to terms of order $`1/M`$. Changes occur only at higher orders.
#### 2 Optimal weights
What is the best choice for the weight $`w(p_T,y,n)`$? In practice, it should be chosen so as to maximize the effect of flow: one should try to obtain a value of $`Q_n`$ as large as possible, since this value will determine the accuracy in the measurement of azimuthal distributions, as we shall see in Sec. IV. From the definition (52), averaging over azimuthal angles and denoting by $`(v_n)_j`$ the value of $`v_n`$ for the corresponding particle, one obtains
$$Q_n=\frac{\underset{j=1}{\overset{M}{}}(v_n)_jw_j}{\sqrt{_{j=1}^Mw_j^2}}\sqrt{\underset{j=1}{\overset{M}{}}(v_n)_j^2},$$
(53)
where we have used a simple triangular inequality, and the fact that the flow coefficients $`(v_n)_j`$ are real. The identity holds when $`w_j=\lambda (v_n)_j`$, where $`\lambda `$ is arbitrary. In other terms, the optimal weight for a particle with given rapidity and transverse momentum is the associated flow coefficient $`(v_n)_j`$ itself.
Of course, since the goal is precisely to measure $`v_n`$, the above discussion does not answer the question of the choice of the optimal weight. However, general properties of the $`v_n`$’s can be used to guess a reasonable choice of $`w`$. Since $`v_n`$ is an odd (resp. even) function of the center of mass rapidity for odd $`n`$ (resp. even $`n`$), so should be $`w`$. Regarding the $`p_T`$ dependence, one may note that at low $`p_T`$, $`v_n`$ generally behaves as $`v_np_T^n`$ . Therefore, it seems natural to choose $`wp_T^n`$ when measuring the $`n^{\mathrm{th}}`$ harmonic. For $`n=1`$, $`Q_n`$ then becomes the sum of transverse momenta, weighted by an odd function of rapidity, which was the definition chosen in . For $`n=2`$, $`Q_n`$ is then equivalent to the transverse momentum sphericity tensor used in .
### F Gaussian limit
In this section, we compare our method to methods previously used in , which rely on the large multiplicity, Gaussian limit. It is well known that, according to the central limit theorem, the distribution of the fluctuations of $`Q`$ around its average value $`Q`$ is Gaussian in the limit of large $`M`$. Up to corrections of order $`1/M`$, the normalized probability of $`Q=Q_x+iQ_y`$ thus reads
$$\frac{\mathrm{d}p}{\mathrm{d}^2Q}=\frac{1}{2\pi \sigma _x\sigma _y}\mathrm{exp}\left(\frac{(Q_xQ)^2}{2\sigma _x^2}\frac{Q_y^2}{2\sigma _y^2}\right),$$
(54)
with $`\sigma _x^2=Q_x^2Q^2`$ and $`\sigma _y^2=Q_y^2`$.
We shall first show that this limit is equivalent to the cumulant expansion to order 4 presented in Sec. III B. Then we shall discuss the relationship with an alternative method to measure flow, which has been used in the literature, and consists in fitting the distribution of $`|Q|`$.
#### 1 Higher harmonics
In the case of the Gaussian distribution (54), one easily calculates the cumulants used in Sec. III B:
$`|Q|^2`$ $`=`$ $`Q^2+\sigma _x^2+\sigma _y^2,`$ (55)
$`|Q|^42|Q|^2^2`$ $`=`$ $`Q^4+2(\sigma _x^2\sigma _y^2)Q^2+(\sigma _x^2\sigma _y^2)^2.`$ (56)
In order to compare these equations with Eqs. (27) and (A10), we need to evaluate the sum $`\sigma ^2\sigma _x^2+\sigma _y^2`$ and the difference $`\sigma _x^2\sigma _y^2`$.
From Eqs. (23) and (25), one obtains
$`\sigma ^2`$ $`=`$ $`|Q|^2Q^2`$ (57)
$`=`$ $`1v_1^2+(M1)e^{i(\varphi _j\varphi _k)}_c,`$ (58)
where the last term is of order unity since $`e^{i(\varphi _j\varphi _k)}_c`$ is of order $`1/N1/M`$. This still holds for the generalized vector (52). One thus recovers Eq. (27).
Let us now calculate the difference:
$`\sigma _x^2\sigma _y^2`$ $`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{j,k=1}{\overset{M}{}}}\left[\mathrm{cos}(\varphi _j+\varphi _k)\mathrm{cos}\varphi _j\mathrm{cos}\varphi _k\right]`$ (59)
$`=`$ $`v_2v_1^2+O(v_2),`$ (60)
where the first two terms in the last equation come from the diagonal terms $`j=k`$, while the remaining term is the contribution of nondiagonal terms. Reporting this expression into Eq. (55), we recover Eq. (A10): higher harmonics reflect a deviation from isotropy in the fluctuations of $`Q`$.
#### 2 Isotropic fluctuations
Neglecting higher harmonics, we may write $`\sigma _x=\sigma _y`$. Then the distribution (54) becomes
$$\frac{\mathrm{d}p}{\mathrm{d}^2Q}=\frac{1}{\pi \sigma ^2}\mathrm{exp}\left(\frac{|QQ|^2}{\sigma ^2}\right).$$
(61)
With this distribution, we expect to recover the results of Sec. III B, where higher harmonics were also neglected.
Indeed, one finds after some algebra, for arbitrary, real $`x`$,
$$\mathrm{ln}I_0(2x|Q|)=\sigma ^2x^2+\mathrm{ln}I_0(2xQ),$$
(62)
to be compared with Eqs. (32) and (35). According to Eq. (58), the extra term $`\sigma ^2x^2`$ is of order unity, in agreement with the statement following Eq. (35) that the correction at order $`x^{2k}`$ is $`O(M^{1k})`$.
Corrections to the central limit theorem are of order $`1/M`$. Thus, expanding Eq. (62) in powers of $`x`$, one obtains identities which are valid up to that order. To order $`x^4`$, we recover the result obtained in the previous section, see Eq. (34), with the same accuracy. To order $`x^{2k}`$ with $`k>2`$, the results obtained in Sec. III B are more accurate since we have seen that the correction is of magnitude $`M^{1k}1/M`$.
#### 3 Distribution of $`|Q|`$
A method for extracting the flow from the data, which was proposed in , consists in plotting the measured distribution of $`|Q|`$. This method led to the first observation of collective flow in ultrarelativistic nucleus-nucleus collisions . It is a simplified version of the method based on the sphericity tensor , which led to the first observation of collective flow at Bevalac . Note that these methods are more reliable than what we call the “standard method” in this paper, in the sense that one need not neglect nonflow correlations.
The distribution of $`|Q|`$ is obtained by integration of Eq. (61) over the phase of $`Q`$:
$$\frac{1}{|Q|}\frac{\mathrm{d}p}{\mathrm{d}|Q|}=\frac{2}{\sigma ^2}\mathrm{exp}\left(\frac{Q^2+|Q|^2}{\sigma ^2}\right)I_0\left(\frac{2|Q|Q}{\sigma ^2}\right).$$
(63)
One must then fit both parameters $`\sigma `$ and $`Q`$ to the data.
If there is no flow, that is $`Q=0`$, the $`|Q|`$ distribution given by Eq. (63) is purely Gaussian:
$$\frac{1}{|Q|}\frac{\mathrm{d}p}{\mathrm{d}|Q|}=\frac{2}{\sigma ^2}\mathrm{exp}\left(\frac{|Q|^2}{\sigma ^2}\right).$$
(64)
The $`|Q|`$ distribution deviates from the Gaussian shape if the flow is strong enough compared to the fluctuation scale, that is for values of $`Q\sigma `$. In particular, the maximum of the distribution is shifted to $`|Q|0`$ if $`Q>\sigma `$. Since $`\sigma `$ is of order 1, using Eq. (23), this condition is equivalent to $`v1/\sqrt{M}`$. Note, however, that one need not assume $`v1/\sqrt{M}`$, as with the methods based on two-particle azimuthal correlations.
If $`Q\sigma `$, i.e. $`v1/\sqrt{M}`$, the shape of the distribution is very close to a pure Gaussian distribution. In fact, the deviations from the Gaussian shape are of order $`Q^4/\sigma ^4`$. This can be seen by expanding Eq. (63) to order $`Q^2`$, which is equivalent to replacing $`\sigma ^2`$ with $`\sigma ^2+Q^2`$ in Eq. (64). Alternatively, one can eliminate $`\sigma `$ and obtain $`Q`$ directly using the following identity, which can be easily derived from Eq. (61):
$$|Q|^42|Q|^2^2=Q^4,$$
(65)
again showing that the deviation is of fourth order in the flow. Knowing that the deviation to the central limit is of order $`1/M`$, one finds that Eq. (65) is equivalent to Eqs. (29) and (34), i.e., to the cumulant expansion to order 4.
The importance of the factor $`1/\sqrt{M}`$ in the definition of $`Q`$, Eq. (22), also appears clearly when fitting Eq. (63) to experimental data. Because of this factor, $`\sigma `$ does not depend on the multiplicity $`M`$ in the limit of large $`M`$, as discussed above. This is especially important when the fit is done using events with different multiplicities $`M`$. If there is no flow, the distribution of $`|Q|`$ is Gaussian with width $`\sigma `$. If $`\sigma `$ depended on $`M`$, the distribution would rather be a superposition of Gaussian distributions with different widths. In this case, the left-hand side of Eq. (65) would be positive, hiding a possible weak flow. This phenomenon probably explains why the first analysis of the E877 Collaboration gives zero values of the flow in some centrality bins.
When fitting Eq. (63) to the data, it is important to fit independently $`Q`$, which reflects the flow, and $`\sigma `$, which also involves two-particle correlations, according to Eq. (58). Assuming that $`\sigma `$ is the same for all Fourier harmonics, as was done by E877 , amounts to neglecting two-particle correlations.
Finally, note that the Gaussian limit can also be applied to the subevent method, yielding interesting results: in particular, the distribution of the relative angle between $`Q_\mathrm{I}`$ and $`Q_{\mathrm{I}\mathrm{I}}`$ is not the same for direct correlations and correlations due to flow .
## IV Differential flow
In this section, we explain how it is possible to perform detailed measurements of azimuthal distributions: typically, one wishes to measure $`v_n`$ for a given type of particle as a function of the rapidity $`y`$ and the transverse momentum $`p_T`$. In the following, we shall call this particle a “proton,” but it can be anything else. We denote by $`\psi `$ its azimuthal angle, and by $`v_m^{}`$ the corresponding differential flow coefficients $`v_m^{}=e^{im\psi }`$. Unlike the standard method, as stated before, we do not make the assumption that all azimuthal correlations are due to flow. As in the case of the integrated flow studied in Sec. III, we get rid of nonflow correlations order by order, by means of a cumulant expansion.
The principle of the method is explained in Sec. IV A. In Sec. IV B, we show that $`v_m^{}`$ can be obtained from the azimuthal correlation between $`\psi `$ and the flow vector $`Q`$. As in the case of integrated flow, the order to which nonflow correlations must be eliminated depends in practice on the number of events available: this is explained in Sec. IV C, where we also estimate the resulting accuracy on $`v_m^{}`$. Our method is compared to traditional methods in Sec. IV D.
### A Principle and orders of magnitude
The differential flow coefficients $`v_m^{}`$ can be obtained only through azimuthal correlations with other particles, typically particles used to estimate the orientation of the reaction plane, which we call “pions” in this section, although they can be anything else.
For instance, correlating the proton with one pion, $`v_1^{}`$ can be derived from the measurement of the two-particle azimuthal correlation
$$e^{i(\psi \varphi _1)}=v_1^{}v_1+O\left(\frac{1}{N}\right),$$
(66)
where $`v_1`$ refers to the pion, and is determined independently. We have used an analogy with Eq. (8). The term $`O(1/N)`$ comes from two-particle nonflow correlations between the proton and the pion. The error made in the determination of $`v_1^{}`$ is thus of order $`1/(Nv_1)`$. Of course, one should correlate the proton to particles with a strong flow, so that $`v_1`$ be as large as possible.
More accurate measurements can be obtained using higher order correlations and a cumulant expansion. For instance, at fourth order, one can eliminate the two-particle nonflow correlation by correlating the proton with three pions and taking the cumulant, by analogy with Eqs. (17) and (18):
$`e^{i(\psi +\varphi _1\varphi _2\varphi _3)}`$ $``$ $`e^{i(\psi +\varphi _1\varphi _2\varphi _3)}e^{i(\psi \varphi _2)}e^{i(\varphi _1\varphi _3)}e^{i(\psi \varphi _3)}e^{i(\varphi _1\varphi _2)}`$ (67)
$`=`$ $`v_1^{}v_1^3+O\left({\displaystyle \frac{1}{N^3}}\right).`$ (68)
More generally, correlating the proton with $`2k+1`$ pions, the connected part of the correlation is of order $`1/N^{2k+1}`$ \[since it corresponds to direct $`(2k+2)`$-particle correlations\], while the contribution of flow is $`v_1^{}v_1^{2k+1}`$. Comparing both terms, the accuracy on $`v_1^{}`$ is thus of order $`1/(Nv_1)^{2k+1}`$. Using Eq. (21), this shows that the accuracy increases with increasing $`k`$, i.e. when using multiparticle correlations.
Higher harmonics, such as $`v_2^{}`$, can be obtained by at least two methods. The first consists in multiplying all the angles by 2 in the equations above, and replacing $`v_1^{}`$ and $`v_1`$ with $`v_2^{}`$ and $`v_2`$, respectively. A second method is to mix two different harmonics, measuring $`e^{i(2\psi \varphi _1\varphi _2)}`$. If the source is isotropic, this quantity is of order $`1/N^2`$ since it involves a direct three-particle correlation. If there is flow, neglecting other sources of correlation for simplicity, $`e^{i(2\psi \varphi _1\varphi _2)}`$ factorizes into $`e^{2i\psi }e^{i\varphi _1}e^{i\varphi _2}=v_2^{}v_1^2`$. Putting everything together, we obtain
$$e^{i(2\psi \varphi _1\varphi _2)}=v_2^{}v_1^2+O\left(\frac{1}{N^2}\right).$$
(69)
One sees that nonflow correlations come into play only at order $`1/N^2`$, rather than $`1/N`$ when comparing the same harmonics as in (66). Nonetheless, they do not disappear. These correlations can also be eliminated order by order using the cumulant expansion, as we shall see in Sec. IV B. Generally, if one correlates the proton with $`2k+m`$ pions, one obtains an accuracy on $`v_m^{}`$ of order $`1/(Nv_1)^{2k+m}`$.
### B Differential flow from correlations with $`Q_n`$
In order to correlate a proton with pions, it is convenient to use the event flow vector $`Q_n`$, Eq. (22). From now on in this section, we choose $`n=1`$, and drop the subscript $`n`$, i.e. we write $`Q`$ and $`v`$ instead of $`Q_1`$ and $`v_1`$. On the other hand, we keep the subscript $`m`$ for the proton $`v_m^{}`$ because several harmonics may be measured. Generalization to arbitrary $`n`$ is straightforward: one simply multiplies all azimuthal angles (of both protons and pions) by $`n`$.
In the standard flow analysis, one usually excludes “autocorrelations” by excluding the “proton” under study from the definition of the event flow vector ; that is, the azimuthal angle $`\psi `$ is not one of the $`\varphi _j`$ in Eq. (22). Within our method, one can still do so, but it is not even necessary. First, autocorrelations will be removed order by order as well as direct correlations, as in the case of the integrated flow in Sec. III B. Furthermore, autocorrelations, if any, can be subtracted exactly if the event flow vector $`Q`$ is defined with unit weight, as in Eq. (22). This subtraction is performed in Appendix C 4. For simplicity, we neglect the corresponding term in this section, unless otherwise specified.
Let us start with the measurement of the first harmonic $`v_1^{}`$. The two-particle azimuthal correlation between the proton and a pion, Eq. (66), can be expressed introducing the vector $`Q`$ defined by Eq. (22). Summing Eq. (66) over all the pions involved in $`Q`$, one obtains the correlation between $`Q`$ and the proton:
$$Q^{}e^{i\psi }=Q\left[v_1^{}+O\left(\frac{1}{Nv}\right)\right].$$
(70)
The value of $`Q`$ must be obtained independently, using the methods discussed in Sec. III.
More accurate measurements, involving correlations of the proton with several pions, are performed using higher order moments, as in Sec. III B. These higher order moments are obtained by weighting the previous expression with powers of $`|Q|^2`$, i.e. by measuring $`|Q|^{2k}Q^{}e^{i\psi }`$. These moments are then decomposed into cumulants. For instance, Eq. (68) becomes
$`|Q|^2Q^{}e^{i\psi }`$ $``$ $`|Q|^2Q^{}e^{i\psi }2Q^{}e^{i\psi }|Q|^2`$ (71)
$`=`$ $`Q^3\left[v_1^{}+O\left({\displaystyle \frac{1}{(Nv)^3}}\right)\right],`$ (72)
through which we define the cumulant $`|Q|^2Q^{}e^{i\psi }`$.
As in the case of integrated flow, the decomposition of higher order moments $`|Q|^{2k}Q^{}e^{i\psi }`$ in cumulants can be represented in terms of diagrams. For instance, the decomposition of $`|Q|^2Q^{}e^{i\psi }`$ is displayed in Fig. 4.
The diagrams in this figure stand for
$`|Q|^2Q^{}e^{i\psi }`$ $`=`$ $`2Q^{}e^{i\psi }|Q|^2+|Q|^2Q^{}e^{i\psi }`$ (73)
$`=`$ $`2Q^{}e^{i\psi }|Q|^2+|Q|^2Q^{}e^{i\psi }`$ (74)
One thus recovers the expression of the cumulant, Eq. (71). More generally, in order to decompose the moment $`|Q|^{2k}Q^{}e^{i\psi }=Q^kQ^{k+1}e^{i\psi }`$, one draws a cross on the left representing the proton, $`k`$ dots on the left and $`k+1`$ dots on the right representing the pions. The graphs combine all possible subsets of the points on the left with subsets of the points on the right containing the same number of elements.
Let us now discuss the measurements of higher harmonics of the proton azimuthal distribution $`v_m^{}`$. In the case $`m=2`$, Eq. (69) gives, summing over the pions involved in $`Q`$,
$$Q^2e^{2i\psi }=Q^2\left[v_2^{}+O\left(\frac{1}{(Nv)^2}\right)\right].$$
(75)
To obtain a better accuracy, one must decompose higher order moments $`|Q|^{2k}Q^2e^{2i\psi }`$ in cumulants. In terms of the diagrammatic representation, the proton is now associated with two crosses, as seen in Fig. 5 for $`k=1`$.
As before, the graphs combine all possible subsets of the points on the left with subsets of the points on the right containing the same number of elements, with the subsidiary condition that the two crosses belong to the same subset. In the left-hand side of Fig. 5, the dot on the left of the dashed line can be associated with any of the three dots on the right. The equation represented by the figure can be written
$`|Q|^2Q^2e^{2i\psi }`$ $`=`$ $`3Q^2e^{2i\psi }|Q|^2+|Q|^2Q^2e^{2i\psi }`$ (76)
$`=`$ $`3Q^2e^{2i\psi }|Q|^2+|Q|^2Q^2e^{2i\psi }`$ (77)
where the last term involves a direct five-particle correlation, and is therefore of order $`M^2\times O(1/N^4)`$. When there is flow, one obtains
$`|Q|^2Q^2e^{2i\psi }`$ $`=`$ $`|Q|^2Q^2e^{2i\psi }3Q^2e^{2i\psi }|Q|^2`$ (78)
$`=`$ $`Q^4\left[2v_2^{}+O\left({\displaystyle \frac{1}{(Nv)^4}}\right)\right].`$ (79)
Cumulants of arbitrary order, for arbitrary harmonics $`v_m^{}`$, can be obtained by expanding in powers of $`x`$ the following generating equation, derived in Appendix C 2:
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{x^{2k+m}}{k!(k+m)!}|Q|^{2k}Q_{}^{}{}_{}{}^{m}e^{im\psi }=\frac{I_m(2x|Q|)(Q^{}/|Q|)^me^{im\psi }}{I_0(2x|Q|)}$$
(80)
where $`I_m`$ is the modified Bessel function of order $`m`$. For $`m=1`$, one recovers Eq. (74) by expanding this equation to order $`x^3`$. For $`m=2`$, one recovers Eq. (76) by expanding this equation to order $`x^4`$.
The cumulants defined by Eq. (80) are related to the differential flow by
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{x^{2k+m}}{k!(k+m)!}|Q|^{2k}Q_{}^{}{}_{}{}^{m}e^{im\psi }=\frac{I_m(2xQ)}{I_0(2xQ)}v_m^{}+\frac{I_m\left(2x/\sqrt{M}\right)}{I_0\left(2x/\sqrt{M}\right)}$$
(81)
The second term corresponds to autocorrelations, and must be included only if the proton is involved in the flow vector $`Q`$. In the case $`m=1`$, one recovers the lowest order formulas (70) and (71) by expanding this equation to order $`x`$ and $`x^3`$, respectively. For $`m=2`$, one recovers Eqs. (75) and (78) by expanding it to order $`x^2`$ and $`x^4`$, respectively.
At order $`x^{2k+m}`$, Eq. (81) gives an accuracy on $`v_m^{}`$ of order $`1/(Nv)^{2k+m}`$, as expected from the discussion of Sec. IV A.
### C Statistical errors
Equation (81) generates an infinite set of equations to measure the differential flow $`v_m^{}`$, since it can be expanded to any arbitrary order $`x^{2k+m}`$. As in the case of integrated flow, the best choice of $`k`$ is the one that yields the best accuracy on $`v_m^{}`$. It results from a compromise between systematic errors stemming from nonflow correlations, which decrease when using higher order cumulants, and statistical errors, which increase with the order $`k`$.
The equation obtained when expanding Eq. (81) to order $`x^{2k}`$ is of the type
$$|Q|^{2k}Q^me^{im\psi }=b_kQ^{2k+m}v_m^{}+O\left(M^{k(m/2)}\right),$$
(82)
where $`b_k`$ is a numerical coefficient of order unity. Neglecting for the moment the error on the integrated flow $`Q`$, this equation gives a systematic error on $`v_m^{}`$
$$(\delta v_m^{})_{\mathrm{syst}}Q^{2km}M^{k(m/2)}.$$
(83)
This systematic error decreases when increasing the order $`k`$. Note that for $`m=1`$, the systematic error should be the same on the differential flow as on its integrated value at a given order. Thus we expect Eqs. (83) and (46) to give the same result, using (23). In doing the comparison, one must pay attention to the fact that the cumulant used for differential flow $`|Q|^{2k}Q^{}e^{i\psi }`$ involves $`2k+2`$ particles while the cumulant for integrated flow $`|Q|^{2k}`$ involves only $`2k`$ particles. Thus, comparing the two “at a given order” means that we must replace $`2k`$ in Eq. (83) by $`2k+2`$ in Eq. (46).
The statistical error on the cumulant (82) is of order $`1/\sqrt{N_{\mathrm{evts}}^{}}`$, where $`N_{\mathrm{evts}}^{}`$ is the number of events containing a proton. This leads to an error on $`v_m^{}`$:
$$(\delta v_m^{})_{\mathrm{stat}}Q^{2km}\left(N_{\mathrm{evts}}^{}\right)^{1/2}$$
(84)
If $`m=1`$, we again recover the result obtained for the integrated flow, provided we replace $`2k`$ by $`2k+2`$ in Eq. (48), and $`N_{\mathrm{evts}}^{}`$ by $`MN_{\mathrm{evts}}`$ (which is the total number of particles involved in the measurement of integrated flow) in Eq. (84).
If $`Q<1`$, the statistical error (84) increases with increasing $`k`$, and the optimal value of $`k`$ is that for which statistical and systematic errors are equivalent, i.e.
$$2km+\frac{\mathrm{ln}N_{\text{evts}}^{}}{\mathrm{ln}M}.$$
(85)
The error on $`Q`$, estimated in Sec. III D, should also be taken into account. However, the measurement of differential flow is done in a limited region of phase space, by definition, so that the corresponding statistics is smaller than for the integrated flow where many more events can be used. It is then safe to assume that the statistical error on $`Q`$ gives a negligible contribution to the error on $`v_m^{}`$.
If $`m=1`$, the previous equation shows that $`k=1`$ is more accurate than $`k=0`$ (the latter value corresponds to the standard method, neglecting correlations) only if $`\mathrm{ln}N_{\mathrm{evts}}^{}/\mathrm{ln}M>2`$, i.e. if the statistics is large enough, typically $`N_{\mathrm{evts}}^{}>10^4`$ for an event multiplicity $`M100`$. For higher harmonics $`m>1`$, the contribution of nonflow correlations are smaller as explained above: thus the lowest order method $`k=0`$ is to be chosen unless a very large number of events is available, typically $`N_{\mathrm{evts}}>10^6`$ for the second harmonic $`m=2`$ if $`M100`$.
### D Relation with previous methods
Previously used methods also study the correlation between the event flow vector (22) with the momentum of the proton. The traditional justification is that, as explained in Sec. III A, the phase $`n\varphi _Q`$ of the event flow vector (22) gives an estimate of the orientation of the reaction plane modulo $`2\pi /n`$. Studying the correlation between $`\psi `$ and $`\varphi _Q`$, one can reconstruct harmonics $`v_n^{}`$, $`v_{2n}^{}`$, $`v_{3n}^{}`$, etc.
The standard analysis relies on a purely angular correlation. One measures the average $`\mathrm{cos}m(\psi \varphi _Q)`$. Neglecting nonflow correlations, this quantity is the product of $`v_m^{}`$ and a resolution factor which is given by an independent measurement. Our method relies on similar averages, weighted by powers of $`|Q|`$:
$$|Q|^{2k+m}\mathrm{cos}m(\psi \varphi _Q)=Q^{m+k}Q^ke^{im\psi }.$$
(86)
In the traditional method, autocorrelations are usually removed explicitly by specifying that the proton under study is not used in constructing $`Q_n`$ in Eq. (22) . However, nonflow direct correlations, which are of the same order of magnitude as autocorrelations, do remain, and limit the sensitivity of the analysis. With our method, autocorrelations can be removed in the same way as in the standard analysis. But we also remove direct correlations, thereby increasing the sensitivity of the measurements.
## V Acceptance corrections
For simplicity, the discussion has been limited so far to an ideal detector, i.e. a detector with an acceptance which is azimuthally isotropic in $`\varphi `$. An actual detector is never perfect, either because its components are of uneven quality, or simply because it does not cover the whole $`\varphi `$ range. In this section, we discuss a simple extension of the method which allows us to work with any detector. More precisely, it allows the detection of deviations from an isotropic source, i.e. flow, with any detector, and the correction is implemented in the same way for all detectors. However, the accuracy on the measurement of $`v_n`$ can be poor if the detector covers only a limited range in $`\varphi `$.
The only modification lies in the definition of the cumulants, for which the expressions given in Secs. III B and IV B are no longer valid. These modified cumulants are defined in Sec. V A for integrated flow and in Sec. V B for differential flow. As we shall see, the analytical expression of higher order cumulants become very lengthy, so that it is more convenient to work directly at the level of generating functions. As an illustration of our method, results of a simple Monte-Carlo simulation are given in Sec. V C.
### A Integrated flow
The key idea is that anisotropies in the detector acceptance can be handled much in the same way as anisotropies of the emitting source. The only difference is that the relevant coordinate system is the laboratory system in the first case, and the system associated with the reaction plane in the second case.
Let us be more specific: until now, we have been working in the coordinate system associated with the reaction plane, i.e. with the emitting source. In this system, we used a cluster expansion to define direct $`k`$-particle correlations, of order $`N^{1k}`$ relative to the uncorrelated $`k`$-particle distribution (Sec. II C). This cluster expansion allowed us to construct the “connected moments” of the distribution of $`Q`$, which were noted as $`Q^kQ^l_c`$ (Appendix B 1), of order $`M^{1kl}`$ relative to the corresponding moment $`Q^kQ^l`$. This decomposition was performed for an arbitrary source, but with an ideal detector.
Exchanging the roles played by the source and the detector, the same reasoning applies if we work with an isotropic source and an imperfect detector, provided we use the coordinate system associated with the detector. We thus define the connected moments exactly in the same way, replacing $`\varphi `$ by the measured $`\overline{\varphi }`$ (see Sec. II A). Similarly, the flow vector $`Q`$ will be denoted $`\overline{Q}`$ when azimuthal angles are measured in the laboratory system, i.e. when $`\varphi _j`$ is replaced by $`\overline{\varphi }_j`$ in the definition (22).
I f the acceptance is not perfect, averages such as $`e^{in\overline{\varphi }}`$ or $`e^{in(\overline{\varphi }_1+\overline{\varphi }_2\overline{\varphi }_3)}`$ no longer vanish. Thus, nondiagonal moments $`\overline{Q}^k\overline{Q}^l`$ with $`kl`$ are also nonvanishing: there is no more cancellation due to isotropy, and all terms must be kept in the cumulant expansion. At order 2, for instance, the cumulants are defined as
$`\overline{Q}^2`$ $``$ $`\overline{Q}^2\overline{Q}^2,`$ (87)
$`|\overline{Q}|^2`$ $``$ $`|\overline{Q}|^2\overline{Q}\overline{Q}^{}.`$ (88)
These cumulants are of the same magnitude as when the acceptance is perfect, i.e. of order unity, while the moments $`\overline{Q}^2`$ and $`|\overline{Q}|^2`$ scale like the multiplicity $`M`$ if the detector is very bad. Note that at this order ($`k+l=2`$), taking the cumulant is equivalent to shifting the distribution of $`\overline{Q}`$ by its average value $`\overline{Q}`$, as proposed in .
Higher order cumulants can be obtained in similar way as for an ideal detector. The only difference is that the simplifications due to isotropy no longer exist. Thus one cannot use expression (B10) for the generating function of the moments; one must use instead the more general expression (B7). The cumulant $`\overline{Q}^k\overline{Q}^l`$ is therefore defined by
$$\underset{k,l}{}\frac{z^kz^l}{k!l!}\overline{Q}^k\overline{Q}^l=\mathrm{ln}𝒢_0(z)=\mathrm{ln}e^{z^{}\overline{Q}+z\overline{Q}^{}}.$$
(89)
Expanding the right-hand side to order $`z^kz^l`$, one obtains the cumulant $`\overline{Q}^k\overline{Q}^l`$ as a function of the measured moments $`\overline{Q}^k^{}\overline{Q}^l^{}`$ with $`k^{}k`$ and $`l^{}l`$. While the moment $`\overline{Q}^k\overline{Q}^l`$ is of magnitude $`M^{(k+l)/2}`$ for a bad detector, the corresponding cumulant $`\overline{Q}^k\overline{Q}^l`$ is of order $`M^{(k+l)/2}N^{1kl}M^{1(k+l)/2}`$.
If the acceptance is not too bad, we assume that relation (35) between the cumulants and the integrated flow is approximately preserved. The integrated flow can then obtained from the cumulants to order 2, 4, 6 using Eqs. (27), (34) and (36), which we write again in the form:
$`Q^2`$ $`=`$ $`|\overline{Q}|^2|1+O(1)\pm \sqrt{{\displaystyle \frac{1+2Q^2}{N_{\mathrm{evts}}}}},`$ (91)
$`Q^4`$ $`=`$ $`|\overline{Q}|^4{\displaystyle \frac{1}{M}}+O\left({\displaystyle \frac{1}{M}}\right)\pm 2\sqrt{{\displaystyle \frac{1+4Q^2+Q^4+2Q^6}{N_{\mathrm{evts}}}}},`$ (92)
$`Q^6`$ $`=`$ $`{\displaystyle \frac{1}{4}}|\overline{Q}|^6{\displaystyle \frac{1}{M^2}}+O\left({\displaystyle \frac{1}{M^2}}\right)\pm {\displaystyle \frac{3}{2\sqrt{N_{\mathrm{evts}}}}},`$ (93)
where, in the right-hand side of each equation, the last three terms stand for autocorrelations, systematic errors due to direct $`2k`$-particle correlations, and statistical errors due to the finite number of events (see Sec. III D), respectively. Note that $`Q`$ denotes the average value of $`Q`$ in the coordinate system associated with the reaction plane, i.e. what we call the “integrated flow”. It must not be mistaken for $`\overline{Q}`$ \[see for instance Eq. (87)\], which denotes the average value in the laboratory coordinate system, and vanishes if the acceptance is perfect. Note also that only the “diagonal cumulants” $`|\overline{Q}|^{2k}`$ (i.e. with $`k=l`$) are related to the flow. These diagonal cumulants could equivalently be written $`|Q|^{2k}`$ since $`Q`$ and $`\overline{Q}`$ differ only by a phase. Other cumulants, $`Q^kQ^l`$ with $`kl`$, are not influenced by the flow and vanish except for statistical and systematic errors. They can therefore be used to estimate the magnitude of errors.
The modified definition of higher order cumulants involves a large number of terms when the detector acceptance is nonisotropic. For instance, the fourth-order cumulant is obtained by expanding Eq. (89) to order $`z^2z^2`$:
$`|\overline{Q}|^4=|\overline{Q}|^4`$ $``$ $`2\overline{Q}\overline{Q}\overline{Q}^22\overline{Q}^{}\overline{Q}^{}\overline{Q}^22|\overline{Q}|^2^2\overline{Q}^2\overline{Q}^2`$ (94)
$`+`$ $`8\overline{Q}\overline{Q}^{}|\overline{Q}|^2+2\overline{Q}^2\overline{Q}^2+2\overline{Q}^{}^2\overline{Q}^26\overline{Q}^2\overline{Q}^{}^2.`$ (95)
This equation replaces Eq. (29) for an imperfect detector. It shows that implementing acceptance corrections order by order can be very tedious since it involves a large number of terms.
It is simpler to work directly with generating functions. Although this might seem to be more complicated, it is not unnatural since the generating functions constructed from experimental data have the same geometrical properties as the data, in particular regarding the detector acceptance. For instance, when the detector is isotropic, so is the generating function Eq. (B10).
One can compute numerically the generating function of the cumulants $`𝒢_0(x,y)`$ at various points in the complex plane, then extract numerically the coefficients at a given order by means of an interpolating polynomial. Let us be more specific: separating the real and imaginary part of the flow vector, we write it as
$`\overline{Q}_x`$ $``$ $`{\displaystyle \frac{1}{\sqrt{_{j=1}^Mw_j^2}}}{\displaystyle \underset{j=1}{\overset{M}{}}}w_j\mathrm{cos}(n\overline{\varphi }_j),`$ (96)
$`\overline{Q}_y`$ $``$ $`{\displaystyle \frac{1}{\sqrt{_{j=1}^Mw_j^2}}}{\displaystyle \underset{j=1}{\overset{M}{}}}w_j\mathrm{sin}(n\overline{\varphi }_j).`$ (97)
The generating function of the cumulants, defined by Eq. (89), is a real-valued function:
$$\mathrm{ln}𝒢_0(x,y)\mathrm{ln}e^{2x\overline{Q}_x+2y\overline{Q}_y},$$
(98)
where we have set $`z=x+iy`$.
According to Eq. (89), the cumulant to order $`2k`$, $`|\overline{Q}|^{2k}`$ is the coefficient of $`(zz^{})^k=(x^2+y^2)^k`$ in the power series expansion of this generating function, up to a factor $`1/(k!)^2`$:
$$\mathrm{ln}𝒢_0(x,y)\underset{k=1}{\overset{\mathrm{}}{}}\frac{|\overline{Q}|^{2k}}{(k!)^2}(x^2+y^2)^k,$$
(99)
where we have kept only the relevant terms in the expansion. The cumulant can be obtained from the tabulated values of $`\mathrm{ln}𝒢_0(x,y)`$ using the interpolation formulas given in Appendix D 1.
### B Differential flow
When measuring differential flow, acceptance corrections can be implemented in the same way as for integrated flow. Flow is extracted using the same formulas as when the detector is perfectly isotropic in azimuth (Sec. IV), without the simplifications allowed by isotropy. Therefore, one must take as the generating function of the cumulants $`𝒞_m(z)`$ the general expression (C7) instead of Eq. (C11). We thus define the cumulants by
$$\underset{k,l}{}\frac{z^kz^l}{k!l!}\overline{Q}^k\overline{Q}^le^{im\overline{\psi }}=𝒞_m(z)\frac{e^{z^{}\overline{Q}+z\overline{Q}^{}+im\overline{\psi }}}{e^{z^{}\overline{Q}+z\overline{Q}^{}}}$$
(100)
where $`\overline{\psi }`$ denotes the azimuthal angle of the proton, measured in the laboratory coordinate system. This equation replaces Eq. (80) for an imperfect detector. Expanding Eq. (100) to order $`z`$ for $`m=1`$, we obtain for instance
$$\overline{Q}^{}e^{i\overline{\psi }}\overline{Q}^{}e^{i\overline{\psi }}\overline{Q}^{}e^{i\overline{\psi }}$$
(101)
We assume that the relation (81) between the cumulants and the differential flow, $`v_m^{}`$, is approximately preserved if the acceptance is not too bad. For $`k=0`$ and $`k=1`$, flow is then related to the cumulants by Eqs. (70) and (71) for $`m=1`$ and by Eqs. (75) and (78) for $`m=2`$. We rewrite these formulas:
$`Qv_1^{}`$ $`=`$ $`Q^{}e^{i\psi }+O\left({\displaystyle \frac{1}{M^{1/2}}}\right)\pm {\displaystyle \frac{1}{\sqrt{N_{\mathrm{evts}}^{}}}},`$ (103)
$`Q^3v_1^{}`$ $`=`$ $`|Q|^2Q^{}e^{i\psi }+O\left({\displaystyle \frac{1}{M^{3/2}}}\right)\pm {\displaystyle \frac{1}{\sqrt{N_{\mathrm{evts}}^{}}}},`$ (104)
$`Q^2v_2^{}`$ $`=`$ $`Q^2e^{2i\psi }+O\left({\displaystyle \frac{1}{M}}\right)\pm {\displaystyle \frac{1}{\sqrt{N_{\mathrm{evts}}^{}}}},`$ (105)
$`Q^4v_2^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}|Q|^2Q^2e^{2i\psi }+O\left({\displaystyle \frac{1}{M^2}}\right)\pm {\displaystyle \frac{1}{\sqrt{N_{\mathrm{evts}}^{}}}},`$ (106)
where, in the right-hand side of each equation, the second term represents the systematic error due to direct particle correlations, while the last term is the statistical error due to the finite number of events. Note that only the cumulants $`Q^kQ^le^{im\psi }`$ with $`l=k+m`$ are related to the flow by Eqs. (V B).
We wish to recall here that the differential flow $`v_2^{}`$ might also have been obtained from the correlation between the azimuth of the proton and the event flow vector $`Q_2`$. As stated in Sec. IV B, the only modification is a multiplication of all angles by 2, so that this does not change Eqs. (100) and (107). Therefore, $`v_2^{}`$ may be deduced from Eqs. (103) and (104) by the simple substitution of $`v_1^{}`$ and $`Q`$ by $`v_2^{}`$ and $`Q_2`$ respectively.
As in the case of integrated flow, the modified definitions of the cumulants quickly involve a large number of terms when going to higher orders. Therefore, it is simpler in practice to extract the cumulants numerically from the generating function. For this purpose, one must tabulate numerically the real and imaginary parts of $`𝒞_m(z)`$:
$`\mathrm{}[𝒞_m(x,y)]`$ $`=`$ $`{\displaystyle \frac{e^{2x\overline{Q}_x+2y\overline{Q}_y}\mathrm{cos}(m\overline{\psi })}{e^{2x\overline{Q}_x+2y\overline{Q}_y}}},`$ (107)
$`\mathrm{}[𝒞_m(x,y)]`$ $`=`$ $`{\displaystyle \frac{e^{2x\overline{Q}_x+2y\overline{Q}_y}\mathrm{sin}(m\overline{\psi })}{e^{2x\overline{Q}_x+2y\overline{Q}_y}}}.`$ (108)
Keeping only the terms with $`l=k+m`$ which are related to the flow, the generating function (100) becomes
$$𝒞_m(z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{|Q|^{2k}Q^me^{im\psi }}{k!(k+m)!}z^kz^{k+m}.$$
(109)
Interpolation methods to calculate the cumulants $`|Q|^{2k}Q^me^{im\psi }`$ as a function of the tabulated values of the generating function, are explained in detail in Appendix D 2.
### C Results of a Monte-Carlo simulation
We have tested our method with a simple Monte-Carlo simulation. Particles have been generated randomly with the distribution
$$\frac{\mathrm{d}N}{\mathrm{d}\varphi }1+2v_1\mathrm{cos}\varphi +2v_2\mathrm{cos}(2\varphi ).$$
(110)
The value of the integrated directed flow, which we tried to reconstruct, was fixed to $`v_1=0.03`$, corresponding roughly (up to a sign) to the value measured at SPS for pions . We have taken various values of $`v_2`$, in order to probe the interference between both harmonics, discussed in Sec. II D.
In a first step, we do not simulate nonflow correlations between the particles. In order to take into account the effect of detector inefficiencies, we have assumed that all particles are detected, except in a blind azimuthal sector of size $`\alpha `$. The simulation has been performed with $`N_{\mathrm{evts}}=200000`$ events, and a multiplicity $`M=200`$ for each event. For simplicity, we have assumed that exactly 200 particles are detected in each event. Fluctuations in $`M`$ should not influence the results, as explained in Sec. III A 3. With these values, the optimal sensitivity for the integrated flow is obtained for $`k=2`$ according to Eq. (51), i.e. by taking the fourth order cumulant. We therefore reconstruct the flow using Eq. (92).
With the values we have chosen, $`Q=v_1\sqrt{M}0.42<1`$, so that traditional methods might fail, as stated before. Within our method, the statistical error on $`v_1`$, calculated with Eq. (50), is of the order of $`0.14\%`$. Since direct correlations between particles are not simulated, the only systematic error comes from detector inefficiencies and the higher harmonic $`v_2`$.
Results are shown in the table below. The table gives the reconstructed $`v_1`$ as a function of the size of the blind angle $`\alpha `$, and the higher harmonic $`v_2`$.
If $`v_2=0`$, the reconstructed value is compatible with the theoretical value within statistical errors, except for the highest value of $`\alpha `$, i.e. when the detector covers only half of the range in azimuth. Therefore, errors due to acceptance imperfections are under good control.
The systematic error from higher harmonics, on the other hand, is far from negligible. The limits of applicability of our method, given by Eq. (A11), are here $`0.43<v_2<0.07`$. We have checked these bounds numerically. The value $`v_2=0.06`$ is very close to the upper bound. However, the corresponding relative error on $`v_1`$ is only 12% with an ideal detector.
In a second step, we simulate nonflow correlations: for simplicity, we do this assuming that particles are emitted in pairs, both particles in a pair having exactly the same azimuthal angle. This would be the case for the two-body decay of a very fast resonance. Taking the same numerical values as above, the standard method, corresponding to Eq. (91), gives $`v_1=7.7\%`$: it fails, as expected, overestimating the flow by more than a factor of 2. On the other hand, the fourth-order formula (92), which eliminates two-particle nonflow correlations, gives $`v_1=3.1\%`$, in much better agreement with the theoretical value.
## VI Summary
We have proposed in this paper a new method for the flow analysis, which is more sensitive than traditional methods to small anisotropies of the azimuthal distributions. In this section, we summarize the procedure which should be followed in practice.
The first step consists in measuring the “integrated flow,” as explained in Sec. III. This corresponds to the problem of the reaction plane determination in the standard flow analysis. One first constructs, event by event, the flow vector $`\overline{Q}_n`$ defined by Eq. (97), where the $`\overline{\varphi }_j`$ are the azimuthal angles of the particles in the laboratory coordinate system. The weight $`w_j`$ is chosen as explained in Sec. III E 2; ideally, it should be taken equal to the differential flow $`v_n(p_T,y)`$, i.e. proportional to $`p_T^n`$, and even (resp. odd) in the rapidity $`y`$ for even (resp. odd) $`n`$. Alternatively, one may choose the simpler version with unit weights (22). The value of $`n`$ depends on the system under study: up to energies of 10 GeV per nucleon, one usually works with $`n=1`$, i.e. with $`Q_1`$ . At SPS, directed flow is so small that a better accuracy is obtained by working directly with the second harmonic, i.e. by constructing $`Q_2`$ . Then, only even harmonics can be measured. Most of this paper has been written assuming $`n=1`$. In order to generalize the results to the case $`n=2`$, one need only multiply all azimuthal angles by 2.
Measuring the integrated flow amounts to measuring the average value of the flow vector, $`Q_n`$, in the coordinate system where the reaction plane is fixed. The average value $`Q_n`$ is of order $`v_n\sqrt{M}`$ (it is even equal to that value if one is working with unit weights), where $`v_n`$ is the Fourier harmonic of order $`n`$, and $`M`$ the number of particles used in the flow analysis. As explained in Sec. III, the integrated flow $`Q_n`$ is obtained from the cumulant $`|Q_n|^{2k}`$, which removes nonflow correlations up to order $`2k`$, the standard method corresponding to the lowest order, $`k=1`$. The value of $`k`$ is chosen so as to obtain the best sensitivity. It results from a balance between systematic and statistical errors, and depends both on the number of events $`N_{\mathrm{evts}}`$ available for the flow analysis, and on the number of particles used to determine the reaction plane in each event, $`M`$. The optimal order $`k`$ is then given by Eq. (51). However, performing measurements with other values of $`k`$ does not cost much and provides a useful comparison.
The cumulant $`|Q_n|^{2k}`$ is a combination of the moments of the distribution of $`\overline{Q}_n`$, i.e. it is expressed as a function of the measured moments $`\overline{Q}_n^l\overline{Q}_n^m`$, with $`lk`$ and $`mk`$. In this paper, we have used the formalism of generating functions to derive the corresponding formulas at arbitrary order. As explained in Sec. V, this is not only an elegant formalism: it is also the simplest way to calculate the cumulants numerically from experimental data. For this purpose, one tabulates the generating function $`𝒢_0(x,y)`$, defined by Eq. (98), at various points in the $`(x,y)`$ plane. In this equation, the brackets denote an average over the whole sample of events. The cumulant $`|Q_n|^{2k}`$ is then obtained by extracting numerically the coefficient in front of $`(x^2+y^2)^k`$ in the power series expansion of $`\mathrm{ln}𝒢_0(x,y)`$, as explained in Sec. V A. The integrated flow $`Q_n`$ is finally obtained from the cumulant using Eqs. (V A).
The value of $`Q_n`$ is the important parameter in the flow analysis, since it determines the accuracy of the reconstruction of azimuthal distributions. If $`Q_n>1`$, the flow can easily be studied with traditional methods, although the present method should give more accurate results. If $`Q_n<1`$, on the other hand, standard methods fail, while our method still works.
The second step in the flow analysis is to perform detailed measurements of the flow coefficient $`v_m^{}`$ for a particle of given rapidity and transverse momentum, i.e. differential flow. The coefficient $`v_m^{}`$ can be obtained from the comparison of the azimuth of the particle under study with an event flow vector, which can be either $`Q_m`$, calculated with the same harmonic, or a $`Q_n`$, calculated with a different harmonic, provided $`m`$ is a multiple of $`n`$. For instance, $`v_2^{}`$ can be measured with respect to $`Q_1`$ or $`Q_2`$, as explained in . We show in Sec. IV B that it is the value of $`Q_n`$ which determines the accuracy on the measurement of $`v_m^{}`$. Therefore, $`n`$ should be chosen so that $`Q_n`$ be as large as possible. For instance, at RHIC where $`v_2`$ is expected to be much larger than $`v_1`$, $`v_2^{}`$ should be measured with $`Q_2`$ rather than with $`Q_1`$, as is already the case at SPS . In the text, we have assumed $`n=1`$. If one uses $`Q_2`$, then $`m`$ must be replaced by $`2m`$ everywhere in our equations.
As the integrated flow, the differential flow $`v_m^{}`$ is obtained from a cumulant $`|Q|^{2k}Q^me^{im\psi }`$ which eliminates nonflow correlations up to an arbitrary order $`2k+m`$, the standard analysis corresponding to the case $`k=0`$. Here again, the best choice of $`k`$ is the one which leads to the smallest error: its value is given by Eq. (85). In order to measure the cumulants, one first tabulates the generating function (107) at various points in the complex plane. The cumulant is then obtained by extracting the coefficient proportional to $`z^kz^{k+m}`$ in the power series expansion of the generating function, as explained in Sec. V B. Finally, the differential flow $`v_m^{}`$ is related to the cumulants by Eqs. (V B).
A limitation of our method at a given order is the possible interplay of higher harmonics in the measurement. For instance, Eq. (18) shows that in the fourth-order cumulant, the second harmonic $`v_{2n}`$ interferes with $`v_n`$. More precisely, $`|v_{2n}|`$ must be small compared with $`Mv_n^2`$ (see Eq. (19)). This limitation means that the method should be used with much care when extracting the directed flow ($`n=1`$) at RHIC and LHC , since it is expected to be much smaller than elliptic flow. On the other hand, in the case $`n=2`$, there should be no problem since $`v_2`$ is much larger than $`v_4`$.
While higher harmonics or statistical errors may limit the use of the method, there is no problem with the acceptance of detectors. As a matter of fact, the required corrections appear in a natural way in the method, at all orders, from a modification of the generating equation which is the same for all detectors. In particular, the sensitivity remains unchanged when acceptance corrections are taken into account, so that choosing the order in the expansion of the generating equation does not depend on that problem.
Most of our results have been established in the limit where azimuthal anisotropies are weak. For this reason, our method seems to be more adapted to ultrarelativistic energies, i.e. at SPS energies and beyond, where $`v_1`$ and $`v_2`$ are usually less than $`0.1`$. In particular, it should be very useful in the forthcoming flow analyses at the Brookhaven Relativistic Heavy Ion Collider.
###### Acknowledgements.
We thank Art Poskanzer for helpful comments on the first version of this paper, and Raimond Snellings and Sergei Voloshin for stimulating discussions. We also thank Aihong Tang for correcting some misprints.
## A Detailed study of the four-particle azimuthal correlation
In Sec. A 1, we calculate the cumulant of the four-particle azimuthal correlation, introduced in Sec. II D. Then, in Sec. A 2, we calculate the fourth order cumulant of the $`Q`$ distribution, introduced in Sec. III B.
### 1 Cumulant of the four-particle correlation
The cumulant of the four-particle azimuthal distribution has been defined by Eq. (17) when the source is isotropic. We set $`n=1`$ for simplicity:
$$e^{i(\varphi _1+\varphi _2\varphi _3\varphi _4)}e^{i(\varphi _1+\varphi _2\varphi _3\varphi _4)}e^{i(\varphi _1\varphi _3)}e^{i(\varphi _2\varphi _4)}e^{i(\varphi _1\varphi _4)}e^{i(\varphi _2\varphi _3)}.$$
(A1)
Here, we want to evaluate the right-hand side of this equation when the source is no longer isotropic.
In order to do so, we expand the four-particle distribution into connected parts, as explained in Sec. II C. Using the diagrammatic representation introduced there, the quantity in Eq. (A1) can be decomposed as in Fig. 6.
The diagrams in Fig. 6 stand for:
$`e^{i(\varphi _1+\varphi _2\varphi _3\varphi _4)}2e^{i(\varphi _1\varphi _3)}^2=`$ $``$ $`v_1^4+2v_1^2e^{i(\varphi _3+\varphi _4)}_c+e^{i(\varphi _1+\varphi _2)}_ce^{i(\varphi _3+\varphi _4)}_c`$ (A2)
$`+`$ $`4v_1e^{\pm i(\varphi _1+\varphi _2\varphi _3)}_c+e^{i(\varphi _1+\varphi _2\varphi _3\varphi _4)}_c.`$ (A3)
Note that the direct two-particle correlations $`e^{i(\varphi _1\varphi _3)}_c`$ are automatically removed. In the isotropic case, only the connected part of the correlation, i.e. $`e^{i(\varphi _1+\varphi _2\varphi _3\varphi _4)}_c`$, remains in the right-hand side of Eq. (A2).
Let us now enumerate the orders of magnitude of the different terms in the right-hand side of Eq. (A2). As stated above, all terms but the last vanish in the isotropic case: indeed, $`e^{i(\varphi _1+\varphi _2\varphi _3)}_c`$ and $`e^{\pm i(\varphi _1+\varphi _2)}_c`$ are not invariant under the transformation $`\varphi _j\varphi _j+\alpha `$, where $`\alpha `$ is any angle. Therefore, it seems reasonable to consider that these terms are proportional to $`v_1`$ or $`v_2`$, depending on whether a factor $`e^{\pm i\alpha }`$ or $`e^{\pm 2i\alpha }`$ appears under the previous transformation. Furthermore, since we consider here connected $`k`$-particle correlations, they behave like $`O(1/N^{k1})`$ \[see Sec. II C\]. More precisely,
$$e^{i(\varphi _1+\varphi _2\varphi _3)}_c=O\left(\frac{v_1}{N^2}\right),e^{\pm i(\varphi _1+\varphi _2)}_c=O\left(\frac{v_2}{N}\right).$$
(A4)
Note that the second term in the right-hand side of Eq. (A2) is smaller than either the first or the third terms.
Finally, the order of magnitude of the right-hand side of Eq. (A1) is $`v_1^4+O(v_2^2/N^2+1/N^3)`$. We have neglected $`v_1^2/N^2`$ since it is smaller than either $`v_1^4`$ or $`1/N^3`$.
### 2 Calculation of the cumulant $`|Q|^4`$
In this section, we derive the order of magnitude of the fourth order cumulant of the $`|Q|^2`$ distribution, defined by Eq. (29). From the definition of the event flow vector Eq. (22), one obtains
$$|Q|^4=\frac{1}{M^2}\underset{j,k,l,m}{}\left(e^{i(\varphi _j+\varphi _k\varphi _l\varphi _m)}e^{i(\varphi _j\varphi _l)}e^{i(\varphi _k\varphi _m)}e^{i(\varphi _j\varphi _m)}e^{i(\varphi _k\varphi _l)}\right).$$
(A5)
In the above sum, one may distinguish nondiagonal terms, when all four indices are different, and diagonal terms, for which at least two indices are equal.
Nondiagonal terms correspond precisely to the cumulant of the four-particle correlation. The corresponding contribution, evaluated in Sec. A 1, must be multiplied by the combinatorial factor $`M(M1)(M2)(M3)M^4`$. With the factor $`1/M^2`$ in front of the sum in Eq. (A5), the contribution of nondiagonal terms to $`|Q|^4`$ is of order $`M^2v_1^4+O(v_2^2+1/N)`$.
We are now going to show that diagonal terms give a contribution at most of the same order as nondiagonal terms. Let us enumerate the various diagonal terms:
1. If $`j=k=l=m`$, each term in the sum is equal to $`1`$. This is the contribution which we call “autocorrelations”. Multiplying by a combinatorial factor $`M`$ and by the factor $`1/M^2`$ in Eq. (A5), the corresponding contribution is exactly $`1/M`$.
2. When three indices are identical while the fourth is different, i.e. in $`4M(M1)`$ cases, the difference in Eq. (A5) reduces to $`e^{i(\varphi _1\varphi _2)}`$. Using Eq. (8), this contribution is of order $`4v_1^2+O(1/N)`$. Although this contribution is a two-particle correlation, it is suppressed by the combinatorial factor: $`v_1^2`$ is much smaller than the term $`M^2v_1^4`$ which appears in the cumulant of the four-particle azimuthal correlation (see Sec. A 1). Therefore, this contribution will be negligible.
3. Let us consider the cases when the indices are equal two by two.
* If $`j=k`$ and $`l=m`$ but $`jl`$, which occurs $`M(M1)`$ times, the difference is given by
$$e^{2i(\varphi _1\varphi _3)}2e^{i(\varphi _1\varphi _3)}^2=v_2^2+e^{2i(\varphi _1\varphi _3)}_c2\left(v_1^2+e^{i(\varphi _1\varphi _3)}_c\right)^2.$$
(A6)
The order of magnitude is then $`v_2^2+O(1/N)`$. Here, we have neglected terms of order $`v_1^2/N`$ and $`1/N^2`$, smaller than $`1/N`$; the term $`v_1^4`$ is smaller by a combinatorial factor $`1/M^2`$ than the similar contribution of nondiagonal terms. Note that the higher harmonic $`v_2`$ contributes here. We shall see below that these higher harmonics can limit the use of our method.
* The $`2M(M1)`$ cases {$`j=m`$ and $`k=l`$ but $`jl`$} or {$`j=l`$ and $`k=m`$ but $`kl`$} yield a contribution $`e^{i(\varphi _1\varphi _3)}^2`$. Its order of magnitude is $`2v_1^4+O(1/N^2)`$, negligible compared to nondiagonal terms.
4. There are two cases when three indices are different:
* If $`j=l`$ or $`j=m`$ or $`k=l`$ or $`k=m`$, while the two remaining indices are different, the contribution is $`e^{i(\varphi _1\varphi _3)}^2`$, to be multiplied by a combinatorial factor $`4M(M1)(M2)`$. Thus, the order of magnitude is $`M[4v_1^4+O(1/N^2)]`$ and this contribution is suppressed by a factor $`1/M`$ with respect to the cumulant of the four-particle correlation.
* If the two identical indices are either $`(j,k)`$ or $`(l,m)`$, the combinatorial factor is $`2M(M1)(M2)`$, which multiplies a term $`e^{\pm i(2\varphi _1\varphi _3\varphi _4)}2e^{i(\varphi _1\varphi _3)}^2`$. Using Eq. (9), the three-particle correlation $`e^{i(2\varphi _1\varphi _3\varphi _4)}`$ can be expanded as
$`e^{i(2\varphi _1\varphi _3\varphi _4)}`$ $`=`$ $`e^{2i\varphi _1}e^{i\varphi _3}e^{i\varphi _4}+e^{i(2\varphi _1\varphi _3)}_ce^{i\varphi _4}+e^{i(2\varphi _1\varphi _4)}_ce^{i\varphi _3}`$ (A8)
$`+e^{2i\varphi _1}e^{i(\varphi _3+\varphi _4)}_c+e^{i(2\varphi _1\varphi _3\varphi _4)}_c`$
$`=`$ $`v_2v_1^2+O\left({\displaystyle \frac{v_1^2}{N}}\right)+O\left({\displaystyle \frac{v_2^2}{N}}\right)+O\left({\displaystyle \frac{1}{N^3}}\right).`$ (A9)
The second term in the difference, $`2e^{i(\varphi _1\varphi _3)}^2`$, gives a contribution of $`2v_1^4+O(1/N^2)`$. Finally, since terms such as $`v_1^4`$, $`v_1^2/N`$ are suppressed because of the combinatorial factor, the contribution in this case is $`2Mv_2v_1^2+O(Mv_2^2/N)+O(M/N^2)`$.
We shall assume that the total multiplicity in the collision $`N`$ and the number $`M`$ of particles used to calculated the flow vector are large and of the same order of magnitude. Then, we find that the contribution of the diagonal terms is $`1/M+v_2^2+2Mv_1^2v_2+O(v_2^2+1/M)`$.
All in all, when we add the contributions of diagonal and nondiagonal terms, we obtain the following result:
$$|Q|^4=\frac{1}{M}M^2v_1^4+2Mv_1^2v_2+v_2^2+O(v_2^2)+O\left(\frac{1}{M}\right).$$
(A10)
The first term in the right-hand side coresponds to autocorrelations, the last two terms are due to nonflow correlations, and the three remaining terms arise from flow. One would like $`M^2v_1^4`$ to be the dominant flow term. However, higher harmonics, i.e. $`v_2`$, also contribute. If $`v_2`$ is large enough, it may even reverse the sign of the contribution of flow to $`|Q|^4`$. This does not happen provided $`v_2`$ lies in the following interval :
$$Mv_1^2(\sqrt{2}+1)<v_2<Mv_1^2(\sqrt{2}1).$$
(A11)
We have checked these bounds with our Monte-Carlo simulation, see Sec. V C.
## B A generating equation for the integrated flow
In this Appendix, we first construct the cumulants of the distribution of $`|Q|`$, which we note $`|Q|^{2k}`$, as a function of the measured moments $`|Q|^{2k}`$ (Secs. B 1 and B 2). Then, we relate the cumulants to the integrated flow $`Q`$ (Sec. B 3), and show how to remove autocorrelations at all orders (Sec. B 4).
### 1 Cluster decomposition of the moments
We have shown in Sec. II C how the $`k`$-particle momentum distribution can be decomposed, in a coordinate frame where the reaction plane is fixed, into a sum of terms involving lower order distributions ($`k^{}`$ particles with $`k^{}<k`$), plus a “connected” term of relative order $`1/N^{k1}`$. This decomposition also applies to the moments of the distribution of the event flow vector $`Q`$ defined by Eq. (22). As pointed out in Sec. III B, moments of order $`k`$ involve $`k`$-particle azimuthal correlations. This allows us to write a series of equations similar to Eqs. (6) and (9):
$`Q`$ $`=`$ $`Q_c`$ (B2)
$`Q^2`$ $`=`$ $`Q_c^2+Q^2_c`$ (B3)
$`QQ^{}`$ $`=`$ $`Q_cQ^{}_c+QQ^{}_c`$ (B4)
$`Q^3`$ $`=`$ $`Q_c^3+3Q_cQ^2_c+Q^3_c`$ (B5)
$`Q^2Q^{}`$ $`=`$ $`Q_c^2Q^{}_c+2Q_cQQ^{}_c+Q^2_cQ^{}_c+Q^2Q^{}_c,\mathrm{etc}.`$ (B6)
In these equations, the subscript $`c`$ denotes “connected” moments. The connected moment of order $`k`$ is of magnitude $`M^{1k/2}`$: a factor $`M^{1k}`$ comes from the fact that it involves direct $`k`$-particle correlations (see Sec. II C), and a factor $`M^{k/2}`$ from the scaling of $`Q`$ with the number of particles like $`\sqrt{M}`$, see Eq. (22).
The expansion of a given moment $`Q^kQ^l`$ in connected parts can be represented graphically by the expansion of a $`(k+l)`$-point diagram into connected diagrams. This is similar to the decomposition of the $`k`$-particle distribution in Figs. 1 and 2. To be more specific, the decomposition of $`Q^kQ^l`$ is represented by drawing $`k`$ dots of one type corresponding to powers of $`Q`$ and $`l`$ dots of another type corresponding to powers of $`Q^{}`$. One then takes all possible partitions of this set of $`k+l`$ points. To each subset of points one associates the corresponding connected moment. The contribution of a given partition is the product of the contributions of each subset. Finally, $`Q^kQ^l`$ is the sum of the contributions of all partitions. Figure 7 represents, as an example, the decomposition of $`Q^2Q^{}`$.
The connected moments can be expressed as a function of the moments by inverting Eqs. (B 1) order by order. However, this procedure is very tedious. An elegant and compact way to express moments of arbitrary order in terms of the connected parts, and to invert these relations, consists in using generating functions. The generating function of the moments is a function of the complex variable $`z`$ which is defined as
$$𝒢_0(z)=e^{z^{}Q+zQ^{}}=\underset{k,l}{}\frac{z^kz^l}{k!l!}Q^kQ^l,$$
(B7)
where $`k`$ and $`l`$ go from $`0`$ to $`+\mathrm{}`$, and the brackets denote an average over many events. It is well known in graph theory that the generating function of connected diagrams is the logarithm of the generating function of all diagrams . Therefore, the generating function of the connected moments is the logarithm of the generating function of the moments :
$$\underset{k,l}{}\frac{z^kz^l}{k!l!}Q^kQ^l_c=\mathrm{ln}𝒢_0(z).$$
(B8)
The normalization coefficient $`1/k!l!`$ has been chosen such that $`Q^kQ^l_c`$ appears with a unit coefficient in the expansion of $`Q^kQ^l`$, as in Eq. (B 1). Expanding Eqs. (B7) and (B8) to order $`z^2z`$, one finds for instance
$`Q^2Q^{}_c`$ $`=`$ $`Q^2Q^{}Q^2Q^{}2QQQ^{}+2Q^2Q^{},`$ (B9)
which can be checked by inverting Eqs. (B 1) order by order. Note that we are working in a coordinate system where the reaction plane correponds to the $`x`$-axis, and is unknown. In this coordinate system, the generating function (B7) is not a measurable quantity.
### 2 Isotropic source
We now consider specifically an isotropic source, i.e. without flow. In that case, the moment $`Q^kQ^l`$ vanishes if $`kl`$. The connected parts $`Q^kQ^l_c`$ enjoy the same property. Therefore, in the diagrammatic expansion, one only retains terms containing as many powers of $`Q`$ as of $`Q^{}`$, i.e. as many dots on the left as on the right. The quantity represented in Fig. 7 does not satisfy this property, and therefore it vanishes. A decomposition with nonvanishing terms is represented in Fig. 3.
Keeping only the terms $`k=l`$, the generating function (B7) becomes
$$𝒢_0(z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{|z|^{2k}}{(k!)^2}|Q|^{2k}=I_0(2|zQ|),$$
(B10)
where $`I_0`$ is the modified Bessel function of order $`0`$. Note that now the generating function $`𝒢_0`$ itself is isotropic, since $`𝒢_0(z)=𝒢_0(ze^{i\alpha })`$. The consequence is that it can be evaluated in the laboratory coordinate system rather than in the coordinate system associated with the reaction plane: it thus becomes a measurable quantity. We define the cumulants through
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{|z|^{2k}}{(k!)^2}|Q|^{2k}\mathrm{ln}𝒢_0(z)=\mathrm{ln}I_0(2|zQ|).$$
(B11)
They coincide with the connected moments $`|Q|^{2k}_c`$ defined in Eq. (B8) if the source is isotropic. Note that for an isotropic system, the raw moment $`|Q|^{2k}`$ is of order unity, as noted in Sec. III B. The corresponding cumulant $`|Q|^{2k}`$ is of order $`M^{1k}`$. Eq. (B11) corresponds to Eq.(32), where we have set $`x=|z|`$.
### 3 Flow
Let us now calculate the cumulants in the case of collisions with flow. Neglecting for simplicity nonflow correlations between particles, we can write $`Q^kQ^l=Q^kQ^{}^l=Q^{k+l}`$. The generating function (B7) thus becomes
$$𝒢_0(z)=e^{(z+z^{})Q}.$$
(B12)
Now, we want to compare with the experimental value of $`𝒢_0(z)`$, which is measured in the laboratory coordinate system where the azimuth of the reaction plane $`\varphi _R0`$. The generating function in this coordinate system is deduced from Eq. (B12) by the substitution $`zze^{i\varphi _R}`$. Averaging the new expression over all possible $`\varphi _R`$, under the assumption that the distribution of $`\varphi _R`$ is uniform, one obtains:
$$𝒢_0(z)=\frac{1}{2\pi }_0^{2\pi }e^{(ze^{i\varphi _R}+z^{}e^{i\varphi _R})Q}d\varphi _R=I_0(2|z|Q).$$
(B13)
Gathering the results obtained in Eqs. (B11) and (B13), we obtain:
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{|z|^{2k}}{(k!)^2}|Q|^{2k}=\mathrm{ln}𝒢_0(z)=\mathrm{ln}I_0(2|z|Q).$$
(B14)
Expanding Eq. (B14) to order $`|z|^{2k}`$, one obtains an equation relating $`Q^{2k}`$ to the cumulant $`|Q|^{2k}`$. However, when writing Eq. (B12), we have neglected direct $`2k`$-particle correlations and autocorrelations. As explained in Sec. B 1, both give a contribution of magnitude $`M^{1k}`$ to the cumulant $`|Q|^{2k}`$. Thus, Eq. (B14) at order $`|z|^{2k}`$ is valid up to a corrrection of order $`M^{1k}`$.
### 4 Removing autocorrelations
Equation (B14) can be somewhat refined. In the case of a $`Q`$ vector defined with unit weights, as in Eq. (22), autocorrelations can be calculated and subtracted explicitly, which is the purpose of this section.
This calculation has already been done in Sec. III for the lowest orders $`k=1`$ and $`k=2`$: we have seen in Eq. (25) that diagonal terms give a contribution 1 in the expansion of $`|Q|^2`$. In this paper, we refer to these diagonal terms as “autocorrelations”. Similarly, they give a contribution $`1/M`$ to the fourth order cumulant $`|Q|^4`$, see Eq. (34) and Appendix A.
To calculate the contribution of autocorrelations to the cumulant at an arbitrary order, we once again make use of the generating function $`𝒢_0(z)`$, Eq. (B7). Neglecting correlations for simplicity, the contributions of the $`M`$ particles to $`𝒢_0(z)`$ factorize, leading to
$$𝒢_0(z)=e^{(2x\mathrm{cos}\varphi +2y\mathrm{sin}\varphi )/\sqrt{M}}_\varphi ^M,$$
(B15)
where we have set $`z=x+iy`$, and the brackets here denote an average over $`\varphi `$. Assuming for simplicity that the $`\varphi `$ distribution is isotropic, one obtains
$$𝒢_0(z)=\left[I_0\left(\frac{2|z|}{\sqrt{M}}\right)\right]^M.$$
(B16)
This is the expression of the generating function if there are only autocorrelations (no direct correlations, no flow). If there is flow, we assume that autocorrelations and flow give additive contributions to the cumulants, which yields instead of Eq. (B14):
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{|z|^{2k}}{(k!)^2}|Q|^{2k}=\mathrm{ln}𝒢_0(z)=\mathrm{ln}I_0(2|z|Q)+M\mathrm{ln}I_0\left(\frac{2|z|}{\sqrt{M}}\right).$$
(B17)
This formula is equivalent to Eq. (35), which we use in Sec. III B. It removes exactly all autocorrelations when the event flow vector $`Q`$ is defined with unit weights, as in Eq. (22).
## C A generating equation for differential flow
In this appendix, we follow closely the same procedure as in Appendix B, applied to differential flow. In Secs. C 1 and C 2, we first construct the relevant cumulants $`|Q|^{2k}Q^le^{im\psi }`$, as a function of the measured moments $`Q^kQ^le^{im\psi }`$. Here, $`\psi `$ denotes the azimuthal angle of the particle under study (which we call a proton), and $`m`$ the order of the harmonic measured for this particle. Then, we relate the cumulants to the integrated flow $`v_m^{}`$ (Sec. C 3), and show how to remove autocorrelations (Sec. C 4).
### 1 Cluster decomposition
A quantity such as $`Q^kQ^le^{im\psi }`$ involves correlations between $`k+l+1`$ particles: $`k+l`$ “pions” (according to the terminology introduced in Sec. IV) and a proton. This quantity can therefore be decomposed, in the coordinate system where the reaction plane is fixed, into a sum of terms involving lower order correlations, plus a connected term of relative order $`1/N^{k+l}`$. For instance, we can write
$`e^{im\psi }`$ $`=`$ $`e^{im\psi }_c`$ (C2)
$`Qe^{im\psi }`$ $`=`$ $`Q_ce^{im\psi }_c+Qe^{im\psi }_c`$ (C3)
$`QQ^{}e^{im\psi }`$ $`=`$ $`Q_cQ^{}_ce^{im\psi }_c+QQ^{}_ce^{im\psi }_c+Q_cQ^{}e^{im\psi }_c+Q^{}_cQe^{im\psi }_c+QQ^{}e^{im\psi }_c,`$ (C4)
where, in the third equation, the last term is of order $`1/N^2`$ relative to the first one. Such decompositions can be represented diagrammatically, in a way similar to the decomposition of $`Q^kQ^l`$ in Appendix B. We choose to represent the proton by $`m`$ crosses on the left, for reasons which will become clear below, when we consider the specific case of an isotropic source. For instance, Eq. (C4) can be represented diagrammatically by Fig. 8.
In order to express in a compact way the relations between the moments $`Q^kQ^le^{im\psi }`$ and the corresponding connected moments $`Q^kQ^le^{im\psi }_c`$, we introduce the following generating function
$$𝒢_m(z)=e^{z^{}Q+zQ^{}}e^{im\psi }=\underset{k,l}{}\frac{z^kz^l}{k!l!}Q^kQ^le^{im\psi },$$
(C5)
Expanding $`𝒢_m(z)`$ to order $`z^kz^l`$, one obtains all the moments $`Q^kQ^le^{im\psi }`$. In order to obtain the generating function of the connected moments, we note that each diagram in Fig. 8 can be written as the product of a connected diagram containing the crosses, i.e. the proton, times an arbitrary diagram (not necessarily connected) involving only pions, which corresponds to the terms $`Q^kQ^l`$ considered in Appendix B. For instance, using Eqs. (B2) and (B4), one can rewrite Eq. (C4) as
$$QQ^{}e^{im\psi }=QQ^{}e^{im\psi }_c+QQ^{}e^{im\psi }_c+Q^{}Qe^{im\psi }_c+QQ^{}e^{im\psi }_c.$$
(C6)
Therefore, the generating function of the diagrams with pions and protons $`e^{z^{}Q+zQ^{}+im\psi }`$ is the product of the generating function of graphs with only pions, i.e. $`𝒢_0(z)`$ defined in Eq. (B7), by the generating function of connected graphs with pions and protons. This latter is therefore:
$$𝒞_m(z)=\underset{k,l}{}\frac{z^kz^l}{k!l!}Q^kQ^le^{im\psi }_c\frac{𝒢_m(z)}{𝒢_0(z)}=\frac{e^{z^{}Q+zQ^{}+im\psi }}{e^{z^{}Q+zQ^{}}}$$
(C7)
As in Eq. (B8), the normalization coefficient $`1/k!l!`$ has been chosen so that $`Q^kQ^le^{im\psi }_c`$ appears with a unit coefficient in the expansion of $`Q^kQ^le^{im\psi }`$.
### 2 Isotropic source
We now consider the particular case of an isotropic source, without flow. Then the moment $`Q^kQ^le^{im\psi }`$ vanishes when $`k+ml`$, and so do the corresponding connected parts. This is the reason why we chose to represent the proton with $`m`$ crosses: in the isotropic case, only diagrams with the same number of points (crosses and dots) on each side of the dashed line do not vanish.
Expanding the generating function (C5) and keeping only the nonvanishing terms, one finds
$`𝒢_m(z)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{|z|^{2k}z^m}{k!(k+m)!}}|Q|^{2k}Q_{}^{}{}_{}{}^{m}e^{im\psi }`$ (C8)
$`=`$ $`I_m(2|zQ|)\left({\displaystyle \frac{zQ^{}}{|zQ|}}\right)^me^{im\psi },`$ (C9)
where $`I_m`$ is the modified Bessel function of order $`m`$.
We define the cumulants $`|Q|^{2k}Q_{}^{}{}_{}{}^{m}e^{im\psi }`$ so that they coincide with the connected moments in Eq. (C7) when the source is isotropic. Using Eq. (B10), this gives
$`𝒞_m(z)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{|z|^{2k}z^m}{k!(k+m)!}}|Q|^{2k}Q_{}^{}{}_{}{}^{m}e^{im\psi }`$ (C10)
$``$ $`{\displaystyle \frac{I_m(2|zQ|)\left(\frac{zQ^{}}{|zQ|}\right)^me^{im\psi }}{I_0(2|zQ|)}}.`$ (C11)
This equation is equivalent to Eq. (80), setting $`x=|z|`$.
### 3 Flow
Finally, we turn to the more general case of collisions with flow. Neglecting for simplicity nonflow correlations between particles, the generating function (C5) becomes
$$𝒢_m(z)=e^{(z+z^{})Q}v_m^{}.$$
(C12)
As explained in Sec. B 3, this quantity is measured in the laboratory coordinate system, therefore one must replace $`z`$ by $`ze^{i\varphi _R}`$ and average the new expression over all possible $`\varphi _R`$. That yields
$`𝒢_m(z)`$ $`=`$ $`v_m^{}{\displaystyle _0^{2\pi }}e^{(ze^{i\varphi _R}+z^{}e^{i\varphi _R})Q}e^{im\varphi _R}{\displaystyle \frac{\mathrm{d}\varphi _R}{2\pi }}`$ (C13)
$`=`$ $`I_m(2|z|Q)\left({\displaystyle \frac{z}{|z|}}\right)^mv_m^{}.`$ (C14)
Using Eq. (B13), the generating function of cumulants (C7) takes the form
$$𝒞_m(z)=\frac{I_m(2|z|Q)}{I_0(2|z|Q)}\left(\frac{z}{|z|}\right)^mv_m^{}.$$
(C15)
Gathering Eqs. (C11) and (C15), we obtain
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{|z|^{2k}z^m}{k!(k+m)!}|Q|^{2k}Q_{}^{}{}_{}{}^{m}e^{im\psi }=𝒞_m(z)=\frac{I_m(2|z|Q)}{I_0(2|z|Q)}\left(\frac{z}{|z|}\right)^mv_m^{}.$$
(C16)
Expanding this equation to order $`|z|^{2k}z^m`$, one obtains a proportionality relation between the cumulant $`|Q|^{2k}Q_{}^{}{}_{}{}^{m}e^{im\psi }`$ and $`Q^{2k+m}v_m^{}`$. Having measured independently the integrated flow $`Q`$, one thus obtains the differential flow $`v_m^{}`$ from the cumulant. As discussed in Sec. IV, the corresponding error from nonflow correlations is of order $`(Q\sqrt{M})^{k(m/2)}`$.
### 4 Removing autocorrelations
In the case when the “proton” is included in the construction of the event flow vector $`Q_n`$, i.e. if $`\psi `$ is one of the angles $`\varphi _j`$ in Eq. (22), the resulting autocorrelations can be removed at the level of the generating function $`𝒞_m(z)`$ in Eq. (C11): this subtraction is similar to that performed in Sec. B 4 for the integrated flow.
Neglecting correlations for simplicity, the generating function of the cumulants, defined by Eq. (C7), becomes
$$𝒞_m(z)=\frac{e^{(2x\mathrm{cos}\psi +2y\mathrm{sin}\psi +im\psi )/\sqrt{M}}_\psi }{e^{(2x\mathrm{cos}\psi +2y\mathrm{sin}\psi )/\sqrt{M}}_\psi },$$
(C17)
where we have set $`z=x+iy`$, and the brackets denote an average over $`\psi `$. Assuming for simplicity that the $`\psi `$ distribution is isotropic, one obtains
$$𝒞_m(z)=\frac{I_m\left(2|z|/\sqrt{M}\right)}{I_0\left(2|z|/\sqrt{M}\right)}\left(\frac{z}{|z|}\right)^m.$$
(C18)
This is the value of the generating function if there are only autocorrelations. If there is flow in addition, we assume that the contributions of autocorrelations and flow are additive. Equation (C16) is then replaced by
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{|z|^{2k}z^m}{k!(k+m)!}|Q|^{2k}Q_{}^{}{}_{}{}^{m}e^{im\psi }=\left(\frac{I_m(2|z|Q)}{I_0(2|z|Q)}v_m^{}+\frac{I_m\left(2|z|/\sqrt{M}\right)}{I_0\left(2|z|/\sqrt{M}\right)}\right)\left(\frac{z}{|z|}\right)^m.$$
(C19)
This equation is equivalent to Eq. (81), setting $`x=|z|`$. This formula removes exactly all autocorrelations when the vector $`Q_n`$ is defined with unit weights.
## D Interpolation formulas
In this Appendix, we give interpolation methods to calculate numerically the cumulants from their generating functions.
### 1 Integrated flow
The cumulants used for the measurement of the integrated flow are defined by Eq. (99). In order to compute numerically the cumulants $`|Q^{2k^{}}|`$ for $`k^{}=1\mathrm{}k`$ from the generating function, one can for instance tabulate the generating function at the following points:
$$G_{p,q}\mathrm{log}𝒢_0(r_0\sqrt{p}\mathrm{cos}\frac{2q\pi }{q_{\mathrm{max}}},r_0\sqrt{p}\mathrm{sin}\frac{2q\pi }{q_{\mathrm{max}}})$$
(D1)
for $`p=1,\mathrm{},k`$ and $`q=0,\mathrm{},q_{\mathrm{max}}1`$. In this equation, $`r_0`$ is a real number which should be chosen small enough for the series expansion to converge rapidly, typically $`r_00.1`$, and $`q_{\mathrm{max}}`$ is the number of angles at which the generating function is evaluated, which should satisfy the condition $`q_{\mathrm{max}}>2k`$.
One then averages over the angle, thereby eliminating nonisotropic terms up to order $`|z|^{2k}`$:
$$G_p\frac{1}{q_{\mathrm{max}}}\underset{q=0}{\overset{q_{\mathrm{max}}1}{}}G_{p,q}.$$
(D2)
Then, the $`G_p`$, with $`p=1,\mathrm{},k`$, are related to the cumulants $`|\overline{Q}|^{2k^{}}`$ with $`k^{}=1,\mathrm{},k`$ by the following linear system of equations:
$$G_p=\underset{k^{}=1}{\overset{k}{}}|\overline{Q}|^{2k^{}}\frac{r_0^{2k^{}}}{(k^{}!)^2}p^k^{}1pk.$$
(D3)
For practical purposes, it is enough to take $`k=3`$, as explained in Sec. III D. In this case, the solution of the above system reads
$`|\overline{Q}|^2`$ $`=`$ $`{\displaystyle \frac{1}{r_0^2}}\left(3G_1{\displaystyle \frac{3}{2}}G_2+{\displaystyle \frac{1}{3}}G_3\right),`$ (D4)
$`|\overline{Q}|^4`$ $`=`$ $`{\displaystyle \frac{2}{r_0^4}}\left(5G_1+4G_2G_3\right),`$ (D5)
$`|\overline{Q}|^6`$ $`=`$ $`{\displaystyle \frac{6}{r_0^6}}\left(3G_13G_2+G_3\right).`$ (D6)
### 2 Differential flow
The cumulants used for the measurement of the harmonic $`v_m^{}`$ are defined from the generating function by Eq. (109). In order to compute numerically the cumulants $`|Q|^{2k^{}}Q^me^{im\psi }`$ for $`k^{}=0,\mathrm{},k`$, from the generating function, we first tabulate the real and imaginary parts of the generating function, defined by Eq. (107), at the following points:
$`X_{p,q}\mathrm{}\left[𝒞_m(r_0\sqrt{p}\mathrm{cos}{\displaystyle \frac{2q\pi }{q_{\mathrm{max}}}},r_0\sqrt{p}\mathrm{sin}{\displaystyle \frac{2q\pi }{q_{\mathrm{max}}}})\right]`$ (D7)
$`Y_{p.q}\mathrm{}\left[𝒞_m(r_0\sqrt{p}\mathrm{cos}{\displaystyle \frac{2q\pi }{q_{\mathrm{max}}}},r_0\sqrt{p}\mathrm{sin}{\displaystyle \frac{2q\pi }{q_{\mathrm{max}}}})\right]`$ (D8)
for $`p=1,\mathrm{},k+1`$ and $`q=0,\mathrm{},q_{\mathrm{max}}1`$. The number of angles $`q_{\mathrm{max}}`$ must satisfy the condition $`q_{\mathrm{max}}>2(k+m)`$, as we see below.
One then multiplies $`𝒞_m(z)`$ by $`z^m`$, takes the real part and averages over azimuthal angles. Provided $`q_{\mathrm{max}}`$ is large enough, one thus eliminates all nonisotropic terms up to order $`z^kz^{k+m}`$ in the generating function:
$$C_p\frac{\left(r_0\sqrt{p}\right)^m}{q_{\mathrm{max}}}\underset{q=0}{\overset{q_{\mathrm{max}}1}{}}\left[\mathrm{cos}\left(m\frac{2q\pi }{q_{\mathrm{max}}}\right)X_{p,q}+\mathrm{sin}\left(m\frac{2q\pi }{q_{\mathrm{max}}}\right)Y_{p,q}\right].$$
(D9)
Then, the values of $`C_p`$ for $`p=1,\mathrm{}k+1`$ are related to the cumulants $`|Q|^{2k^{}}Q^me^{im\psi }`$ for $`k^{}=0,\mathrm{},k`$ by the following linear system of equations:
$$C_p=\underset{k^{}=0}{\overset{k}{}}|\overline{Q}|^{2k^{}}\overline{Q}^me^{im\overline{\psi }}\frac{r_0^{2(k^{}+m)}p^{k^{}+m}}{k^{}!(k^{}+m)!},1pk+1.$$
(D10)
Taking $`k=1`$ is sufficient for most purposes, as shown in Sec. IV C. For $`m=1`$, the solution of this system is
$`\overline{Q}^{}e^{i\overline{\psi }}`$ $`=`$ $`{\displaystyle \frac{1}{r_0^2}}\left(2C_1{\displaystyle \frac{1}{2}}C_2\right),`$ (D11)
$`|\overline{Q}|^2\overline{Q}^{}e^{i\overline{\psi }}`$ $`=`$ $`{\displaystyle \frac{1}{r_0^4}}\left(2C_1+C_2\right),`$ (D12)
while for $`m=2`$,
$`\overline{Q}^2e^{2i\overline{\psi }}`$ $`=`$ $`{\displaystyle \frac{1}{r_0^4}}\left(4C_1{\displaystyle \frac{1}{2}}C_2\right),`$ (D13)
$`|\overline{Q}|^2\overline{Q}^2e^{2i\overline{\psi }}`$ $`=`$ $`{\displaystyle \frac{1}{r_0^6}}\left(6C_1+{\displaystyle \frac{3}{2}}C_2\right).`$ (D14)
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# Measuring anisotropic scattering in the cuprates
## I Introduction
The unusual transport properties in the normal state of the cuprates continue to attract widespread interest. Many aspects differ significantly from transport in conventional metals iye\_1992a . The resistivity is linear in temperature gurvitch\_1987a with scattering which is apparently electronic rather than electron-phonon in origin bonn\_1993a . The Hall effect shows strong temperature dependence iye\_1992a . On doping with zinc, both the resistivity and inverse Hall angle show a rather simple Matthiessen’s rule chien\_1991a , suggestive of additive scattering mechanisms. However, their temperature dependences differ significantly. This has led to the suggestion that the resistivity and the inverse Hall angle are controlled by very different scattering rates (the so-called ‘two-lifetime’ behavior) anderson\_1991a . This view seems to be confirmed by optical measurements kaplan\_1996a which suggest two independent scattering rates in these materials. The magnetoresistance does not follow Kohler scaling behavior but appears also to be controlled by the Hall angle scattering rate harris\_1995a . Attempts to understand this behavior may be divided into three categories.
Anderson anderson\_1991a has argued that a non-Fermi liquid description must be invoked to understand transport in the cuprates. In Anderson’s picture the electron decays into holons and spinons which separately control resistivity and magneto-transport respectively. A phenomenological transport equation was introduced by Coleman, Schofield and Tsvelik schofield\_1996a to capture a model of transport where the two lifetimes are controlled by independent fluids of particles.
A second approach argues that the magnetic field plays a special role in the transport process. In the model of Kotliar, Sengupta and Varma kotliar\_1996a a singular skew scattering term is invoked. In the picture by Lee and Lee lee\_1997a only recombined slave-fermion and slave-boson particles can interact with the true applied magnetic field. The slave particles themselves are insensitive to the true field because of the large fluctuations in the fictitious gauge field. In both of these scenarios the measured cyclotron frequency would appear to be temperature dependent—a result at odds with current optical Hall experiments drew\_1999a .
The third possibility—the ‘anisotropic scattering’ scenario—envisages a conventional metallic state characterized by an electron quasiparticle scattering rate which depends strongly on the particle momentum. This picture was first suggested by Cooper carrington\_1992a and has been further elaborated upon by Stojković and Pines stojkovic\_1996a , and also by Yakovenko and coworkers yakovenko\_1998a ; yakovenko\_1999a . Pines et al. argue that antiferromagnetic fluctuations in the normal state strongly scatter regions on the Fermi surface which are connected by the antiferromagnetic wave vector. Regions away from these ‘hot-spots’ are weakly scattered. However, Hlubina and Rice hlubina\_1995a have claimed that the cold regions would tend to ‘short-circuit’ the hot spots, thereby leading to conventional transport at low temperatures. Ioffe and Millis ioffe\_1998a have introduced a new variant on the anisotropic scattering picture which is partly motivated by the short-circuiting problem and also inspired by the direct measurements of electron lifetimes. They use angle-resolved photoemission data to argue that quasiparticles are very strongly scattered over most of the Fermi surface. The only long-lived quasiparticles are found at the zone diagonals, where they decay with a weak $`T^2`$ Fermi liquid form. This is known as the ‘cold spot’ model and has been used to understand a wide variety of experiments (see for example Van der Marel vanderMarel\_1999a and Xiang and Hardy xiang\_2000a ).
In this paper we consider this third scenario of anisotropic scattering and treat a minimal model which we believe is illustrative of the strengths and weaknesses of this approach (similar in spirit to some treatments of the cuprate superconducting state maki\_1995a ). Whilst more sophisticated treatments exist, either with a specific relaxation mechanism stojkovic\_1996a ; rosch\_1999a or in the limit of extremely anisotropic in-plane scattering ioffe\_1998a , there are advantages to this model’s simplicity. We can obtain analytic expressions for a wide range of transport properties over the full range of anisotropy. Our calculations illustrate the conclusions of Ioffe and Millis in the limit of strong anisotropy, but also allow us to see how this evolves from the more familiar isotropic metal. This is important as the approach to the isotropic limit may be more relevant to experiments on overdoped materials where the two-lifetime behavior is less apparent. We also consider the transition to the intermediate and high field regimes of magnetoresistance within this picture, and non-linear effects at high electric field.
Perhaps the key discrepancy between the ‘cold-spot’ model and the transport properties of the cuprates is the in-plane orbital magnetoresistance. The model predicts a large magnetoresistance with a distinct temperature dependence in contrast with current experiments. This was was noted in Ioffe and Millis’ original paper ioffe\_1998a and is reproduced in passing here. It has been argued that proper inclusion of vertex corrections could account for this discrepancy narikiyo\_2000a ; kontani\_1999a . The most important new work in this paper is to consider the orbital magnetoresistance for currents moving along the $`c`$-axis. The $`c`$-axis conductivity in a quasi-two dimensional metal may be computed without vertex corrections millis\_2000a and so should be a robust feature of this scenario. We show how an in-plane magnetic field affects the $`c`$-axis conductivity and argue that this experimental geometry provides a key test of the model. Indeed experiments exist on Tl-2201 hussey\_1996a but only in the overdoped regime where the two lifetimes become less distinct. Nevertheless our complete solution of an anisotropic scattering model allows us to obtain a reasonable fit to the data provided we include both scattering anisotropy and bandstructure effects.
Our main conclusions are that within our model the in-plane magnetoresistance remains an outstanding problem if this model is to fit the experiment—not just from its magnitude and temperature dependence but from the scale at which deviations from the weak-field regime occur. While a proper treatment of vertex corrections may correct this for in-plane properties, the $`c`$-axis magnetoresistance is insensitive to these corrections and provides a robust test of the model. Experiments here have mainly been done on overdoped cuprates and our detailed fit to experiments are consistent with this model. High field $`c`$-axis measurements on optimally doped cuprates should be made to test unambiguously this transport phenomenology.
This paper is planned as follows. In section II we present our models of scattering and bandstructure which are the ingredients for the analytic calculation of transport properties. The results are to be found in Section III. These are displayed in the general case, and also in the limit of large scattering anisotropy. Both in-plane and out-of-plane magnetoresistivity are shown, and compared with experimental observations in section IV. Non-ohmic in-plane conductivity is also calculated. Conclusions are drawn in section V.
## II Anisotropic transport
Our main assumption in this work is to use a Boltzmann equation within the relaxation time approximation as a description of transport. We write the collision integral, $`I[f]`$, as $`\mathrm{\Gamma }_𝐤[f(𝐤)f^{(0)}(𝐤)]`$ and impose a phenomenological scattering rate, $`\mathrm{\Gamma }_𝐤`$, which varies smoothly around the in-plane Fermi surface and has square symmetry.
We do not speculate on the underlying cause of the quasiparticle relaxation rate but, if this scattering were due to scattering from a soft bosonic mode at finite $`𝐐`$ (antiferromagnetic spin fluctuations for example), then the validity of the relaxation time approximation might be questionable rosch\_2000a . This is because such scattering would tend to equilibrate quasiparticles only between points on the Fermi surface connected by the spanning vector ($`𝐐`$) of the mode. By contrast, the relaxation time approximation assumes that the quasiparticle is equilibrated uniformly around the Fermi surface. This will not affect our calculation of c-axis response if the bosonic mode is 2-dimensional, but it could change the in-plane physics. This is presumably the reason why vertex corrections could be important for in-plane transport. However, including this process correctly would tend to make the quasiparticle distribution even more anisotropic when driven out of equilibrium. As we will see, the anisotropy of the out-of-equilibrium distribution of quasiparticles is precisely what is probed by the in-plane magnetoresistance. Increasing this anisotropy is likely to increase further the in-plane magnetoresistance.
Within the relaxation time approximation we now introduce two forms of scattering rate to model the anisotropy.
### II.1 The cold-spot model
Following Ioffe and Millis ioffe\_1998a , the angle-resolved photoemission results suggest that there are points on the Fermi surface where the electron is a fairly long-lived excitation: points where the Fermi surface intersects the Brillouin zone diagonal. Elsewhere on the Fermi surface, only broad features in energy are seen in the photoemission. Within a Fermi liquid picture one would then infer that electron-like quasiparticles are strongly scattered everywhere except for these special points. We can capture this by the following form for our scattering rate
$$\mathrm{\Gamma }(k_{},\theta ,k_{})=\mathrm{\Gamma }_f\mathrm{cos}^22\theta +\mathrm{\Gamma }_s\mathrm{sin}^22\theta .$$
(1)
Here $`\theta `$ is the azimuthal angle between the in-plane momentum wavevector $`k_{}`$ and the $`a`$-axis direction of the crystal, as shown in Fig. 1.
This is a simple model which nevertheless allows us to explore the transition from isotropic ($`\mathrm{\Gamma }_f=\mathrm{\Gamma }_s`$) to extremely anisotropic scattering ($`\mathrm{\Gamma }_f\mathrm{\Gamma }_s`$). The latter limit corresponds to the case studied by Ioffe and Millis, where, at the planar zone diagonals \[$`\theta =(2n+1)\pi /4`$\] the scattering rate has a minimum. In their model, $`\mathrm{\Gamma }_s`$ is just the Fermi-liquid scattering rate, $`\mathrm{\Gamma }_s=1/\tau _{FL}T^2`$. The scattering around the Fermi-surface in the limit of strong anisotropy is illustrated in Fig. 2.
It will often be convenient to express the scattering rate in terms of an anisotropy parameter, $`\alpha `$, where Eq. 1 is written as
$$\mathrm{\Gamma }(k_{},\theta ,k_{})=\mathrm{\Gamma }_0\left(1+\alpha \mathrm{cos}4\theta \right),$$
(2)
with
$`\mathrm{\Gamma }_0`$ $`=`$ $`(\mathrm{\Gamma }_f+\mathrm{\Gamma }_s)/2,`$ (3)
$`\alpha `$ $`=`$ $`{\displaystyle \frac{(\mathrm{\Gamma }_f\mathrm{\Gamma }_s)}{(\mathrm{\Gamma }_f+\mathrm{\Gamma }_s)}}.`$ (4)
In this form $`\alpha =0`$ gives us the isotropic scattering case whilst $`\alpha =1`$ or $`1`$ gives large anisotropy with cold spots on the zone diagonals or zone axes respectively. We will express results in either notation ($`\mathrm{\Gamma }_f,\mathrm{\Gamma }_s`$ or $`\alpha ,\mathrm{\Gamma }_0`$) depending on applicability to available data and ease of interpretation, and are always able to calculate across the full range of $`\alpha `$.
### II.2 The hot-spot model
Alternatively it is has been proposed that hot spots exist near the $`(\pi ,0)`$ points of the Brillouin zone where strong scattering occurs, perhaps due to antiferromagnetic spin fluctuations. This leaves quasiparticles on the rest of the Fermi surface with a longer lifetime. To capture this within our phenomenology we add the scattering lifetimes around the Fermi surface so
$$\mathrm{\Gamma }_𝐤^1=\tau (k_{},\theta ,k_{})=\tau _f\mathrm{cos}^22\theta +\tau _s\mathrm{sin}^22\theta .$$
(5)
Here we expect $`\tau _f`$ to be smaller than $`\tau _s`$ so that we have a model with strong-scattering hot spots and extended cold regions with less scattering (see Fig. 2).
### II.3 Bandstructure
The next component of the model is the bandstructure. We assume that the in-plane dispersion is isotropic, which is not unreasonable for the cuprates. By contrast the dispersion in the out-of-plane direction is known to have significant dependence on the in-plane momentum xiang\_1996a ; andersen\_1994a ; liechtenstein\_1996a . This is due to the local chemistry of the copper-oxide planes as emphasized by Xiang xiang\_1998a . Transport between planes occurs principally through the copper 4s orbitals. Since these orbitals are circularly symmetric in the plane, their overlap amplitude with the Cu 3d$`_{x^2y^2}`$–O 2p hybrids has d-wave symmetry as illustrated in Fig. 3.
The resulting tunneling term is therefore proportional to the square of this amplitude and can be represented by the form $`t_{}\mathrm{cos}^2(2\theta )`$. The most important effect this has on $`c`$-axis transport stems from the fact that the longest-lived quasiparticles on the zone diagonals are precisely those that have a vanishing probability of $`c`$-axis hopping. For this reason we ignore other features of the $`c`$-axis dispersion which may lead to additional small tunneling probabilities at other points on the Fermi-surface such as those due to a body-centered tetragonal unit cell considered elsewhere dragulescu\_1999a ; vanderMarel\_1999a . Our bandstructure is then given by
$$ϵ(𝐤_{},\theta ,𝐤_{})=ϵ(k_{}^2)2t_{}\left(1+\gamma \mathrm{cos}4\theta \right)\mathrm{cos}(k_{}c).$$
(6)
To illustrate the importance of this effect and to allow for other tunneling mechanisms we parameterize the degree of tunneling anisotropy by $`\gamma `$, which we allow to vary from 0 to 1. When $`\gamma `$=1, we have complete elimination of c-axis transport for quasiparticles with in-plane momentum along the zone diagonals. The frequency dependence of this end-point of the model has been studied by van der Marel vanderMarel\_1999a . Allowing both the degree of scattering anisotropy ($`\alpha `$) and $`c`$-axis dispersion ($`\gamma `$) to vary across their full range turns out to be crucial in understanding the out-of-plane transport and is a novel aspect of this work. It unifies various other approaches stojkovic\_1996a ; ioffe\_1998a ; hussey\_1996a which focus on the extremes of one or other of the ranges of $`\alpha `$ and $`\gamma `$.
## III Results
Using the standard representation of the quasiparticle distribution, $`f=f^{(0)}+\psi _ϵf^{(0)}`$, we write the Boltzmann equation (for uniform fields and temperature gradients) as abrikosov\_1980a
$`_t\psi +\mathrm{\Gamma }_𝐤\psi {\displaystyle \frac{e}{\mathrm{}}}(𝐄+𝐯_𝐤\times 𝐁)_𝐤\psi `$ $`=`$
$`e𝐄𝐯_𝐤`$ $``$ $`{\displaystyle \frac{ϵ_𝐤\mu }{T}}𝐯_𝐤T.`$ (7)
We will solve this equation in various limits in the subsequent sections. We will only be interested in d.c. properties in this paper, however our results for the cold-spot model (Eq. 1) are easily extended into the frequency domain. After Fourier transforming with respect to time, we see that
$`_t\psi +\mathrm{\Gamma }(\theta )\psi `$ $``$ $`i\omega \psi +(\mathrm{\Gamma }_f\mathrm{cos}^22\theta +\mathrm{\Gamma }_s\mathrm{sin}^22\theta ),`$
$``$ $`(\mathrm{\Gamma }_f+i\omega )\mathrm{sin}^22\theta +(\mathrm{\Gamma }_s+i\omega )\mathrm{cos}^22\theta .`$
We can therefore obtain the finite frequency results for the cold-spot model from our expressions at d.c. simply by the prescription
$$\mathrm{\Gamma }_f\mathrm{\Gamma }_f+i\omega ,\mathrm{\Gamma }_s\mathrm{\Gamma }_s+i\omega .$$
(9)
### III.1 In-plane transport in weak magnetic fields
We first focus on leading order response in an in-plane electric field. For a circular Fermi surface we may write Eq. 7 as
$$\mathrm{\Gamma }(\theta )\psi \omega _c\frac{\psi }{\theta }=eEv_F\mathrm{cos}\theta ,$$
(10)
where $`\omega _c=eBv_F/(\mathrm{}k_F)`$ and $`v_F`$ is the in-plane Fermi velocity.
We may solve this equation by Jones-Zener expansion for in-plane currents with a weak magnetic field along the $`c`$-axis. The currents, and hence conductivity elements are calculated using the relation
$$𝐣_\mu =\frac{eE}{4\pi ^3\mathrm{}}\frac{𝐯_\mu k_F}{v_F}\psi 𝑑\theta 𝑑k_{},$$
(11)
where current is flowing in direction $`\mu `$ with velocity $`𝐯_\mu `$. We expand $`\sigma _{xx}`$ and $`\sigma _{xy}`$ to order $`B^2`$ and the resistivity is then found by simple matrix inversion.
If we consider first hot-spot scattering (Eq. 5), we can see that this scattering rate can be eliminated as a viable model for the cuprates. Each of the properties evaluated is dominated by the long scattering time, which is the short-circuiting effect proposed by Hlubina and Rice hlubina\_1995a . The in-plane conductivities are given by
$`\sigma _{xx}^{(0)}`$ $`=`$ $`{\displaystyle \frac{e^2k_Fv_F(\tau _f+\tau _s)}{4\pi \mathrm{}c}},`$ (12)
$`{\displaystyle \frac{\sigma _{xy}}{\sigma _{xx}^{(0)}}}`$ $`=`$ $`\mathrm{tan}\mathrm{\Theta }_H=\omega _c{\displaystyle \frac{(3\tau _f^2+2\tau _f\tau _s+3\tau _s^2)}{4(\tau _f+\tau _s)}}+O[B^3],`$
$`{\displaystyle \frac{\mathrm{\Delta }\sigma _{xx}}{\sigma _{xx}^{(0)}}}`$ $`=`$ $`\omega _c^2{\displaystyle \frac{(21\tau _f^234\tau _f\tau _s+21\tau _s^2)}{8}}+O[B^4].`$ (14)
Within this model, all computed quantities are dominated by $`\tau _s`$ in the limit of strong anisotropy ($`\tau _s\tau _f`$). This model is unable to reproduce the two-lifetime phenomenology of the cuprates so we will not consider it further.
However, for the cold spot model we obtain the following in terms of $`\mathrm{\Gamma }_f`$ and $`\mathrm{\Gamma }_s`$:
$`\sigma _{xx}^{(0)}`$ $`=`$ $`{\displaystyle \frac{e^2v_Fk_F}{2\pi \mathrm{}c}}{\displaystyle \frac{1}{\sqrt{\mathrm{\Gamma }_f\mathrm{\Gamma }_s}}},`$ (15)
$`{\displaystyle \frac{\sigma _{xy}}{\sigma _{xx}^{(0)}}}`$ $`=`$ $`\mathrm{tan}\mathrm{\Theta }_H={\displaystyle \frac{\omega _c}{2}}\left({\displaystyle \frac{1}{\mathrm{\Gamma }_f}}+{\displaystyle \frac{1}{\mathrm{\Gamma }_s}}\right)+O[B^3],`$ (16)
$`{\displaystyle \frac{\mathrm{\Delta }\sigma _{xx}}{\sigma _{xx}^{(0)}}}`$ $`=`$ $`\omega _c^2\left[{\displaystyle \frac{5}{8}}\left({\displaystyle \frac{\mathrm{\Gamma }_f}{\mathrm{\Gamma }_s^3}}+{\displaystyle \frac{\mathrm{\Gamma }_s}{\mathrm{\Gamma }_f^3}}\right){\displaystyle \frac{1}{8}}\left({\displaystyle \frac{1}{\mathrm{\Gamma }_f^2}}+{\displaystyle \frac{1}{\mathrm{\Gamma }_s^2}}\right)\right]+O[B^4].`$
This results in an in-plane orbital weak-field magnetoresistance
$$\frac{\mathrm{\Delta }\rho _{xx}}{\rho _{xx}^{(0)}}=\omega _c^2\frac{(\mathrm{\Gamma }_f\mathrm{\Gamma }_s)^2}{8\mathrm{\Gamma }_f^3\mathrm{\Gamma }_s^3}(5\mathrm{\Gamma }_f^2+7\mathrm{\Gamma }_f\mathrm{\Gamma }_s+5\mathrm{\Gamma }_s^2),$$
(18)
which vanishes in the isotropic limit ($`\mathrm{\Gamma }_s=\mathrm{\Gamma }_f`$). This reflects the well known result that there is no orbital magnetoresistance in an isotropic metal. In the high anisotropy ($`\alpha 1`$) limit, where $`\mathrm{\Gamma }_f\mathrm{\Gamma }_s`$, we find the results shown in table 1.
These results illustrate the way in which the cold-spot model reproduces the phenomenology of the cuprates as shown by Ioffe and Millis ioffe\_1998a . They take $`\mathrm{\Gamma }_f`$ to be large and temperature independent, while $`\mathrm{\Gamma }_s`$ is assumed small and proportional to $`T^2`$. The geometric mean then gives the linear in temperature resistivity while the inverse Hall angle is proportional to $`\mathrm{\Gamma }_s`$ and hence varies as $`T^2`$. The results are also illustrative of the problems of this model. Adding Zn impurities to the $`\mathrm{CuO}_2`$ has been interpreted as adding a unitary scatterer which should add a temperature independent term to the scattering rates. Adding a constant elastic scattering term to the two rates appearing in the geometric mean which forms the resistivity (Eq. 15) will not reproduce the Matthiessen’s rule behavior seen in experiment chien\_1991a . Instead it will change both the temperature dependence and the intercept of the resistivity.
The second difficulty is the large magnetoresistance that this model predicts. While experiments suggest that the magnetoresistance is proportional to the Hall angle squared harris\_1995a with a constant of proportionality which can be of order 1, the cold-spot model predicts that these two quantities differ by $`\mathrm{\Gamma }_f/\mathrm{\Gamma }_s`$—a large, temperature dependent factor. It is tempting to appeal to new physical processes to explain this (such as vertex corrections). However the cause of this large magnetoresistance is intimately related to the appearance of distinct temperature dependences of the resistivity and inverse Hall angle. Magnetic fields couple to the anisotropy of $`\psi `$ around the Fermi surface as can be seen from Eq. 10 and so at each order in the magnetic field, the effect of the anisotropy in the model is amplified. This allows the inverse Hall angle to differ from the resistivity but consequently has an even more dramatic effect on the magnetoresistance which is higher order in magnetic field. We do not attempt to offer a solution for this puzzle though given the systematics in transport behavior across a wide range of cuprates it seems unlikely that this is due to fine tuning. (By contrast in the charge-conjugation phenomenology schofield\_1996a the magnetoresistance is proportional to the square of the Hall angle and is small since this model is isotropic around the Fermi surface.)
### III.2 Beyond the weak-field regime
Thus far, in keeping with other workers, we have shown that in-plane transport is governed by a different scattering rate (formed from the combination of $`\mathrm{\Gamma }_f`$ and $`\mathrm{\Gamma }_s`$) at each order in the magnetic field. Here we go beyond the weak-field regime and show numerically that no new combinations emerge beyond second order in the magnetic field until the magnetoresistance saturates. The scattering rate combination for weak-field magnetoresistance provides the scale for magnetic field effects throughout the intermediate field region.
The saturation field may be obtained from a high field expansion of the Boltzmann equation—expanding the solution in powers of $`1/\omega _c`$. We find that the magnetoresistance saturates at a value of
$$\frac{\mathrm{\Delta }\rho _{xx}(B\mathrm{})}{\rho _{xx}(B=0)}=\frac{\left(\sqrt{\mathrm{\Gamma }_f}\sqrt{\mathrm{\Gamma }_s}\right)^2}{2\sqrt{\mathrm{\Gamma }_s\mathrm{\Gamma }_f}}.$$
(19)
The high field limit is reached when $`\omega _c\mathrm{\Gamma }_f`$ so the fastest rate controls the transition to the high field limit. This is because saturation requires that quasiparticles can completely orbit the Fermi surface and it is the hot regions which limit this processtyler\_1998a . This is also the scale for Landau quantization. The approach to the saturated value is shown in Fig. 4. The magnetic field is plotted as a function of $`\omega _c/\mathrm{\Gamma }_0`$ since in the cold spot model this is dominated by $`\mathrm{\Gamma }_f`$ and so shows that this sets the scale for saturation.
\[As an aside, we can obtain the universal form of the magnetoresistance in an isotropic metal. Of course, this is formally zero but by dividing by the saturated value in the limit of small anisotropy we can obtain a finite result. This is also shown in Fig. 4.\]
While we cannot solve Eq. 10 analytically except in the low and high field limits, the equation is readily solved numerically. The fact that we require a solution of $`\psi `$ which is periodic around the Fermi surface means that we can look for a Fourier series solution. Our numerical results reproduce the analytical results of the previous section. In addition, we find the remarkable result that the magnetoresistance is a function of a single dimensionless parameter until it saturates. This scaling parameter governs the magnetoresistance even beyond the weak-field ($`B^2`$) regime of the magnetoresistance so we may write
$$\frac{\mathrm{\Delta }\rho _{xx}(B)}{\rho _{xx}(0)}=f(\omega _c\tau _{\mathrm{mr}}),\omega _c\mathrm{\Gamma }_f,$$
(20)
where $`\tau _{\mathrm{mr}}`$ is defined by the weak-field magnetoresistance such that $`\mathrm{\Delta }\rho _{xx}/\rho _{xx}=(\omega _c\tau _{\mathrm{mr}})^2`$. For $`\mathrm{\Gamma }_f\mathrm{\Gamma }_s`$ we have $`\tau _{\mathrm{mr}}^25\mathrm{\Gamma }_f/8\mathrm{\Gamma }_s^3`$.
We illustrate this universality in the inset to Fig. 4. We see that as a function of this combination of scattering rates, the magnetoresistance follows a universal curve until the saturation value is approached. An intermediate field region emerges for fields $`1<\omega _c\tau _{\mathrm{mr}}<\omega _c/\mathrm{\Gamma }_f`$, where the magnetoresistance adopts a new power law $`B^{0.33}`$. Thus our numerical treatment suggests that no new scattering rate combinations are produced in the Jones-Zener expansion until the magnetoresistance saturates. We have confirmed this is the case analytically at fourth order in the magnetic field. This scaling also means that in our model there is an absolute value of magnetoresistance above which one must see a deviation from a $`B^2`$ dependence. Our model predicts that with a magnetoresistance of only 0.5% one would expect at least a 10% deviation from a $`B^2`$ form. Comparing this with experimental work tyler\_1998a on optimally doped Tl-2201 in pulsed magnets, we find a deviation from quadratic dependence is only perhaps being seen in optimally doped Tl-2201 at a magnetoresistances of order 3%. (The observed deviation from the weak-field region is more consistent with the Coleman-Schofield-Tsvelik phenomenology which predicts that this should occur when the magnetoresistance is at or less than about 10% tyler\_1998a .)
### III.3 $`c`$-axis transport
In-plane magnetoresistance in the cold-spot model is at odds with the current experiments. It has been argued that including vertex corrections could correct this narikiyo\_2000a ; kontani\_1999a . As we have argued in the introduction, $`c`$-axis transport is then of fundamental importance in testing the cold-spot model. This is because to lowest order in $`t_{}^2`$ the $`c`$-axis conductivity is related to the convolution of two in-plane spectral functions mckenzie\_1999a . These are the quantities which may be inferred from the angle-resolved photoemission that inspired the cold-spot model. Thus, independent of a picture of quasiparticles for in-plane transport, a cold-spot model should be able to account for the $`c`$-axis d.c. conductivity millis\_2000a . In particular we must show that this leads to a consistent picture of $`c`$-axis orbital magnetoresistance (magnetic field in the plane and electric field out of the plane). The most systematic study of these effects has been performed by Hussey et al. hussey\_1996a in overdoped Tl-2201 and this motivates the following analysis. The two-lifetime behavior is much less apparent in the overdoped material so we will not be able to use the $`\mathrm{\Gamma }_f\mathrm{\Gamma }_s`$ asymptotics. Instead we need a full solution for any degree of anisotropy.
We will now calculate the $`c`$-axis magnetoresistance, $`\mathrm{\Delta }\rho _{zz}/\rho _{zz}`$. To obtain this, we first calculate the relevant components of the conductivity matrix, $`\sigma _{ij}`$ and invert. In this experimental geometry, $`\sigma `$ has zero components $`\sigma _{xy}`$,$`\sigma _{xz}`$,$`\sigma _{yx}`$ and $`\sigma _{zx}`$ due to the $`𝐁`$ field being in-plane. $`\sigma _{xx}`$ is equal to the zero-magnetic field value of $`\sigma _{yy}`$. We can expand the other terms in magnetic field:
$$\sigma _{\nu \mu }=\sigma _{\nu \mu }^{(0)}+\underset{n=1}{\overset{n=3}{}}\sigma _{\nu \mu }^{(n)},$$
(21)
where the superscripts refer to the order of effect in magnetic field. Symmetry under time reversal means, of course, that $`\sigma _{zz}^{(1)}`$ and $`\sigma _{zz}^{(3)}`$ and the Hall terms $`\sigma _{yz}^{(0)}`$ and $`\sigma _{yz}^{(2)}`$ are zero. In addition the Hall term $`\sigma _{yz}^{(1)}`$ is small ($`t_{}^2`$). Under these conditions, the $`c`$-axis magnetoresistivity in weak magnetic field simplifies to
$`\rho _{zz}`$ $`=`$ $`\rho _{zz}^{(0)}+\mathrm{\Delta }\rho _{zz}^{(2)}\mathrm{\Delta }\rho _{zz}^{(4)}+O[B^6],`$
$`{\displaystyle \frac{\mathrm{\Delta }\rho _{zz}^{(2)}}{\rho _{zz}^{(0)}}}`$ $`=`$ $`{\displaystyle \frac{\sigma _{zz}^{(2)}}{\sigma _{zz}^{(0)}}},`$
$`{\displaystyle \frac{\mathrm{\Delta }\rho _{zz}^{(4)}}{\rho _{zz}^{(0)}}}`$ $`=`$ $`{\displaystyle \frac{(\sigma _{zz}^{(2)}\sigma _{zz}^{(0)}\sigma _{zz}^{(4)})}{\sigma _{zz}^{(0)\mathrm{\hspace{0.17em}2}}}}.`$ (22)
Here we have adopted the sign convention of Hussey et al. hussey\_1996a and Drǎgulescu et al. dragulescu\_1999a , expecting the fourth-order magnetoresistivity to be negative.
To obtain the conductivity, we follow the same general procedure as for the in-plane case, solving the Boltzmann equation using a Jones-Zener expansion. We now introduce the angle, $`\varphi `$, which is the in-plane angle between the B field and the $`a`$-direction. $`\theta `$ becomes the azimuthal coordinate relative to the direction of the B field. The solutions to the Boltzmann equation are now
$$\psi (𝐤)=\frac{eE}{\mathrm{}\mathrm{\Gamma }(\theta +\varphi )}\frac{ϵ}{k_{}}+\underset{n=1}{\overset{\mathrm{}}{}}\psi ^{(n)}(𝐤),$$
(23)
where, to lowest order in $`t_{}`$,
$$\psi ^{(n)}(𝐤)=\frac{eB^n}{\mathrm{}\mathrm{\Gamma }(\theta +\varphi )}\left(\frac{1}{\mathrm{}}v_F\mathrm{sin}\theta \frac{\psi ^{(n1)}}{k_{}}\right).$$
(24)
The ratio of $`\sigma _{zz}^{(2)}/\sigma _{zz}^{(0)}`$, when expanded in powers of $`\alpha `$ yields a zero-order term which is, as Hussey et al. hussey\_1996a first showed
$$\frac{\sigma _{zz}^{(2)}}{\sigma _{zz}^{(0)}}=\frac{c^2e^2v_F^2}{2\mathrm{}^2\mathrm{\Gamma }_0^2}+O[\alpha ^2].$$
(25)
Similarly, for the quadratic term in the $`c`$-axis conductivity:
$$\frac{\sigma _{zz}^{(4)}}{\sigma _{zz}^{(0)}}/\left(\frac{\sigma _{zz}^{(2)}}{\sigma _{zz}^{(0)}}\right)^2=\frac{6+3\gamma ^2+2\gamma \mathrm{cos}(4\varphi )}{2(2+\gamma ^2)},$$
(26)
which in the limit of simple inter-plane tunneling ($`\gamma 0`$) gives 3/2, again as shown by Hussey et al.
The Hall conductivity, $`\sigma _{yz}`$, is found to vary as $`t_{}^2`$. In the limit of isotropic in-plane scattering the out-of-plane Hall angle is identical to the in-plane one schofield\_2000a
$$\text{lim}_{\alpha 0}\frac{\sigma _{yz}^{(1)}}{\sigma _{zz}^{(0)}}=\frac{eBv_F}{\mathrm{}k_F\mathrm{\Gamma }_0}.$$
(27)
This is not true in general however. Nevertheless its dependence on $`t_{}^2`$ means that the Hall conductivity does not appear, at leading order, in the out-of-plane magnetoresistance.
Rather than display the full expression for every quantity we have calculated, we show graphically two quantities which Hussey et al. examined experimentally in single-crystal Tl<sub>2</sub>Ba<sub>2</sub>CuO<sub>6</sub>. There a key observation was the four-fold variation of the $`c`$-axis magnetoresistance as the magnetic field was rotated in-plane, i.e. that
$$\mathrm{\Delta }\rho _{zz}^{(4)}=\overline{\rho }_{zz}^{(4)}+\stackrel{~}{\rho }_{zz}^{(4)}\mathrm{cos}(4\varphi ),$$
(28)
where $`\overline{\rho }_{zz}^{(4)}`$ and $`\stackrel{~}{\rho }_{zz}^{(4)}`$ were both found to be positive. In fact it is straightforward to show that this 4-fold modulation at fourth order in B is a simple consequence of square symmetry. Hussey et al. initially analyzed their angle-dependent magnetoresistance results purely in terms of anisotropic, in-plane scattering (as parameterized here by $`\alpha `$). We first look at the offset part in the fourth order magnetic field term in the magnetoresistance. We examine $`\frac{\overline{\rho }_{zz}^{(4)}/\rho _{zz}^{(0)}}{(\mathrm{\Delta }\rho ^{(2)}/\rho _{zz}^{(0)})^2}`$, a quantity found experimentally to be roughly independent of temperature and approximately equal to 0.6 $`(0.6\pm 0.1)`$. We show it for various values of $`\gamma `$, as a function of $`\alpha `$ in Figure 5. This puts a constraint on the temperature dependence of $`\alpha `$, an issue which is explored in section IV where we try to fit magnetoresistance data to our model.
The second quantity studied experimentally is the amplitude of the modulation in magnetoresistance at order $`B^4`$. The quantity evaluated is $`\stackrel{~}{\rho }_{zz}^{(4)}/\overline{\rho }_{zz}^{(4)}`$, which is positive when the resistance is maximum with $`𝐁`$ aligned along the zone diagonal (Eq. 28). This too is plotted across the range of scattering anisotropy, $`\alpha `$, for various bandstructures, $`\gamma `$, in Figure 6.
Hussey et al. found it to vary with temperature, dropping from about 0.15 at 50K to 0.05 at about 140K, and these limits are shown as dotted lines in Figure 6. In our calculations, $`\stackrel{~}{\rho }_{zz}^{(4)}/\overline{\rho }_{zz}^{(4)}`$ is positive when $`\alpha `$ is negative, corresponding to a situation where the hot spots and cold spots have effectively exchanged their usual positions. This can be explained intuitively by identifying the maxima in the Lorentz force, $`𝐯\times 𝐁`$, at points on the Fermi surface orthogonal to the in-plane magnetic field. Then, loosely, the parts of the Fermi surface least affected by a magnetic field are those with $`𝐯_F`$ parallel to $`𝐁`$. Thus, naively, one would expect the minima in the magnetoresistance to occur when the field is along the cold-spot directions. In fact experimentally this orientation gives a maxima in the magnetoresistance. Turning on anisotropy in the $`c`$-axis dispersion $`\gamma `$, one can sufficiently inhibit c-axis transport through the cold spots in the positive $`\alpha `$ regime that resistance maxima can be seen for $`𝐁`$ along with cold spots along in keeping with the experiments.
This effect has been attributed elsewhere dragulescu\_1999a to the curvature of the real Fermi surface in a model of isotropic scattering. Here we have shown that the two effects will work together in the regime of positive bandstructure anisotropy. Once again, however, we see that there is a constraint on the allowed degree of variation of scattering anisotropy with temperature. In Table 2 we show the form of these 4th order quantities shown in Figures 5 and 6 as well as the size of the 2nd order magnetoresistivity and in-plane properties. This will be useful when fitting to data later.
### III.4 Non-ohmic conductivity
We have also calculated the non-ohmic, in-plane conductivity of the system in zero magnetic field. Hlubina hlubina\_1998a has suggested that this is a useful test of the model of Ioffe and Millis. The Boltzmann equation (Eq. 7) is now solved by Jones-Zener expansion in $`𝐄`$ and is given by
$$\mathrm{\Gamma }_𝐤\psi \frac{e}{\mathrm{}}𝐄_𝐤\psi =e𝐄𝐯_𝐤.$$
(29)
From the 4-fold symmetry of the system, it is straightforward to deduce that the currents in the $`x`$ and $`y`$ directions must have the following form at third order in the electric field:
$`j_x`$ $`=`$ $`\sigma _0E_x+\sigma _1E_x^3+\sigma _2E_y^2E_x,`$ (30)
$`j_y`$ $`=`$ $`\sigma _0E_y+\sigma _1E_y^3+\sigma _2E_x^2E_y.`$ (31)
We then find that
$`\sigma _0`$ $`=`$ $`{\displaystyle \frac{e^2k_Fv_F}{2\pi \mathrm{}c\mathrm{\Gamma }_0(1\alpha ^2)^{1/2}}},`$ (32)
$`\sigma _1`$ $`=`$ $`{\displaystyle \frac{e^4v_F}{4\pi \mathrm{}^3k_Fc\mathrm{\Gamma }_0^3}}{\displaystyle \frac{\alpha ((1+\alpha )(\alpha ^2+2\alpha +2))}{(1\alpha ^2)^{7/2}}},`$ (33)
$`\sigma _2`$ $`=`$ $`{\displaystyle \frac{e^4v_F}{4\pi \mathrm{}^3k_Fc\mathrm{\Gamma }_0^3}}{\displaystyle \frac{\alpha ((\alpha ^39\alpha ^2+4\alpha 6)}{(1\alpha ^2)^{7/2}}},`$ (34)
which all simplify or reduce to zero appropriately in the isotropic limit, $`\alpha 0`$. Furthermore, in the limit of strong anisotropy, $`\alpha 1`$, we find $`\sigma _1=\sigma _2`$, as demonstrated by Hlubina hlubina\_1998a .
The most obvious consequence of non-linear response in the electric field is that the current no longer flows parallel to the electric field. Defining the parallel current ($`j_{||}`$) as being the response parallel to the applied field, and the transverse current as the current component perpendicular to the applied field ($`j_{}`$) we find that
$$j_{}=\frac{1}{4}E^3\left(\sigma _1\sigma _2\right)\mathrm{sin}4\varphi ,$$
(35)
where $`\varphi `$ is the angle between the in-plane electric field and the $`a`$-axis. Using our expressions for $`\sigma _1`$ and $`\sigma _2`$ (Eqs. 33 and 34) we may write
$`{\displaystyle \frac{j_{}}{j_{||}}}`$ $`=`$ $`\left({\displaystyle \frac{eE}{\mathrm{}k_F\mathrm{\Gamma }_0}}\right)^2{\displaystyle \frac{\alpha \left(2+\alpha ^3\right)}{2\left(1\alpha ^2\right)^3}}\mathrm{sin}4\varphi ,`$ (36)
$`=`$ $`\underset{\alpha 1}{lim}\left({\displaystyle \frac{v_d}{v_F}}\right)^2{\displaystyle \frac{5}{16}}\left({\displaystyle \frac{\mathrm{\Gamma }_f}{\mathrm{\Gamma }_s}}\right)^3\mathrm{sin}4\varphi .`$ (37)
Here $`v_d=eE/2m\mathrm{\Gamma }_0`$ is the Drude drift velocity of a fast decaying quasiparticle. So, although the $`v_d/v_F`$ is generally tiny, nonlinear effects are indeed strongly enhanced by anisotropy.
### III.5 Thermal transport
In this treatment we have not speculated on the anisotropy of the energy relaxation rate around the Fermi surface. However if we make the assumption that the energy relaxation follows the quasiparticle relaxation rate and only varies around the Fermi surface (as opposed to variations away from the Fermi surface) we can calculate all of the thermal transport coefficients. Without reproducing the details of the calculation we can see immediately from the form of the Boltzmann equation (Eq. 7) that temperature gradients drive the quasiparticle distribution in exactly the same way as electric fields except for the usual factor of $`(ϵ_𝐤\mu _F)/T`$. This is seen in the right-hand-side of Eq. 7. This factor is isotropic around the Fermi surface and does not introduce any further angular dependence. Thus we can conclude that all currents proportional to temperature gradients will have exactly the same dependence on scattering rate as those driven by electric fields (for example $`\sqrt{\mathrm{\Gamma }_f\mathrm{\Gamma }_s}`$ in the absence of a magnetic field). Within this approximation the thermal conductivity will obey the Wiedemann-Franz law and the diffusion thermopower will have the usual linear temperature dependence and be independent of the scattering rate. To account for the unusual systematics in the measured thermopower of the cuprate metals obertelli\_1992a in this model one would need to add a significant anisotropic energy dependence of the scattering rate that differs from the anisotropy of the transport relaxation rate.
## IV Application to experimental data
In order to apply the analytic results of our model to experimental data, we have made a comparison of magnetoresistivity results with measurements on overdoped samples of single-crystal Tl<sub>2</sub>Ba<sub>2</sub>CuO<sub>6</sub>, made by Tyler tyler\_1998b and Hussey hussey\_1996a . We choose this system because its $`c`$-axis magnetoresistance is well characterized. We have argued that $`c`$-axis properties are the most robust quantities in this model since they are not affected by vertex corrections, nor do they rely on a quasiparticle picture. We then identify combinations of measured $`c`$-axis quantities which directly probe the degree of anisotropy with minimal dependence on unknown parameters. Finally we address the in-plane transport features. We also identify combinations of experimental quantities which allow a comparison with theory involving the fewest assumptions on unknown parameters. In particular we can constrain the degree of anisotropy in the $`c`$-axis dispersion ($`\gamma `$) and the range of scattering anisotropy ($`\alpha `$) from the $`c`$-axis magnetoresistance.
First we consider the overall magnitude of the fourth order magnetoresistance by comparing this to the second order term. This is shown in Fig. 5(b). When we compare this to our model \[Fig. 5(a)\] we see that to obtain the observation of an essentially temperature independent result of around 0.6, the degree of scattering anisotropy must be limited to $`\alpha <0.5`$. Furthermore, if $`\alpha `$ is to have any temperature dependence at all then $`\gamma `$ should be approaching 1 where the gradients in Fig. 5(a) are smallest.
We have already remarked on a second conclusion from a comparison between theory and experiment. In order for the maxima in the $`c`$-axis magnetoresistance to occur when the field is along the zone-diagonals, $`c`$-axis transport along these directions must be suppressed. This is illustrated in Fig. 6(a) where we see that with $`\alpha >0`$ we require $`\gamma >0`$ to obtain a ratio of $`\stackrel{~}{\rho }_{zz}^{(4)}/\overline{\rho }_{zz}^{(4)}`$ of the correct sign. Furthermore the experimental bounds on this ratio \[illustrated in Fig. 6(b)\] confirm that the degree of scattering anisotropy ($`\alpha `$) must be less than about 0.6. In no sense then are we in the limit of strong anisotropy in this overdoped material.
A further observation comes from the temperature dependence shown in Fig. 6(b). Generically the quantities shown in Figs. 5(b) and 6(b) should be temperature dependent so the observation that only the second of these has any significant variation with temperature restricts $`\alpha `$ to vary between about 0.2 and 0.4 with a band structure anisotropy $`\gamma 0.9`$. Also, in order to fit the trends in the data, we see that the overdoped state must become less anisotropic as the temperature is lowered. Proximity to a zero temperature critical point would suggest the opposite should be true but perhaps here elastic processes are begining to dominate.
In order to predict other experimental quantities, we then need to know the temperature dependence of $`\mathrm{\Gamma }_0`$. This too can be found from the second order c-axis magnetoresistance, when data is compared with the functional form in table 2, here fixing $`\gamma =1`$. In addition it can also be independently infered from the in-plane resistivity but now relying on the assumptions of Boltzman transport theory.
With these two temperature dependencies known, we may try to predict in-plane properties. We show plots of the cotangent of the Hall angle, $`\mathrm{cot}\mathrm{\Theta }_H`$, and the in-plane magnetoresistance $`\mathrm{\Delta }\rho _{xx}^{(2)}/\rho _{xx}^{(0)}`$, together with experimental data from Hussey and Tyler in Figs. 7 and 8. We see by comparing the experimental data with the parameters extracted from other measurements that the data and theory follow the same trends and are of the correct order of magnitudes. The two different fits on each plot represent the two different methods of infering $`\mathrm{\Gamma }_0`$: either using the in-plane resistivity or the out-of-plane magnetoresistance. The difference in these two methods is that the in-plane measurements would be expected to be sensitive to the transport lifetime, as opposed to the quasiparticle lifetime in the case of out-of-plane measurements. Surprisingly it appears that the quasiparticle lifetime is longer than the transport lifetime.
Overall we see that for overdoped Tl-2201 a consistent picture emerges of a Fermi-liquid-like metal with weak scattering anisotropy around the Fermi surface but a strong anisotropy in the $`c`$-axis dispersion. A more rigourous test of the model would be obtained by performing a similar comparison on optimally doped Tl-2201 ($`T_c=85`$K). There the two-lifetime behavior of in-plane magnetotransport is clearly apparent with very good fits to data being found with simple power laws tyler\_1998b :
$`\rho _{xx}`$ $`=`$ $`20+1.56(T/\mathrm{K})\mu \mathrm{\Omega }\mathrm{cm},`$ (38)
$`(B/\mathrm{T})\mathrm{cot}\theta _H`$ $`=`$ $`300+1.7810^2(T/\mathrm{K})^2,`$ (39)
for $`100\mathrm{K}<T<300`$K. In the limit of high in-plane anisotropy we can combine these measurements to obtain the magnetoresistance. We find
$$\frac{\mathrm{\Delta }\rho _{xx}}{\rho _{xx}^{(0)}}=\frac{5}{2}\left(\frac{e^2k_F^4}{\pi ^2c^2B^2}\right)\left[\frac{\rho _{xx}^{(0)}}{\mathrm{cot}^2\theta _H}\right]^2.$$
(40)
This would predict a magnetoresistance of about 0.04 at 10T and 130K (so already out of the weak-field regime)—a factor of 40 greater than currently seen. We have used in-plane properties to predict the magnetoresistance so uncertainties about vertex corrections are present in this estimate.
Again, using the in-plane properties, we would expect that the out-of-plane magnetoresistance would be 0.1 in fields as low as 3T with a fourth order magnetoresistance of 0.001 in an in-plane field of about 5T. These effects should be measurable but again rely on using in-plane transport lifetimes to provide a measure of in-plane quasiparticle lifetimes. However, we can make an unambiguous prediction within this model: that the positions of the maxima in resistivity at fourth order should move from the directions (seen in the overdoped system) to the directions in the optimally doped materials. This is seen in Fig. 6(a) where for large $`\alpha `$ we see that bandstructure is unable to prevent the sign of the modulation term, $`\stackrel{~}{\rho }_{zz}^{(4)}/\overline{\rho }_{zz}^{(4)}`$, from being negative.
## V Conclusion
We have presented a thorough investigation of a minimal model for transport in a quasi-2D system, where we allow for anisotropic scattering in-plane and an anisotropy in the out-of-plane dispersion, both fully variable. Such a model has been proposed for the normal state of the cuprate superconductors with either strong scattering hot spots and weak scattering cold regions around the Fermi surface or hot regions with cold spots. We have studied the transport properties of both types of models and illustrated how short-circuiting in the hot-spot model makes this inconsistent with transport measurements on the cuprates. For the cold-spot model we have computed in- and out-of-plane magnetoresistivity, in-plane non-ohmic conductivity and thermal conductivities for an arbitrary degree of anisotropy.
We find that the in-plane magnetoresistance in this model is too large when compared with experiments on the cuprates, in keeping with other work ioffe\_1998a . In addition the in-plane magnetoresistance should be universal beyond the weak field limit with a well defined deviation from a $`B^2`$ dependence at weaker fields than currently observed. It has been argued that vertex corrections may account for these discrepancies between the model and experiment. Here we have focussed on the out-of-plane magnetoresistance for which there are no vertex corrections in the quasi-2D limit. We have completely characterized the magnetoresistance to fourth order in this geometry. We have compared our model with experiments on overdoped Tl2201 and show that the model can be reasonably well fit to the experiments. This requires only weak scattering anisotropy but a high degree of bandstructure anisotropy in the $`c`$-axis dispersion.
A better test of this model would be to compare it to optimally doped Tl-2201 where the two-lifetime behavior is very clearly seen in in-plane magnetotransport measurements. This model would predict that strong angle dependent $`c`$-axis magnetoresistance should be observed and the positions of the minima in the optimally doped system should be rotated by $`\pi /4`$ from those of the overdoped material. A quantitative analysis of the degree of in-plane anisotropy could also be made. The extent to which vertex corrections control the in-plane transport properties could then be assessed.
## Acknowledgements
The authors would like to thank C. Bergemann, D. R. Broun, N. E. Hussey, S. R. Julian, D. E. Khmelnitskii, M. W. Long, G. G. Lonzarich, A. J. Millis, A. Rosch and T. Xiang for useful discussions. We are grateful for the hospitality of the Center for Materials Theory, Rutgers University and the Newton Institute, Cambridge where some of this work was done and also for the support of the Theory of Condensed Matter Group in Cambridge. KGS was funded by the EPSRC, AJS by the Royal Society.
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# 1 Introduction and Summary
## 1 Introduction and Summary
In the recent past, there has been much attention to models with large extra dimensions. The surge in activity surrounding this idea owes its origin to the belief that the existence of extra dimensions (beyond four) seems to be a crucial ingredient for the unification of gravity with gauge forces. The initial goal of large radius, $`rM_P^1`$, compactification schemes is to weaken the hierarchy between the electroweak scale, and the four dimensional gravity scale, $`M_P`$. The idea is that the matter content of the Standard Model of elementary particles (SM) is confined to $`(3+1)`$-dimensions, as suggested by , while gravity lives in the whole $`D`$-dimensional space ($`D>4`$). Upon compactification, the hierarchy problem is solved by lowering the fundamental scale of gravity, $`M_{}`$, down to TeV through a model-dependent relation between $`M_{}`$ and $`M_P`$. The compactification mechanisms suggested so far, can be classified into two broad categories: models with tensor product of our four dimensional world with the internal space (in line with the original Kaluza-Klein ideology), and models with warp product between the same.
On compactifying down to four dimensions, one may in general get new degrees of freedom added to the SM spectrum. The new states can be purely from the gravitational sector, or have Standard Model KK excitations in addition (depending on whether the SM interactions are written directly in four dimensions, using the induced metric, or written fully in $`D`$ dimensions). In any case, the new states might lead to detectable modifications of the existing accelerator data and cosmological observations . The phenomenology of both categories has been investigated in great pursuit during the past two years, and we refer to a summary of the recent findings in . The usual way to avoid such new contributions to the prevailing scenario at low energies is often either by decoupling the particles by making their masses very heavy (beyond the present reach of accelerators, say $``$ TeV), or by imposing judicious bounds on their couplings and masses. Recently, it was suggested that the heavy masses could be realized naturally (without fine-tuning), utilizing only certain geometrical properties of the internal manifold, namely that the masses arising from compactification are exponentially large being related to the volume of the internal hyperbolic manifold.
In this work we scrutinize the criteria for chosing the internal manifold, in both cases of tensor and warp product compactifications, based on geometrical and topological arguments, such that the unwanted KK contributions are avoided. We focus on compact, connected, and smooth internal manifolds with scalar curvature, bounded from below $`\kappa (d1)K`$, where $`K`$ is a constant. We consider, for generality, a gravity theory coupled to a Dirac spinor in the presence of a gauge theory. This consideration, has a final aim to be applied to the SM. However, in order to keep the discussion simple and sufficiently general (model independent) we shall not concern ourselves with finer details like localization mechanisms, the issue of obtaining chiral fermions starting from odd dimensions,…etc. Instead, after performing the general analysis of various bounds on the KK masses, one may specialize to the case of the SM.
Most of our analysis relies on the following basic observations concerning the spectrum of Riemannian manifolds, and the Dirac operator on spin manifolds. The main fact is that the spectrum of $`DD\text{/}`$ and that of the Laplacian on compact Riemmanian manifolds is discrete, bounded from below, and the eigenvalues (counted with multiplicity) are ordered: $`0=\lambda _0\lambda _n\lambda _{n+1}`$. Moreover, there exist lower bounds on $`\lambda _1`$ of a Laplacian acting on a scalar in the compact manifold. In addition, there are upper bounds on the eigenvalues which sets a ceiling to how heavy they can be lifted. These translate into lower bounds on the 4-dimensional tree-level masses of particles arising from compactification. For spinors, the classic theorem of Lichnerowicz enables us to impose similar bounds, upper and lower, and altogether exclude tree level massless fermions for certain internal manifolds. To sum up, we use topological considerations to comment on bosonic KK zero modes, while we use geometrical arguments to impose bounds on fermions and massive KK modes.
The paper is organized as follows: in section 2, we summarize our conventions, and state our requirements for choosing the internal manifold so that we are able to produce a phenomenologically reliable scenario in a rather model independent way. The implementation of these demands is carried on the successive sections. In section 3, we relate the eigenvalues of the Laplacian on the internal space with the tree-level masses in four dimensions, and discuss geometrical upper and lower bounds for massive bosons, and topological conditions for the massless ones. In section 4, we comment on massless spinors using Lichnerowicz theorem, and point out the existence of curvature-dependent upper and lower bounds on massive ones. In order to satisfy all of our requirements, stated in the second section, we are able to rule out certain choices of the compactification manifolds. Finally, we summarize our conclusions in section 5.
## 2 Conventions and Set up
As mentioned in the introduction, we consider Einstein’s gravity coupled to a Dirac spinor and a Yang-Mills gauge theory on a $`D`$-dimensional manifold $`W=M_4Y`$,
$$S=\underset{W}{}d^Dx\sqrt{g}\left[\frac{1}{G}+\frac{1}{4}F^2+i\overline{\psi }\widehat{DD\text{/}}_A\psi \right]$$
where $`M_4`$ is a four manifold, which we eventually identify as our 4-dimensional world, and $`Y`$ is a compact $`(D4)`$-dimensional manifold. $`\widehat{DD\text{/}}_A`$ is the twisted Dirac operator on $`W`$ and $`F`$ is the YM field strength (For a detailed treatment of an analogous setting in six and ten dimensions, and conventions, we refer to ). Consequently, the various fields on $`W`$ will be decomposed as the following: scalars on $`W`$ will be scalars on both $`M_4`$ and $`Y`$; a vector on $`W`$ will be a vector on $`M_4`$ and scalar on $`Y`$ or vice versa; the graviton on $`W`$ will appear as a graviton on $`M_4`$ and scalar on $`Y`$ or vice versa, or as a vector on both submanifolds. Finally, a spinor on the parent manifold will decompose as a spinor on both $`M_4`$ and $`Y`$. It is perhaps worth mentioning that a spinor defined on $`W`$ which is a fibre product of $`M_4`$ and $`Y`$, does not necessarily split into spinors defined on the two submanifolds individually, as it does in the Cartesian tensor product case. However, in the special case of a warp product, the fibration being trivial, this decomposition once again holds.
Our analysis includes both the tensor product decomposition (the usual compactification case), and the warp product decomposition . Let us recall that for the tensor product, $`W=M_4Y`$, the inherited metric is $`\widehat{g}=g_4+g_Y`$ where $`g_4`$ and $`g_Y`$ are the metrics on $`M_4`$ and $`Y`$ respectively. Whereas for a warp product, $`W=M_4_^+Y`$, the resultant inherited metric is of the form $`\widehat{g}=f^2g_4+g_Y`$, where $`f`$ is a smooth map $`f:Y^+`$. In this work we choose the warp factor to be $`f=e^{\frac{1}{2}\varphi }`$ as in . The warp factor is to be consistently determined by solving Einstein’s equations.
Our main requirements, for the theory resulting after compactification, can be outlined as the following:
With respect to the gravity sector, we want to end up with one massless graviton, no additional massless gauge bosons, and no massless scalars<sup>3</sup><sup>3</sup>3Massless vectors may enhance the gauge symmetry, and gravitational interactions mediated by scalars violate the equivalence principle..
The zeroeth KK mode(s) of the Dirac spinor is massless in four dimensions. If specialized to the SM, this translates into the requirement that fermion masses are exclusively due to Higgs mechanism.
The masses of KK excitations of various fields are naturally heavy.
## 3 Bounds on bosonic KK masses
The starting point of our analysis is to examine the masses in the gravity sector. Looking at the linearized Einstein’s equations on $`W`$<sup>4</sup><sup>4</sup>4 We use the linearized equations for obvious reasons, nevertheless we include for completeness the decomposition of $``$ on $`M_4`$ and $`Y`$,
$$=_4+\kappa $$
for tensor product; and for warp product
$$=\frac{1}{f^2}\left\{_48f\mathrm{\Delta }f6f^2\right\}+\kappa $$
where $`f`$ is the warp factor, which we may assume to be $`f=e^{\frac{1}{2}\varphi }`$., $`\widehat{\mathrm{\Delta }}h_{\overline{k}\overline{l}}=T_{\overline{k}\overline{l}}`$ ($`\overline{k}`$’s are the indices on $`W`$), one can relate the spectrum of $`\mathrm{\Delta }_Y`$ with the tree-level masses of the various fields on $`M_4`$. The parent Laplacian<sup>5</sup><sup>5</sup>5The “hat” superscript refers to quantities defined on $`W`$., $`\widehat{\mathrm{\Delta }}`$, decomposes as
$$\widehat{\mathrm{\Delta }}=\mathrm{\Delta }_4+\mathrm{\Delta }_Y$$
(1)
in the tensor product case, and as
$$\widehat{\mathrm{\Delta }}=e^{\varphi (y)}\mathrm{\Delta }_4+\mathrm{\Delta }_Y\frac{1}{2}\left(_l\varphi (y)\right)^l$$
(2)
in the warp product case ($`l`$’s are the indices on $`Y`$).
### 3.1 Massless bosons
A necessary condition to meet the first demand ($`i`$) is to select $`Y`$ with the appropriate betti numbers<sup>6</sup><sup>6</sup>6 The number of zero modes of the Laplacian (or equivalently the dimension of the space harmonic $`p`$-forms) on a compact manifold $`Y`$ are given by the $`p`$-th betti numbers $`b_p(Y)`$ of the manifold.. Since $`b_0(Y)=1`$ for any general connected manifold $`Y`$, we are guaranteed to end up with one massless graviton on $`M_4`$. $`b_1(Y)=0`$ would ensure that no new massless vector bosons, as well as massless scalars are produced in $`M_4`$ after compactification. However, this is not the case for a general $`Y`$. For example, a circle has $`b_1(S^1)=1`$ and for a torus, $`T^d`$; $`b_1=d`$, both of which therefore admit massless 1-forms. $`S^d`$s are in general suitable ambient spaces for performing such compactifications, since the have $`b_1(Y)=0`$ for $`d>1`$. Other possible alternatives are Calabi-Yau’s, $`K3`$’s, suitable orbifolds of the type $`T^d/Z_n`$, compact hyperbolic manifolds for $`d3`$. In general, for spaces having $`b_10`$ quotienting by an appropriate discrete isometry, often leaves us with $`b_1=0`$. This topological classification is insufficient when harmonic spinors are discussed to meet the demand ($`ii`$). For, in that case, the curvature of the manifold (a geometric parameter) plays the decisive role.
This analysis of the massless sector applies to both compactification schemes– tensor or warp.
### 3.2 Massive bosons: Lower bounds
Having considered the massless fields, which are the zero modes of the Laplacian on the internal space, we now turn our attention to the first massive excitations. As we mentioned in the introduction, there has been an extensive study for the first non-zero eigenvalue of the Laplacian on Riemannian manifolds. Rigorous bounds, particularly for manifolds with scalar curvature bounded from below by $`(d1)K`$ (where $`K`$ is constant and $`d`$ is the dimension of the manifold), have been established. Any way, assuming a slowly varying $`\kappa `$ makes it possible to replace it by $`(d1)K`$ in the context of discussing mass bounds and scales. It may be noted that, the eigenspectrum being strictly ordered, only the lowest massive states are relevant to our analysis, because if we achieve to decouple these, then all the higher modes will be automatically eliminated in the effective four dimensional theory.
Let $`Y`$ be a compact manifold, and $`\lambda _1`$ the first non-zero eigenvalue of $`\mathrm{\Delta }_Y\varphi _n=\lambda _n\varphi _n`$, where $`\varphi _n`$ is a scalar. Then
$$\lambda _1+\text{max}\{(d1)K,0\}\frac{\pi ^2}{4\sigma ^2}$$
(3)
where $`\sigma `$ is the diameter of the manifold. It is worthwhile to note from (3) that the fundamental parameter for masses arising from compactification is $`\sigma `$, and not generically the volume of the manifold, as it is commonly thought<sup>7</sup><sup>7</sup>7This can be easily understood by observing that it is possible to change the spectrum of the Laplacian by deforming the manifold, and yet keeping its volume fixed. However, the relation between $`M_{}`$ and $`M_P`$ will always involve the volume of $`Y`$.. However, in certain cases one can proceed a step further and relate $`\sigma `$ to the volume of the manifold, and hence rewrite the bounds in terms of the volume instead (e.g. in $`S^d`$ and certain compact hyperbolic manifolds). The inequality (3) translates effectively into a statement about the bounds on the 4-dimensional masses<sup>8</sup><sup>8</sup>8Here, and elsewhere, we use the rest frame when referring to massive states. of the lowest excitations, $`m_1^2`$. It is obvious that when the Ricci curvature, $`\kappa `$, of $`Y`$ is non-negative, then one recovers the standard scenario: $`\lambda _1\pi ^2/(4\sigma ^2)`$, where in the standard KK scenario, as in , $`\sigma `$ is identified with the diameter of the compactification circle(s).
At this point, we note that the explicit expression of the bounds will depend on the nature of the product between the two manifolds– a tensor or a warp product. In the tensor product case, the bound (3) on first scalar excitations (in the 4-dimensional effective theory) will remain unaltered,
$$m_1^2\frac{\pi ^2}{4\sigma ^2}\text{max}\{(d1)K,0\}$$
A natural choice would be $`\sigma ^1M_{}`$ (say $`𝒪(\text{TeV})`$). In the case of $`\kappa `$ bounded from below by a negative constant (i.e. not everywhere positive), the bound will involve the infimum of the curvature (or the curvature itself, if constant or slowly varying), and in order to achieve $`m_1^2\text{TeV}^2`$ we need
$$\kappa |(d1)K|(\frac{\pi ^2}{4}1)\text{TeV}^2$$
(4)
It is remarkable that satisfying this bound on the curvature requires no fine-tuning at all<sup>9</sup><sup>9</sup>9For large extra dimensions models to do anything with string/M-theory, one must have $`d7`$.. As was noticed in , some manifolds with negative scalar curvature, like compact hyperbolic ones, may have attractive features like exponentially large KK masses. We would like to speculate that negatively curved internal spaces may also be favoured (beside the string inspired Ricci flat compactifications), because they support the existence of massless spinors, as will be shown in the next section.
In the warp product case, it is not straightforward to comment at this level. Mainly, because the eigenvalue of $`\mathrm{\Delta }_4`$, $`\stackrel{~}{m}_1^2`$, will be $`y`$-dependent in the $`D`$-dimensional theory,
$$\stackrel{~}{m}_1^2e^{\varphi (y)}\left[\frac{\pi ^2}{4\sigma ^2}\text{max}\{(d1)K,0\}\frac{1}{2}\left(_l\varphi (y)\right)^l\right]$$
(5)
and the integration over the $`y`$ coordinates, in order to get the effective four dimensional mass, becomes non-trivial in the the presence of the extra terms. It is clear from (2) that both the warp factor and the term $`_l\varphi (y)^l`$ (which should be understood as the gradient of the wave function in the internal space) will change the interpretation of the effective four dimensional mass, and hence it is not straightforward to make a statment about the bound. In fact, it is not only the warp factor dependence on $`y`$ which matters here, but also the gradient of the wave function of the field in the internal space. This conclusion is in contrast with the bounds on graviton excitations discussed in .
Whereas the third demand, of the scalar sector in the theory<sup>10</sup><sup>10</sup>10This bound applies on any massless scalar field in the theory, whether or not in the gravity sector., can be met in the tensor product case (by chosing $`Y`$ with an appropriate diameter) it seems difficult to be fulfilled without further model dependent details, in the warp product case.
The same arguments can be carried over to the case of vectors and rank two tensors arising from gravity sector, as a consequence of the linearization procedure. In principle, one could expect curvature dependent additions to the equations of motion, as can be read off from (13), but those additional terms are dropped off because they involve $`𝒪(h^2)`$. In this context, we would like to point out that the vector degrees of freedom, resulting from the metric decomposition, can not in general be eliminated by a gauge choice, and their coupling to a typical scattering amplitude are comparable to those of graviton exchange . Hence, it is important to make them very massive and weakly coupled. On the other hand, the curvature dependent terms in (13) will appear in the YM equations of motion, and it will be difficult to draw conclusions<sup>11</sup><sup>11</sup>11 Unfortunately, no rigorous bounds for Laplacians on 1-forms, relevant to our discussion, have been worked out., apart from the special case when $`d=2`$ where the bounds for the 1-forms are the same as in the case of scalars . In any case, from (13), it can be argued that the bounds for these fields are of the same order as in (3).
### 3.3 Massive bosons: Upper bounds
Finally, we would like to add the following remark. Although the first non-zero eigenvalue of $`\mathrm{\Delta }`$ is bounded from below, it is not possible in general to push it to an infinitely heavy scale. There exists an upper bound which depends on the same parameter $`\sigma `$. For example, if the Ricci curvature $`0`$ then the $`j`$th eigenvalue is bounded from above by
$$\lambda _j\frac{2j^2}{\sigma ^2}d(d+4)$$
(6)
And in the case when Ricci curvature is bounded from below by a negative number ($`K<0`$), then the upper bound becomes more complicated, and includes the (lower bound of) the curvature. For $`d2`$ the bound is
$$\lambda _j\frac{(2l1)^2}{4}K+\frac{4\pi ^2j^2}{\sigma ^2}(1+2^{l1})^2$$
(7)
for $`d=2l`$, $`l=1,2,\mathrm{}`$, and
$$\lambda _jl^2K+\frac{4(1+\pi ^2)j^2}{\sigma ^2}(1+2^{2l2})^2$$
(8)
for $`d=2l+1`$, $`l=1,2,\mathrm{}`$. Concerning 1-forms, the same upper bound holds for $`\lambda _1^{(1)}`$, because $`\lambda _1^{(1)}\lambda _1`$ (without any further assumption concerning the curvature).
In the tensor product case, the bounds (6,7,8), remain unaltered on the effective four dimensional masses. However, a more careful treatment, as discussed in section (3.2), is required in the warp product case.
## 4 Bounds on fermionic KK masses
The Dirac operator on the spin manifold $`W`$ acting on spinors is $`\widehat{DD\text{/}}=\gamma ^{\overline{k}}(X)_{\overline{k}}`$, where $`\gamma ^{\overline{k}}(X)`$’s are the $`D`$-dimensional Gamma matrices in curved space written in terms of the Veilbeins on $`W`$. It decomposes in the tensor product case as
$$\widehat{DD\text{/}}=DD\text{/}_4+DD\text{/}_Y$$
(9)
where $`DD\text{/}_Y=\gamma ^l(y)_l`$, $`DD\text{/}_4=\gamma ^\mu (x)_\mu `$; $`\gamma ^l(y)`$, and $`\gamma ^\mu (x)`$ being the Gamma matricies of the 4 and $`(D4)`$ dimensional spaces respectively ($`\mu `$’s run over $`M_4`$). In the warp product case, the $`DD\text{/}`$ splits differently from the one in (9), and has the form
$$\widehat{DD\text{/}}=e^{\frac{1}{2}\varphi (y)}DD\text{/}_4+DD\text{/}_Y$$
(10)
where $`DD\text{/}_4`$ and $`DD\text{/}_Y`$ are the same operators defined previously. Hence, the four dimensional fermion masses are related to the eigenvalues of the Dirac operator on the internal manifold<sup>12</sup><sup>12</sup>12We define of the spinor (mass$`)^2`$ as an eigenvalue of the squared Dirac operator.. And in particular, the observed massless fermions in four dimensions are nothing but the zero modes of $`DD\text{/}_Y`$ (which lie in $`\text{ker}DD\text{/}_Y`$). It has been shown by Lichnerowicz that not all manifolds admit harmonic (massless) spinors. The argument is based on the relation between the squared Dirac operator and the scalar curvature,
$$DD\text{/}_Y^2=^{}+\frac{1}{4}\kappa $$
(11)
where $`\kappa `$ is the scalar curvature of $`Y`$, $`^{}`$ is the adjoint of $``$, and $`^{}`$ is the connection Laplacian (a positive operator). We use this theorem not only to identify candidate manifolds in which the demand ($`ii`$) can be realized, but also to set geometrical bounds so that ($`iii`$) is met.
### 4.1 Massless fermions
Recall that a spinor $`\psi `$ is said to be harmonic iff $`DD\text{/}\psi =0`$, i.e $`\psi \text{ker}DD\text{/}`$, It is helpful to remember that $`\text{ker}DD\text{/}=\text{ker}DD\text{/}^2`$, and that this space is finite dimensional, and this space is identified with our space of massless fermions, as metioned above. It has been shown in that the existence of harmonic spinors depend strongly on the scalar curvature of the manifold, and in particular, massless spinors do not exist on manifolds with a positive scalar curvature<sup>13</sup><sup>13</sup>13As an example, massless spinors do not exist on a sphere.. This no-go theorem applies also to cases where the scalar curvature is non-negative everywhere, and not necessarily constant. In addition, the formula (11) shows that fermion mass square is bounded from below by the curvature since the operator $`^{}`$ is positive.
According to the above argument, meeting the second demand, namely supporting massless spinors (which will eventually acquire mass only through Higgs mechanism), rules out the entire class of manifolds with positive curvature, unless they have further discrete isometries. If one wants to relax the second demand, by having the above mass term, then a careful attention should be paid in order not to spoil gauge invariance. For instance, a direct mass term in the action for the SM fermions is not gauge invariant. So, adding such a tree-level mass term by hand, and yet keeping gauge invariance, will be at the cost of doubling (or increasing) the number of degrees of freedom (and hence the number of SM generation) depending on the dimension of the spinor representation in the $`D`$-dimensional space.
The above price has to be paid anyway, on either negatively or positively curved manifolds, when one goes beyond $`D=6`$. It can be shown that $`D=6`$ is the maximum dimension possible to end up with a 4-dimensional Weyl spinor (starting from a 6-dimensional Weyl), without extra degress of freedom. Starting from an irreducible spin represenation of $`SO(1,D1)`$, and after some algebra, the resulting number of four dimensional Weyl spinors is: $`(2^{(D5)/2}\times n)`$ for $`D`$ odd (although special care should be taken in order to define spinors in odd dimensions ), and $`(2^{D/23}\times n)`$ for $`D`$ even, where $`n`$ is the number of zero modes of the Dirac operator in the internal space (on a compact manifold, the eigenstates are all square-integrable). This number $`n`$ depends on the coupling of spinors to background fields. Hence, the number of the zeroeth KK fermionic modes will increase, possibly leading to a variant number of flavours.
A common way to get rid of the above extra spinorial degrees of freedom is to use a localization mechanism as first proposed by . These mechanisms rely on the existence of more that one zero mode of the Dirac operator in the internal space, such that at least one of them is not normalizable. Therefore, a necessary condition for applying such scenarios is to have a non-compact internal space, because all the modes of a given Dirac operator on a compact space are normalizable. Hence, the recently discussed mechanisms , break down for the compact manifold, and an extra care is needed for dealing with the additional modes (specially the zero ones). Finally, we mention that the arguments contained here, concerning the zero KK modes, are generic in the sense that they do not depend on whether the product is tensor or warp.
### 4.2 Massive fermions: Lower bounds
As it is the case for the Laplacian, the eigenvalues of the Dirac operator on a compact space are discrete. Therefore, the eigenvalues of the squared Dirac operator are discrete and positive, and in addition any eigenvalue, $`\nu _q^2`$,is bounded from below by the curvature , including $`\nu _1^2`$,
$$\nu _1^2\frac{d}{4(d1)}\lambda _1$$
(12)
where $`\lambda _1`$ is the first eigenvalue of the Yamabe operator,
$$L\frac{4(d1)}{d2}\mathrm{\Delta }_Y+\kappa $$
with $`\mathrm{\Delta }_Y`$ being the positive Laplacian acting on functions. The implication of the appearance of the Laplacian once again in this fermionic context is that there will be an input from the bosonic spectrum (as it transpires from (12) and (3) above) in setting the bound on the massive fermionic excitations. Therefore, the bounds on spin 1/2 and spin 0 masses are not totally independent
$$\nu _1^2\frac{d}{d2}\left[\frac{\pi ^2}{4\sigma ^2}\text{max}\{(d1)K,0\}\right]+\kappa $$
So, for positive curvature ($`K>0`$)
$$\nu _1^2\left(\frac{d}{d2}\right)\frac{\pi ^2}{4\sigma ^2}+\kappa $$
and for negative curvature ($`K<0`$)
$$\nu _1^2\frac{d}{d2}\left[\frac{\pi ^2}{4\sigma ^2}+(d1)K\right]+\kappa $$
In the the tensor product case, the above bounds read the same for the 4-dimensional masses, $`\mu _q`$’s. Thus by choosing $`\sigma ^1𝒪(\text{TeV})`$, we find that it is natural to achieve $`\mu _1^2\text{TeV}^2`$. In the case $`\kappa 0`$, all $`\mu _q^2\text{TeV}^2`$ without any specific value of the curvature. However, when $`\kappa <0`$ the curvature should satisfy an inequality similar to (4)
$$\kappa |(d1)K|\left(\frac{\pi ^2d}{d2}1\right)\text{TeV}^2$$
It is again remarkable that $`\mu _1^2\text{TeV}^2`$ can be naturally achieved having all our mass parameter of the same order of the compactification mass scale. As can be seen from the above inequalities, both $`\sigma `$ and the curvature explicitly enter the expressions of the bounds, and hence set the compactification mass scale.
### 4.3 Massive fermions: Upper bounds
Again, as in for $`\mathrm{\Delta }`$ eigenmodes, an upper bound on $`\nu _q^2`$ exists,
$$\nu _q^2Cq^{2/d}$$
where $`C`$ is a constant that depends only on the geometry of $`Y`$ (even in the presence of a gauge field) . Again here we find restrictions, though not as explicit as in (6,7,8), which limit our freedom in pushing up the KK masses arbitrarily high.
All the above observations, concerning both the upper and lower bounds, have been done in the tensor product case though the comments on zero modes apply equally to both types. However, if the product is warp, then the bounds and fermion masses will be dressed by the factor $`e^{\varphi /2}`$, as seen from (10), and similar arguments to the ones carried in section (3.2) apply.
## 5 Conclusions
We considered, on general grounds, a model of Einstein gravity coupled to a Dirac spinor and a Yang-Mills gauge theory on $`W=M_4Y`$, where $`Y`$ is a compact internal manifold with a scalar curvature bounded from below, and $`M_4`$ is our four dimensional world. Both tensor product and warp product are discussed. Bounds and estimates on the masses of the effective four dimensional theory at the classical level have been pointed out. Topological restrictions in choosing the internal manifold have been identified in order to avoid having certain bosonic massless modes in the four dimensional spectrum. In addition an upper bound on the curvature of $`Y`$ has been proposed, in the case of non-positive curvature. Geometrical upper and lower bounds have been presented for both bosons and fermions masses. In the tensor product case, the characteristic compactification mass scale for bosons is the diameter of the internal manifold, $`\sigma ^1`$, along with $`|K|`$ when the curvature $`\kappa <0`$. For fermions, the compactification mass scale is always set by $`\sigma `$ and the curvature, and this is due to an input from the bosonic spectrum in setting the bound on the massive fermionic excitations. Therefore, there is an interplay between spin 1/2 and spin 0 sectors.
For both fermions and bosons, it turns out that having the masses of the lowest excitations $`\text{TeV}`$ is naturally achieved by taking all the dimensionful parameters, arising from compactification, to be $`𝒪(\text{TeV})`$ (no fine-tuning required). In the warp product case, no direct bounds can be applied for massive states without the knowledge of both the specific shape of the warp factor and the field dependence on the internal space, though we are able to implement general estimates. “Zero-mode” arguments can be applied to both kinds of products. From the analysis conducted in this work, we conclude that non-positively curved internal manifolds with $`b_0=1`$ and $`b_1=0`$ are strongly favoured for phenomenological purposes.
Finally, a comment about non-compact internal manifolds: it has been argued that the spectrum of the Laplacian on non-compact spaces of finite volume has a discrete sector. Moreover, it has been shown recently that suitable choice of spin structure also leads to a discrete spectrum of the Dirac operator for non-compact hyperbolic manifolds of finite volume. One can therefore contemplate analyzing similar bounds for such non-compact spaces, along the same lines as we have executed in this work, and discuss their phenomenological implications .
Acknowledgement. One of us (R.T.) would like to thank Seif Randjbar-Daemi and George Thompson for useful discussions.
### Appendix
The expressions of the Laplacian acting on various tensors have been worked out in , and for convenience we list here the relevant expressions.
$`\mathrm{\Delta }\alpha `$ $`=`$ $`_i^i\alpha ={\displaystyle \frac{1}{\sqrt{g}}}(\sqrt{g}g^{ik}_k)\alpha `$
$`(\mathrm{\Delta }\alpha )_r`$ $`=`$ $`_i^i\alpha _r_r^h\alpha _h`$ (13)
$`(\mathrm{\Delta }\alpha )_{kl}`$ $`=`$ $`_i^i\alpha _{kl}+_k^h\alpha _{hl}+_l^h\alpha _{kh}2_{ki,lh}\alpha ^{ih}`$
where $`i,j,\mathrm{}etc=1,\mathrm{},`$dimension of the manifold on which the tensors and Laplacians are defined. $`_k^h`$ and $`_{ki,lh}`$ are Ricci and Riemann tensors respectively.
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# 1 Introduction
## 1 Introduction
Operating a future electron–positron linear collider (LC) in the LEP2 energy region but with much higher luminosity would allow a variety of precision tests of the electroweak sector to be performed, see for example Ref. . For instance, a 100 fb<sup>-1</sup> scan of the $`W^+W^{}`$ threshold with longitudinally polarized electrons and positrons would offer an opportunity to measure the $`W`$-boson mass with an error of $`6\mathrm{MeV}`$, , with negligible uncertainty from QCD interconnection effects, see for example Ref. . Such a threshold scan could also potentially provide a precise measurement of the $`W`$ decay width, $`\mathrm{\Gamma }_W`$.
At present, the most precise determination of $`\mathrm{\Gamma }_W`$ comes from the indirect measurement at the Tevatron using the ratio of dilepton $`Z`$ and $`W`$ events :
$$\mathrm{\Gamma }_W(\mathrm{CDF}+\mathrm{D0},\text{indirect})=\mathrm{\hspace{0.33em}2.171}\pm \mathrm{\hspace{0.33em}0.027}(\text{stat.})\pm \mathrm{\hspace{0.33em}0.056}(\text{sys.})\mathrm{GeV}$$
(1)
i.e. with an overall precision of approximately $`60\mathrm{MeV}`$. However it is very important to perform a direct determination of this key parameter of the Standard Model as well. The accuracy of the direct measurement of the $`W`$ width is still not very high. Recently the CDF collaboration at the Tevatron have reported the value :
$$\mathrm{\Gamma }_W(\mathrm{CDF},\text{direct})=\mathrm{\hspace{0.33em}2.055}\pm \mathrm{\hspace{0.33em}0.100}(\text{stat.})\pm \mathrm{\hspace{0.33em}0.075}(\text{sys.})\mathrm{GeV}$$
(2)
from measurements of the transverse mass spectrum in leptonic $`W`$ decays. Finally, the LEP experiments have made a preliminary measurement of the $`W`$ width from the line shape in $`W`$ pair production :
$$\mathrm{\Gamma }_W(\mathrm{LEP2},\text{direct})=\mathrm{\hspace{0.33em}2.19}\pm \mathrm{\hspace{0.33em}0.15}(\text{stat.}+\text{sys.})\mathrm{GeV}$$
(3)
Another related method, not yet exploited, would be to perform a precision scan of the $`WW`$ cross section in the threshold region, for example at a future linear collider.
In this paper we focus on another method which is based on previous studies in Ref. . This exploits the high sensitivity of soft photon radiation in pair production of $`W`$ bosons in $`e^+e^{}`$ collisions to the $`W`$ width. Since the event rate is $`𝒪(\alpha )`$ relative to the total $`WW`$ cross section, this method is potentially limited by statistics. However it has the advantage of avoiding problems such as the effect of beamsstrahlung or beam energy spread on line shapes. Whether it will ultimately be statistically competitive with a threshold scan in the precise determination of $`\mathrm{\Gamma }_W`$ will require a dedicated analysis, which is beyond the scope of our studies here. However, our results do suggest that such further investigation would certainly be worthwhile.
The paper is organized as follows. In the next section we discuss the overall structure of the $`e^+e^{}W^+W^{}\gamma `$ differential cross section, emphasizing the various energy regimes for the photon radiation. We present analytic results for the cross section in the soft-photon limit, which allows us to identify the factorizable and non-factorizable contributions. In Section 3 we discuss the various ways of enhancing the non-factorizable contributions, which contain the bulk of the $`\mathrm{\Gamma }_W`$ dependence, by imposing angular cuts on the final-state particles. We illustrate our results by numerical calculations, considering the various possible leptonic and hadronic decay channels. We also mention briefly the analogous results for the $`\gamma \gamma W^+W^{}\gamma `$ process. Section 4 contains our conclusions.
## 2 Bremsstrahlung radiation pattern in $`𝑾^\mathbf{+}𝑾^{\mathbf{}}`$ production
The general formalism for calculating the soft radiation pattern in processes involving the production and decay of unstable particles can be found in Refs. . It is well known that heavy unstable particles such as the $`W`$ boson can radiate before and after their decay. The relative intensity of the two contributions, and consequently the overall structure of the radiation pattern, depends sensitively on the relative size of the emission time-scale and the particle lifetime, see e.g. . In particular, in the second reference in one can find a semi-classical explanation of how the radiation pattern allows the relative distance between the $`W`$-boson decay vertices to be probed.<sup>1</sup><sup>1</sup>1This phenomenon resembles an old idea to use soft-photon radiation for measuring the time delay in nuclear reactions.
We begin by recalling the main properties of the differential distribution for the radiation of a soft photon with momentum $`k^\mu `$ in the process
$$e^+(q_1)e^{}(q_2)W^+(p_1)W^{}(p_2)\mathrm{\hspace{0.33em}4}\text{fermions}(p_3,p_3^{},p_4,p_4^{})\left[+\gamma (k)\right].$$
(4)
It is well known that when unstable particles are produced one is forced to perform a Dyson resummation, which leads to the regularization of the singularities in the propagators $`1/(M_{1,2}^2M_W^2)1/D_{1,2}`$, where $`D_{1,2}=M_{1,2}^2M_W^2+iM_W\mathrm{\Gamma }_W`$, and $`M_{1,2}`$ are the invariant masses of the $`W`$ bosons. However such resummation leads in general to the breaking of gauge invariance through higher-order contributions picked up by Dyson resummation. Thus this standard perturbative approach does not produce an acceptable gauge independent answer. The problem can be avoided by working in the so-called ‘pole-scheme’ . The physical picture behind the pole-scheme is as follows. Any process involving unstable particles can be viewed as a consequence of several subprocesses: production, which is a hard process with a short time-scale $`𝒪(1/M_W)`$; propagation over a typically larger time $`𝒪(1/\mathrm{\Gamma }_W)`$; and decay, which is again a hard process with a time scale $`𝒪(1/M_W)`$. Technically, in the perturbative expansion gauge invariance is guaranteed only order by order. Dyson resummation mixes different orders of the perturbative expansion, and thus breaks gauge invariance. In order to restore it one has to re-expand the amplitudes again in some physical parameter other than the coupling constant, in a way that does not produce singularities. An appropriate small parameter is $`\mathrm{\Gamma }_W/M_W`$. It is constructed as a ratio of two physical scales: the scale of production and decay, $`M_W`$, and the scale of propagation, $`\mathrm{\Gamma }_W`$. It should be noted that this is a somewhat simplified picture, since sometimes there are additional small parameters present in the problem (like the relative velocity, $`\beta `$, close to threshold, or $`M_W^2/s`$ at ultra-relativistic energies, etc.). Then the above mentioned estimates may change, but the arguments remain similar.
From the above considerations one can estimate the accuracy of calculations performed in the pole-scheme. When examining the process (4) one distinguishes three energy domains classified by the distance in energy from threshold, $`\mathrm{\Delta }E=\sqrt{s}2M_W`$, compared to the relevant scales of the process, $`\mathrm{\Gamma }_W`$ and $`M_W`$:
* Relativistic region, $`\mathrm{\Delta }EM_W`$, where the accuracy is $`𝒪(\mathrm{\Gamma }_W/M_W)`$.
* ‘Far-from-threshold’ region, $`\mathrm{\Gamma }_W\mathrm{\Delta }EM_W`$, where the accuracy is $`𝒪(\mathrm{\Gamma }_W/\mathrm{\Delta }E)`$.
* Threshold region, $`\mathrm{\Delta }E\mathrm{\Gamma }_W`$, where the accuracy is $`𝒪(1)`$, and the pole-scheme expansion breaks down.
The pole-scheme approach to processes involving unstable particles has been used to calculate the full $`𝒪(\alpha )`$ correction to the pair production of $`W`$ bosons in $`e^+e^{}`$ collisions . In this paper we use the results of Refs. as a basis for the calculations. Because of the way the pole-scheme is constructed, one can classify all the radiative corrections into two types: factorizable, which act inside separate hard subprocesses (production and decay); and non-factorizable, which interconnect various hard subprocesses. Here we will concentrate on the real photon radiation from $`W`$-pair production in the LEP2/LC energy region $`170500\mathrm{GeV}`$. Again, there are three regimes for photon radiation.
* Hard photon radiation, $`\omega M_W`$, when the photon wavelength is of the same order as the hard process time-scale. The photon can be assigned to one of the hard subprocesses. Alternatively one can say that the photons radiated from different stages of the process do not interfere with each other. The radiation is exclusively factorizable.
* Soft radiation, $`\omega \mathrm{\Gamma }_W`$, when the photon wavelength is much larger than the propagation distance. In this case the photons cannot distinguish the details of the process and are radiated coherently from all stages of the process. Both factorizable and non-factorizable contributions are important.
* Semi-soft radiation, $`\omega \mathrm{\Gamma }_W`$, when the photon wavelength is of the same order as the distance between the $`W`$ decay vertices. In this case both factorizable and non-factorizable contributions are important. However photons are not radiated coherently from all stages of the process.
From this classification one can see that at $`\omega \mathrm{\Gamma }_W`$ there is a transition from a regime in which various subprocesses do not interfere with each other to a regime in which the photon does not distinguish details of the process. This is this transition that is of interest to us in this paper, since it is where the photon spectrum has maximum sensitivity to $`\mathrm{\Gamma }_W`$.
Note that when $`\omega \mathrm{\Gamma }_W`$ the photon is soft with respect to the hard scale of the process, $`\omega M_W`$, but not with respect to the soft scale of the process, $`\mathrm{\Gamma }_W`$. As a consequence the cross section has certain factorization properties. The hard ($`M_W`$-scale) part of the amplitude factorizes just as in the conventional soft-photon approximation, but the soft ($`\mathrm{\Gamma }_W`$-scale) part does not always factorize in the conventional way. This is why we call photons with energy $`\omega \mathrm{\Gamma }_W`$ ‘semi-soft’, rather than simply ‘soft’. Making use of the factorization properties, the complete radiation distribution in this semi-soft regime can be written as an interference of semi-soft currents with the hard parts of the amplitudes in the following way:
$$d\sigma =d\sigma _{\text{Born}}\frac{d\stackrel{}{k}}{(2\pi )^32k_0}\left[2\text{Re }\left(_0_+^{}+_0_{}^{}+_+_{}^{}\right)+_0_0^{}+_+_+^{}+_{}_{}^{}\right].$$
(5)
Here $`d\sigma _{\text{Born}}`$ is the Born cross section in the pole approximation. The currents $`_0`$ and $`_\pm `$ correspond to the radiation from the production and decay stages respectively. The first three terms are non-factorizable contributions, consisting of final–final, $`_+_{}^{}`$, and initial–final, $`_0_+^{}+_0_{}^{}`$, state interferences. The last three terms are the factorizable contributions corresponding to the production and decay parts. The gauge-invariant semi-soft currents $`_0`$ and $`_\pm `$ are given by
$`_0^\mu `$ $`=`$ $`+e\left[{\displaystyle \frac{p_1^\mu }{kp_1}}{\displaystyle \frac{p_2^\mu }{kp_2}}{\displaystyle \frac{q_1^\mu }{kq_1}}+{\displaystyle \frac{q_2^\mu }{kq_2}}\right],`$
$`_+^\mu `$ $`=`$ $`e\left[{\displaystyle \frac{p_1^\mu }{kp_1}}+Q_{f_3}{\displaystyle \frac{p_3^\mu }{kp_3}}Q_{f_3^{}}{\displaystyle \frac{p_{3}^{}{}_{}{}^{\mu }}{kp_3^{}}}\right]{\displaystyle \frac{D_1}{D_1+2kp_1}},`$
$`_{}^\mu `$ $`=`$ $`+e\left[{\displaystyle \frac{p_2^\mu }{kp_2}}+Q_{f_4}{\displaystyle \frac{p_4^\mu }{kp_4}}Q_{f_4^{}}{\displaystyle \frac{p_{4}^{}{}_{}{}^{\mu }}{kp_4^{}}}\right]{\displaystyle \frac{D_2}{D_2+2kp_2}}.`$ (6)
The factors $`Q_f,Q_f^{}`$ are the electric charges of the final-state fermions, with $`Q_fQ_f^{}=1`$. Recall that the integration over the invariant masses of the unstable particles eliminates the pre-factors $`D_{1,2}/(D_{1,2}+2kp_{1,2})`$ in the factorizable terms. In this case the semi-soft currents become the usual soft-photon ones, and factorization takes place with respect to both scales of the process, hence the name ‘factorizable’. In the non-factorizable contributions, however, non-trivial pre-factors survive, and complete factorization with respect to both scales does not take place. An important consequence of this non-factorization is that for hard photons the non-factorizable contribution is suppressed because of the photon energy dependence in the pre-factors, see also Refs. . Thus non-factorizable contributions are important only for soft and semi-soft photons.
The qualitative picture described above is illustrated quantitatively in Fig. 1, which shows the photon energy spectrum, $`1/\sigma _{\text{Born}}\omega d\sigma /d\omega `$, is shown as a function of the photon energy, $`\omega `$, in the semi-soft region.<sup>2</sup><sup>2</sup>2Here and below we use the results and parameter values of Ref. for numerical calculations. In particular we use the Standard Model $`W`$-boson width $`\mathrm{\Gamma }_W=2.082\mathrm{GeV}`$. In this example the $`e^+e^{}`$ CMS energy is $`\sqrt{s}=184\mathrm{GeV}`$ and a purely leptonic $`(\mu ^+\nu _\mu )(\tau ^{}\overline{\nu }_\tau )`$ final state is chosen. A photon ‘isolation’ cut is also applied. This requires that in the CMS frame the direction of the radiated photon is separated by at least $`50^{}`$ from the directions of all the experimentally observed charged particles (i.e. the initial-state $`e^\pm `$ and the final-state $`\mu `$ and $`\tau `$ leptons). By imposing these ‘no-flight’ zones around the charged particles we avoid the quasi-collinear-singularities inherent in the currents in (2). <sup>3</sup><sup>3</sup>3In practice, the collinear singularities are regulated by non-zero fermion masses, see below.
We can see from Fig. 1 that the (negative) non-factorizable contribution to the cross section is indeed suppressed for hard photons, with the damping occurring in the semi-soft regime of the photon energy, $`\omega \mathrm{\Gamma }_W2\mathrm{GeV}`$. In the same photon energy region the dependence of the factorizable contribution on $`\omega `$ is practically flat. This leads to a peaking behaviour of the complete spectrum in the semi-soft region, $`\omega \mathrm{\Gamma }_W`$. Thus by comparing the measured photon bremsstrahlung distribution in this region with the theoretical prediction regarded as a function of $`\mathrm{\Gamma }_W`$, one can in principle determine the $`W`$-boson width, as advocated in Ref. .
In order to gain some quantitative insight on how the method may work in practice, two issues are important:
* How pronounced is the shape of the relevant part of the photon spectrum? In other words, how large is the interesting (strongly $`\mathrm{\Gamma }_W`$ dependent, see Fig. 1) non-factorizable contribution with respect to the factorizable contribution? The relevant parameter here is
$$\alpha (\sqrt{s},\text{cuts})=.\frac{(d\sigma _{\text{nf}}/d\omega )}{(d\sigma _{\text{fact}}/d\omega )}|_{\omega 0},$$
(7)
which depends on the system of cuts chosen and the collider energy, $`\sqrt{s}`$.
* How large is the statistics for a particular choice of cuts? The relevant parameter is the corresponding Born cross section restricted by a particular system of cuts, $`d\sigma _{\text{Born}}(\sqrt{s},\text{cut})`$.
To illustrate how the ratio of non-factorizable to factorizable radiation, $`\alpha (\text{no cuts},\sqrt{s})`$, is influenced by the photon isolation cuts, we show in Fig. 2 the ratio as a function of the CMS energy, $`\sqrt{s}`$, with and without cuts. In order to use a logarithmic scale we plot the absolute values of the ratios, and indicate their sign in the legend of the plot. Here and in what follows we label all quantities by two letters, which specify the decay channel of each of the $`W`$’s, $`L`$ for leptonic and $`H`$ for hadronic, and a subscript, which specifies the system of cuts applied to the kinematics. In this case we consider a purely leptonic final state (in this example $`\mu ^+\tau ^{}`$), thus the label is $`LL`$. The factorizable radiation for the ‘no-cut’ case depends on the masses of the charged fermions through collinear logarithms.
In Fig. 2, in addition to the combined factorizable/non-factorizable effect, $`LL^{\text{nf/fact}}`$, we show separately the initial-final, $`LL^{\text{if/fact}}`$, and final-final, $`LL^{\text{ff/fact}}`$, state ratios. Note that they have opposite signs.<sup>4</sup><sup>4</sup>4This agrees with the observations of Ref. where gluon radiation in $`e^+e^{}t\overline{t}bW^+\overline{b}W^{}`$ was discussed. One can see from the figure that the final-final part of the non-factorizable correction scales as $`\alpha E^4`$ with the CMS energy. In fact the power-counting arguments of Ref. are applicable to the parameter $`\alpha `$, with a small modification due to the fixed rather than integrated photon energy which does not however change the result. For the initial-final part of the interference the energy scaling is different: $`\alpha E^2`$. One can also see that initial-final state interference dominates the non-factorizable effects at high energies,<sup>5</sup><sup>5</sup>5This is again in agreement with observations of Ref. . and thus the complete non-factorizable radiation contribution also scales as $`E^2`$. If one does not apply any cuts, the ratio of non-factorizable to factorizable contributions is small, below $`3\%`$. This is mainly due to the enhancement of the collinear logarithms in the factorizable part of the radiation. However if one keeps the photon well separated from the charged particles, and thus well away from the collinear regions, as in the $`LL_\gamma `$ ratio in Fig. 2 then the collinear logarithms are suppressed and the ratio increases considerably. For example, in Fig. 2 $`\alpha 1030\%`$ for the $`LL_\gamma `$ $`>50^{}`$ cut, depending on the collider energy.
## 3 Enhancement mechanisms for different external states
In Ref. two processes were considered in detail: gluon radiation in $`e^+e^{}t\overline{t}`$ and photon radiation in $`\gamma \gamma W^+W^{}`$. Both cases were considered for collision energies close to threshold. The process we are interested in, (4), differs from the studies of in two respects.
First, we consider higher collision energies, which are experimentally more relevant. Moreover the pole expansion, which we use in our calculations, does not apply at threshold. On the other hand, as we have already seen, at higher-energies non-factorizable effects become relatively small and therefore the sensitivity to $`\mathrm{\Gamma }_W`$ is less . Without any cuts, the non-factorizable corrections, (5), to distributions inclusive with respect to angles scale as $`E^2`$ (initial-final state interference) or $`E^4`$ (final-final state interference) relative to the Born cross section. As a result of this scaling behaviour, at $`\sqrt{s}=184\mathrm{GeV}`$ the ratio of non-factorizable to factorizable contributions to the photon spectrum is $`𝒪(1\%)`$ or smaller. The main objective of our studies is to enhance the non-factorizable effects by applying angular cuts.
A second difference with the study of Ref. is that there the effects of initial state radiation were not fully considered. In terms of Eq. (5), the analysis of Ref. was concerned with only one of the non-factorizable effects, from final-final state interference. As we have already noted, for $`e^+e^{}W^+W^{}+\gamma `$ the situation is more complicated, with three interference contributions: two initial-final and one final-final state.
In Ref. two ways to enhance the non-factorizable effects were proposed. First, it was suggested that certain angular asymmetry properties of initial-final state interference, absent in final-final state interference, could be used to construct observables to which initial-final state interference does not contribute. Moreover, in an observable was constructed which has no contributions from factorizable radiation, using the fact that the factorizable correction does not depend on the angles of the produced particles, at least at threshold (see Ref. for more details).
Unfortunately, the construction of such observables is impractical. Because of $`t`$-channel neutrino exchange, there are always spin-charge correlations present, even in the threshold Born cross section. Initial-final spin correlations induced by the $`W`$ propagators lead to an asymmetry in the factorizable part of the radiation, as well as in the final-final state interference contributions. Therefore the method proposed in Ref. does not appear to be workable. Note that this effect originates in the $`(va)`$ structure of the charged weak current. The Born DPA cross section does not violate $`P`$-parity because only pole residues are calculated. However the type of helicity-charge correlation described above does survive. Technically, anti-symmetric tensors, $`ϵ_{\mu \nu \rho \sigma }`$, induced by the axial current do not contribute linearly to the matrix element (no $`P`$-violation), but only quadratically, via the interference of axial contributions from various stages of the process (the helicity-charge asymmetry).<sup>6</sup><sup>6</sup>6Note that in $`W`$ pair production in $`\gamma \gamma `$ collisions the asymmetry in the threshold Born cross section is absent, because there is no $`(va)`$ structure at the production stage. Therefore, all the results of Ref. remain valid for the $`\gamma \gamma W^+W^{}`$ case. The asymmetry is also absent in $`e^+e^{}ZZ`$ production, because of Bose symmetry. On the other hand the asymmetry will be present, even without the $`(va)`$ coupling of initial-state fermions (as in QED for example), if the initial-state fermions are polarized.
Another idea discussed in was based on the ‘angular ordering’ effect. During the last two decades such angular ordering effects have been intensively discussed in the context of the QCD cascades, see for example . The phenomenon itself has been well known for QED in cosmic ray physics from the middle of the 1950s as the so-called Chudakov effect . To recall the physics of angular ordering, we consider the radiation pattern of soft photons produced by a relativistic $`e^+e^{}`$ pair. If we split the radiation into pieces associated with the $`e^{}`$ and $`e^+`$, and then integrate over the azimuthal angle about, say, the $`e^{}`$ direction, the $`e^{}`$ contribution vanishes for polar angles greater than the $`e^+e^{}`$ opening angle. In particular this implies that the radiation vanishes for collinear $`e^+`$ and $`e^{}`$. In other words, for such a configuration the emitted photon probes only the total electric charge of the $`e^+e^{}`$ pair, which is zero. The suppression of radiation is caused by the (destructive) interference between the emission off the $`e^{}`$ and $`e^+`$, see the second reference in for more details. Because in the present context it is the $`W`$ lifetime that controls this interference pattern, we expect to observe angular ordering behaviour (or not) according to the size of the ratio $`\omega /\mathrm{\Gamma }_W`$. It is therefore clear that the largest effect of non-factorizable radiation relative to factorizable radiation will correspond to the case of collinear oppositely charged particles. In that case factorizable radiation is as important as non-factorizable radiation.<sup>7</sup><sup>7</sup>7 Note that due to the celebrated Low-Kroll-Barnett soft bremsstrahlung theorem the non-classical short-distance-induced corrections to the angular ordering behaviour arise only on the level of quadratic in $`\omega /M`$ terms, see Ref. .
In the case of $`W`$ pair production with $`\mu \nu _\mu \tau \nu _\tau `$ decay in the threshold region there are four radiating charged particles: two initial state fermions, $`e^\pm `$, and two final state fermions, $`\mu ^+,\tau ^{}`$. Corresponding to this there are three non-factorizable interferences: two initial-final, and one final-final state interference. Clearly it is impossible to generate a large effect from all three interferences simultaneously. Indeed, if the $`e^+`$ and $`\mu ^+`$ are collinear and the $`e^{}`$ and $`\tau ^{}`$ are collinear, then the $`\mu ^+`$ and $`\tau ^{}`$ are anti-collinear. Far above threshold, the directions of the $`W`$-boson momenta start to play a role as well. Thus in general one has many cases when some oppositely charged particles are collinear and others are not, leading to a non-trivial interplay between the various interference terms in (5). Rather than choose particular fixed configurations, for which the statistics will be small, it is more efficient to look for angular cuts (no-flight zones) on the various particles such that the interesting (i.e. most $`\mathrm{\Gamma }_W`$ dependent) events are favoured but not overly restricted. We shall not in the present study make any serious attempt to optimize these cuts; rather we will present some illustrative examples pending more detailed Monte Carlo analyses.
In summary, angular cuts (no-flight zones) will be applied to the final state particles (leptons and quarks) and the photon in order to maximize the angular ordering effect, and thus the sensitivity of the photon spectrum to the $`W`$ width. As explained above, the basic idea is to keep oppositely charged particles quasi-collinear, and the photon as far from them as possible. There is an additional requirement motivated by detector considerations. The final state particles should not be too close to the beam direction otherwise the event cannot be unambiguously identified as $`W^+W^{}\gamma `$. We therefore require all final state particles to be produced at polar angles greater than $`5^{}`$ from the beam direction.<sup>8</sup><sup>8</sup>8We are grateful to G. Wilson for clarification of various experimental issues related to $`W`$ studies at a future linear collider.
### 3.1 Leptonic-leptonic final state
The simplest case to analyse is when both $`W`$’s decay leptonically:
$$e^+(q_1)+e^{}(q_2)\mu ^+(p_3)+\tau ^{}(p_4)+\mathrm{\hspace{0.33em}2}\nu +\gamma (k).$$
(8)
The first topology we will consider is when the two final-state charged particles are close to the beam direction. In this case the initial-final state interference gives a large effect. In addition, the photon should be far from the beam direction:
$$\mathrm{}(q_{1,2}k)>50{}_{}{}^{}.$$
(9)
The charged final-state leptons with momenta $`p_3`$ and $`p_4`$ can each be either collinear to the initial-state positron or to the electron. ‘Collinear’ is here defined as being produced with polar angle between $`5^{}`$ and $`10^{}`$ with respect to the beam direction:
$$\mathrm{}(q_1p_{3,4})(5{}_{}{}^{},10{}_{}{}^{}),\text{or}\mathrm{}(q_2p_{3,4})(5{}_{}{}^{},10{}_{}{}^{}).$$
(10)
We therefore have four possible cases, which we label
$$LL_{++},LL_+,LL_+,LL_{\times \times },$$
(11)
corresponding to ($`p_3q_1`$ and $`p_4q_1`$), ($`p_3q_1`$ and $`p_4q_2`$), and ($`p_3q_2`$ and $`p_4q_1`$) correspondingly. $`LL`$ refers to the fact that both $`W`$ bosons decay leptonically. In the last case $`LL_{\times \times }`$ we demand only that the final-state leptons are collinear with the electron and positron beams, without tracing the electric charge flow.
The second class of cuts we will consider is when two final-state particles are quasi-collinear. In this case it is the final-final state interference that produces a large effect. We first demand that all final-state particles are observable
$$\mathrm{}(q_{1,2}p_{3,4})>5{}_{}{}^{},$$
(12)
and then that the final-state charged particles are collinear
$$\mathrm{}(p_3p_4)<10{}_{}{}^{},$$
(13)
and the photon is far from all charged particle directions
$$\mathrm{}(p_{3,4}k)>50{}_{}{}^{},\mathrm{}(q_{1,2}k)>50{}_{}{}^{}.$$
(14)
Here there is only one possible case, which we label
$$LL_\text{f}.$$
(15)
The optimization parameters $`\alpha (\sqrt{s},\text{cuts})`$ and $`\sigma _{\text{Born}}(\sqrt{s},\text{cuts})`$ are shown in Fig. 3 as functions of the CMS energy, $`\sqrt{s}`$, for all five possible leptonic cuts. We see that at low energies the most pronounced shape of the photon spectrum is achieved for the $`LL_+`$ case, i.e. $`\mu ^+`$ collinear with incoming $`e^{}`$ and $`\tau ^{}`$ collinear with incoming $`e^+`$. Then the ratio of the non-factorizable to and factorizable contributions is positive and can even exceed 1, in the lower energy domain. However, the Born cross-section for this set of cuts is very small, which makes this case statistically disadvantageous. At high energies, the outgoing fermion (antifermion) prefers to follow the direction of the incoming fermion (antifermion), and hence both the $`LL_+`$ and $`LL_{\times \times }`$ configurations have large Born cross sections, $`\sigma _{\text{Born}}`$. In terms of the shape parameter, the $`LL_+`$ cut is as good as $`LL_{\times \times }`$. At lower energies, however, $`LL_+`$ becomes more advantageous in terms of shape, but less advantageous in terms of statistics. In fact, referring back to Fig. 2, we see that the original $`LL_\gamma `$-cut is as good in terms of shape as $`LL_{\times \times }`$ at low energy, but much better statistically since it corresponds to a much larger angular acceptance for the final-state charged particles. At higher energies it is still as good in terms of shape, but becomes statistically very poor.
The conclusion is that depending on the energy and statistics available, one can choose different systems of cuts as the preferred ones. There is clearly no unique ‘best cut’ for all energies and all statistics.
### 3.2 Leptonic-hadronic final state
We next consider the case when the $`W^+`$ decays leptonically and the $`W^{}`$ decays hadronically. There is one charged lepton and two jets present in the final state:
$$e^+(q_1)+e^{}(q_2)\mu ^+(p_3)+q(p_4)+\overline{q}^{}(p_4^{})+\nu +\gamma (k).$$
(16)
We again start with the case when the charged primary fermions are close to the beam directions. Just as in the lepton-lepton case, we demand that the photon is far from the beams, $`\mathrm{}(q_{1,2}k)>50^{}`$ and the final-state lepton with momenta $`p_3`$ is either collinear to the initial-state positron or electron, $`\mathrm{}(q_1p_3)(5{}_{}{}^{},10{}_{}{}^{})`$ or $`\mathrm{}(q_2p_3)(5{}_{}{}^{},10{}_{}{}^{})`$. The two quarks with momenta $`p_4`$ and $`p_4^{}`$ should also be collinear to the initial-state particles:
$$\mathrm{}(q_1p_4)(5{}_{}{}^{},20{}_{}{}^{}),\text{or}\mathrm{}(q_2p_4)(5{}_{}{}^{},20{}_{}{}^{}),$$
(17)
and
$$\mathrm{}(q_1p_4^{})(5{}_{}{}^{},20{}_{}{}^{}),\text{or}\mathrm{}(q_2p_4^{})(5{}_{}{}^{},20{}_{}{}^{}).$$
(18)
An important difference here is that one cannot measure the charge of the jet experimentally. Thus in general the following combinations are available:
$$LH_{+(20)},LH_{+(02)},LH_{(20)},LH_{(02)},$$
$$LH_{+(11)},LH_{(11)},LH_{+(\times \times )},LH_{(\times \times )},LH_{\times (\times \times )},$$
where $`LH`$ denotes the leptonic-hadronic final state. The first subscript indicates the direction of the lepton with respect to the positron momentum, and the two subscripts in parenthesis indicate the number of jets collinear with the positron and electron. For example, $`LH_{+(20)}`$ means that the final-state lepton is collinear with the positron, as are both jets. $`LH_{(\times \times )}`$ means that the final-state lepton is collinear with the electron, and the two jets are collinear to either the positron or electron. Thus, in general, the number of different cases is rather large compared to the purely leptonic final state. For the energies we are interested in, $`\sqrt{s}=170500\mathrm{GeV}`$, however, the situation simplifies somewhat, because the kinematics are such that the two jets coming from the decay of the $`W`$ boson cannot in fact satisfy the collinearity selection criterion. Thus
$$LH_{+(20)}=LH_{+(02)}=LH_{(20)}=LH_{(02)}=0.$$
The following cases survive
$$LH_{+(11)}=LH_{+(\times \times )}LH_{+\times },LH_{(11)}=LH_{(\times \times )}LH_\times ,LH_{\times (\times \times )}LH_{\times \times },$$
(19)
where we have introduced the modified notation $`LH_{+\times }`$, $`LH_\times `$ and $`LH_{\times \times }`$.
The second class of cuts again corresponds to the situation when two final-state particles are collinear. In this case the final-final state interference gives a large effect. We demand that all final-state particles are observable
$$\mathrm{}(q_{1,2}p_{3,4})>5{}_{}{}^{},\mathrm{}(q_{1,2}p_4^{})>5{}_{}{}^{},$$
(20)
at least two final-state particles are quasi-collinear
$$\mathrm{}(p_3p_4)<10{}_{}{}^{},\text{or}\mathrm{}(p_3p_4^{})<10{}_{}{}^{},$$
(21)
and the photon is far from all charged particles
$$\mathrm{}(p_3k)>50{}_{}{}^{},\mathrm{}(p_4k)>50{}_{}{}^{},\mathrm{}(p_4^{}k)>50{}_{}{}^{},\mathrm{}(q_{1,2}k)>50{}_{}{}^{}.$$
(22)
Thus there is again only one possible choice
$$LH_\text{f}.$$
(23)
We show in Fig. 4 the optimization parameters $`\alpha (\sqrt{s},\text{cuts})`$ and $`\sigma _{\text{Born}}(\sqrt{s},\text{cuts})`$ as functions of the CMS energy, $`\sqrt{s}`$, for all possible $`LH`$-cuts. Again, we see that the $`LL_\times `$-cut is the best in terms of the shape of the spectrum, but at the same time it is the worst in terms of statistics. $`LH_f`$ seems to be statistically the best overall throughout the energy region under consideration. From the point of view of the shape of the spectrum, $`LH_f`$ is not worse than any other system of cuts for higher energies. Note that for cuts restricting the jets to be quasi-collinear with the collider beams, the energy behaviour of the shape parameter $`\alpha `$ is more complicated than in the case of leptonic-leptonic final states. As mentioned above, if one of the quarks is quasi-collinear with the electron, the other one is automatically quasi-collinear with the positron. The effects coming from the two corresponding interferences have opposite signs. This can even lead to a change of sign of the combined effect at different collider energies.
### 3.3 Hadronic-hadronic final state
Finally we consider the case when both $`W`$ bosons decay hadronically, with four jets present in the final state.
We again start from the case when the charged particles (i.e. jets) are collinear with the beam direction, with the photon well separated from the beam, $`\mathrm{}(q_{1,2}k)>50^{}`$. The quarks with momenta $`p_{3,4}`$ and $`p_{3,4}^{}`$ are required to be collinear with the initial leptons, $`\mathrm{}(q_{1,2}p_{3,4})(5{}_{}{}^{},20{}_{}{}^{})`$ and $`\mathrm{}(q_{1,2}p_{3,4}^{})(5{}_{}{}^{},20{}_{}{}^{})`$. Since one cannot measure the charge of the jet experimentally. only the following combinations are available:
$$HH_{22},HH_{13},HH_{\times \times },$$
where the subscript denotes the number of jets that are collinear with the initial-state positron or electron. For example, $`HH_{13}`$ means that there is one jet collinear with the positron, $`q_1`$, and three jets collinear with the electron, $`q_2`$. Again, the kinematics are such that at LEP2 energies not all of these cases are non-zero:
$$HH_{13}=0,$$
and in fact only one case survives:
$$HH_{22}=HH_{\times \times }.$$
(24)
The second class of cuts corresponds to when two final-state particles (jets) are quasi-collinear. In this case final-final state interference gives a large effect. We first demand that all final state particles are observable
$$\mathrm{}(q_{1,2}p_{3,4})>5{}_{}{}^{},\mathrm{}(q_{1,2}p_{3,4}^{})>5{}_{}{}^{},$$
(25)
at least two final state particles are collinear
$$\mathrm{}(p_3p_4)<10{}_{}{}^{},\text{or}\mathrm{}(p_3p_4^{})<10{}_{}{}^{},\text{or}\mathrm{}(p_3^{}p_4)<10{}_{}{}^{},\text{or}\mathrm{}(p_3^{}p_4^{})<10{}_{}{}^{},$$
(26)
and the photon is far from all of the charged particles
$$\mathrm{}(p_{3,4}k)>50{}_{}{}^{},\mathrm{}(p_{3,4}^{}k)>50{}_{}{}^{},\mathrm{}(q_{1,2}k)>50{}_{}{}^{}.$$
(27)
There is again only one possible case:
$$HH_\text{f}.$$
(28)
The optimization parameters $`\alpha (\sqrt{s},\text{cuts})`$ and $`\sigma _{\text{Born}}(\sqrt{s},\text{cuts})`$ are shown in Fig. 5 for the two possible cuts $`HH_{\times \times }`$ and $`HH_f`$. As in the previous cases, $`HH_f`$ cut is better statistically, but $`HH_{\times \times }`$ is better from the point of view of the shape of the photon spectrum.
### 3.4 Photon-photon colliders
In recent years there has been a growing interest in high-energy photon colliders, using Compton back-scattering of laser light off the lepton beams at linear colliders to produce high-intensity, high-energy beams of photons, see e.g. Ref. . Using $`\gamma \gamma `$ collisions to produce pairs of $`W`$ bosons offers certain advantages over the $`e^+e^{}`$ case. First, the cross section is an order of magnitude larger. Second, ISR effects are absent in this case and so kinematic reconstruction of the $`WW`$ final state is in principle more precise.
It is straightforward to extract the predictions for photon radiation in $`\gamma \gamma W^+W^{}`$ from our study of the more complicated $`e^+e^{}`$ case. In particular, our results for the final-final state interferences $`LL_\text{f}`$, $`LH_\text{f}`$ and $`HH_\text{f}`$ can be applied directly to the $`\gamma \gamma `$ case. Moreover, as we have already explained, in the case of $`W^+W^{}`$ production in photon-photon collisions one can study observables integrated over the photon angle, to which factorizable corrections do not contribute, see Ref. .<sup>9</sup><sup>9</sup>9Note that there is a typo in Eq. (20) of the first reference in . The normalization coefficients in front of the integrals in the first and second terms should be interchanged. The final result given by Eq. (22) is unchanged. This enables us to utilize more events and makes studies in $`\gamma \gamma `$ collisions potentially more statistically powerful than in the $`e^+e^{}`$-case.
## 4 Concluding remarks
A precision measurement of the total $`W`$ decay width presents a challenge for present and future experiments. Line-shape measurements are made difficult by the presence of neutrinos in the final state in the case of leptonic decay modes, and of hadronization corrections in the case of hadronic decays. The indirect measurement at hadron colliders, which uses the ratio of $`W`$ and $`Z`$ leptonic events, has an inherent uncertainty from parton distributions in the theoretical calculation of the total cross sections. It seems to be quite a challenging task to perform a precise direct measurement of the total $`W`$-width, independent of decay modes (and of the $`Z`$ measurements).
As discussed in Ref. , running a future linear collider in the ‘LEP2’ energy region may provide a unique opportunity for a high-precision measurement of the $`W`$ mass and width. The ’traditional’ way of measuring $`\mathrm{\Gamma }_W`$ is from a threshold scan of the total $`WW`$ cross section. Though statistically powerful, this method is not without problems. The uncertainties caused by beam-induced effects (beamsstrahlung, intrinsic energy spread, etc.) could be potentially large. Moreover, the threshold strategy requires operating a linear collider at energy scan points well below threshold where the $`W^+W^{}`$ cross section is very small.
In this paper we have argued that the soft-photon radiation spectrum could also be used to obtain information on $`\mathrm{\Gamma }_W`$. We emphasize that this is an independent method — in effect one is measuring the non-factorizable interference to the cross section, whose magnitude is controlled by the relative size of the photon energy and the $`W`$ width. The method is in principle very clean, requiring only a precise measurement of the soft (i.e. of order few GeV) photon spectrum in $`W^+W^{}\gamma `$ events. However, as we have seen, the effect in the inclusive distribution is very small and therefore is likely to be limited by statistics. On the other hand, we have shown that one can enhance the effect by employing angular cuts on the final-state particles. We have considered various different topologies and different $`W`$ decay channels. Both the sensitivity to the non-factorizable contributions and the overall number of events in the various channels are rather strongly dependent on the collision energy, and it should be possible to develop an optimal strategy given the parameters and running conditions of a future linear collider.
Our study necessarily falls short of any firm conclusion about the competitivity of our method, compared to the threshold scan for example, in determining $`\mathrm{\Gamma }_W`$. At the very least, our method offers a complementary measurement, with completely different systematics. The next step would be to perform a detailed Monte Carlo study including detector and, where appropriate, hadronization effects. Among the questions that such a study could answer are: what is the efficiency for detecting very soft photons? can such photons be measured in the presence of hadronic jets? are the isolation and collinearity cuts we have used realistic? For a given collider energy it should be straightforward to estimate the number of soft photon events for each of the different topologies and decay channels, and by comparing this with the theoretical predictions, to estimate the statistical error on $`\mathrm{\Gamma }_W`$. The results of our work suggest that a more detailed study is definitely worth pursuing.
Acknowledgements
We thank T. Sjöstrand, L. Stodolsky and G. Wilson for useful discussions. VAK thanks the Leverhulme Trust for a Fellowship. The work of APC is supported in part by the UK Particle Physics and Astronomy Research Council. This work was also supported by EU Fourth Framework Programme “Training and Mobility of Researchers”, contract FMRX-CT98-0194 (DG 12 - MIHT).
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# A Rational Surgery Formula for the LMO Invariant
## 1. Introduction
In their paper “Wheels, Wheeling, and the Kontsevich Integral of the Unknot”, B-N, S. Garoufalidis, L. Rozansky and D. P. Thurston \[BGRT\] made two conjectures; the “Wheels” conjecture about the value of the Kontsevich integral of the unknot, and the “Wheeling” conjecture about the relationship between the two natural products on the space of uni-trivalent diagrams. We quote here a paragraph from \[BGRT\], explaining a part of their motivation:
> If \[the Wheeling conjecture\] is true, one would be able to use it along with \[the Wheels conjecture\] and known properties of the Kontsevich integral (such as its behavior under the operations of change of framing, connected sum, and taking the parallel of a component as in \[LM2\]) to get explicit formulas for the Kontsevich integral of several \[…\] knots and links. …Likewise, using \[these conjectures\] and the hitherto known or conjectured values of the Kontsevich integral, one would be able to compute some values of the LMO 3-manifold invariant \[LMO\], using the “Århus integral” formula of \[Å-I–III\].
The “Wheels” and “Wheeling” conjectures are now theorems \[Ko, Moc, HV, BLT, Th\], and this seems a good time to proceed with the plan outlined above. Thus the purpose of our note is to use Wheels and Wheeling and the Århus integral to obtain some explicit formulas for the values of the Kontsevich integral and the LMO invariant on various simple knots, links and 3-manifolds, as well as a general formula for the behavior of the LMO invariant under rational surgery over links (previously such formulas existed only for integral surgery).
### 1.1. The main result: a rational surgery formula for the LMO invariant
The LMO invariant $`\widehat{Z}^{\mathrm{LMO}}(M)`$ of a rational homology 3-sphere $`M`$, which is given by surgery on some regular integrally framed link $`L`$ on $`S^3`$ (we write, $`M=S_L^3`$; “regular” means that the linking matrix of $`L`$ is non-singular), is a properly normalized version of the integral of a certain renormalized version of the Kontsevich integral $`Z(L)`$ of $`L`$. Thus following \[Å-I–III\], we write (see alternative form at Equation (33))
(1) $`\stackrel{ˇ}{Z}(L)`$ $`:=`$ $`\nu ^XZ(L)𝒜(_X),`$
(2) $`\text{Å}_0(L)`$ $`:=`$ $`{\displaystyle \sigma \stackrel{ˇ}{Z}(L)𝑑X};Z_\pm :=\text{Å}_0(^{\pm 1})`$
(3) $`\widehat{Z}^{\mathrm{LMO}}(M)`$ $`:=`$ $`Z_+^{\varsigma _+(L)}Z_{}^{\varsigma _{}(L)}\text{Å}_0(L).`$
In these equations:
* $`X=(x_i)`$ is the set of components of $`L`$, $`𝒜(_X)`$ is the space of chord diagrams<sup>1</sup><sup>1</sup>1 Throughout this paper all the spaces of diagrams that we will consider are graded and we will automatically complete them with respect to the grading, so as to allow infinite sums of terms of increasing degree. whose skeleton consists of $`|X|`$ directed circles colored by the elements of $`X`$ modulo the usual $`4T`$, STU, AS and IHX relations, but not divided by the framing independence relation. For general background about chord diagrams and related topological and Lie algebraic issues see \[B-N\].
* $`\nu =Z()𝒜()𝒜()`$ is the Kontsevich integral of the zero-framed unknot $``$, regarded as an element of the space $`𝒜()𝒜()`$ of chord diagrams whose skeleton is a single directed circle (modulo the same relations as above), or a single directed line (modulo the same).
* $`\nu ^X`$ is the $`|X|`$’th tensor power of $`\nu `$, regarded as an element of $`𝒜()^X`$, the $`X`$’th tensor power of $`𝒜()`$. It acts on the Kontsevich integral $`Z(L)`$ of $`L`$ using the usual “stick in anywhere” action $`𝒜()^X𝒜(_X)𝒜(_X)`$.
* $`\sigma :𝒜(_X)𝒜(_X)`$ is the formal PBW linear isomorphism between the space $`𝒜(_X)`$ and the space $`𝒜(_X)`$ (denoted $`^{\text{links}}(X)`$ in \[Å-II\]) of $`X`$-marked uni-trivalent diagrams modulo AS, IHX and $`X`$-flavored link relations (see \[Å-II, Section 5.2\]). The map $`\sigma `$ is most easily defined as the inverse of the symmetrization map $`\chi :𝒜(_X)𝒜(_X)`$. If $`X`$ is a singleton, we often suppress it and write $`𝒜()`$, $`𝒜()`$, $`\sigma :𝒜()𝒜()`$, etc. We note that $`\sigma `$ is a homonymous variant of a better known isomorphism $`\sigma :𝒜(_X)𝒜(_X)`$, where $`𝒜(_X)`$ is the same as $`𝒜(_X)`$ but with the directed circle skeleton components replaced by directed lines and $`𝒜(_X)`$ is the standard space of $`X`$-marked uni-trivalent diagrams modulo AS and IHX, not reduced by link relations.
* $``$ is the key ingredient of “formal integration”. It can be either the LMO-style “negative dimensional integral” $`^{(m)}`$, or, in the case when $`M`$ is a rational homology sphere, the “formal Gaussian integration” $`^{FG}`$ of \[Å-I–III\]. (The equality of these two integration theories is in \[Å-III\]; a sticky leftover from \[Å-I–III\], that $`^{FG}`$ is well defined modulo link relations, is our Proposition 2.2). We note that $``$ is valued in $`𝒜(\mathrm{})`$, the space of trivalent diagrams modulo AS and IHX (that is, unitrivalent diagrams with no univalent vertices), and that $`𝒜(\mathrm{})`$ is a commutative algebra under disjoint union.
* $`^{\pm 1}`$ is the $`\pm 1`$-framed unknot, and $`\varsigma _+(L)`$ and $`\varsigma _{}(L)`$ are respectively, the numbers of positive and negative eigenvalues of the covariance matrix of $`Z(L)`$, which is the linking matrix of $`L`$ (see \[Å-I, Definition 2.8 and Claim 1.10\]).
It is rather easy to show (see Section 4.1) that if a rational homology 3-sphere $`M`$ is given by surgery on some rationally framed link $`L`$, then its LMO invariant $`\widehat{Z}^{\mathrm{LMO}}(M)`$ can be computed using exactly the same formulas (1)–(3), only replacing the input $`Z(L)`$ (which is not defined for rationally framed links) by some extension thereof, which we will also denote by $`Z(L)`$. Our main result in this paper is a precise formula for this “rationally-framed Kontsevich integral”. We start with some definitions.
###### Definition 1.1.
(\[BGRT\]) Let $`\mathrm{\Omega }𝒜()`$ be given by
(4)
$$\mathrm{\Omega }=\mathrm{exp}_\stackrel{\Gamma }{}\underset{m=1}{\overset{\mathrm{}}{}}b_{2m}\omega _{2m},$$
where the modified Bernoulli numbers $`b_{2m}`$ are defined by the power series expansion
(5)
$$\underset{m=0}{\overset{\mathrm{}}{}}b_{2m}x^{2m}=\frac{1}{2}\mathrm{log}\frac{\mathrm{sinh}x/2}{x/2}$$
(so that $`b_2=1/48`$, $`b_4=1/5760`$, etc.) and $`\omega _{2m}`$ is the $`2m`$-wheel, the degree $`2m`$ uni-trivalent diagram made of a $`2m`$-gon with $`2m`$ legs (so that $`\omega _2=`$, $`\omega _4=`$, …, with all vertices oriented counterclockwise).
###### Definition 1.2.
(\[Å-I–III\]) For any element $`A𝒜(_X)`$, define a map
$$\widehat{A}=_A:𝒜(_X)𝒜(_X)$$
to act on diagrams $`B𝒜(_X)`$ by gluing all legs of $`A`$ to some subset of legs of $`B`$ with matching labels. Likewise, define a pairing
$$A,B_X:=(_AB)_{X0}:=\left(\text{sum of all ways of gluing the }x\text{-marked legs of }A\text{ to the }x\text{-marked legs of }B\text{, for all }xX\right).$$
We note that $`\widehat{A}=_A`$ and $`A,`$ are well defined even for arguments $`B𝒜(_X)`$, provided $`_A`$ annihilates all $`X`$-flavored link relations. This is equivalent to saying that $`A`$ is invariant with respect to $`x`$ for all $`xX`$, where “invariance” is defined in Figure 1.
###### Remark 1.3.
Strictly speaking, $`_AB`$ and $`A,B`$ are ill-defined if both $`A`$ and $`B`$ contain struts (see \[Å-II, Section 2.2\]), so that closed vertex-free loops can be formed by gluing them together. We will not encounter this problem in this paper.
###### Remark 1.4.
Formally $`_A`$ acts as if it were a differential operator on multi-variable polynomials in symbols indexed by $`X`$, according to an operator obtained from $`A`$ by replacing each label $`x`$ by $`_x`$ (see \[Å-II, Section 2\]). Thus while dealing with specific diagrams we will sometimes use notation as follows:
$$\begin{array}{c}\text{ }\text{}\text{ }\end{array}$$
Next we define the Dedekind symbol $`S(p/q)`$ of a reduced rational number $`p/q`$. A comprehensive source of information about the Dedekind symbol is \[KM\], where (among other things) the well-definededness and equivalence of the definitions below is discussed.
###### Definition 1.5.
The Dedekind symbol $`S(p/q)`$ of a reduced rational number $`p/q`$ is defined by the properties
$$S(x)=S(x),S(x+1)=S(x)\text{and}S\left(\frac{p}{q}\right)+S\left(\frac{q}{p}\right)=\frac{p}{q}+\frac{q}{p}+\frac{1}{pq}3\mathrm{sign}(pq).$$
Equivalently, $`S(p/q)`$ is defined by its relation $`S(p/q)=12\mathrm{sign}(q)s(p,q)`$ with the Dedekind sum $`s(p,q)`$, which is given by either of two formulas
$$s(p,q):=\underset{k=1}{\overset{|q|1}{}}((\frac{k}{q}))((\frac{kp}{q}))=\frac{1}{4|q|}\underset{k=1}{\overset{|q|1}{}}\mathrm{cot}\left(\frac{k\pi }{q}\right)\mathrm{cot}\left(\frac{kp\pi }{q}\right),$$
where $`((x))`$ is the sawtooth function $`((x)):=xx1/2`$. See Table 1 on page 1.
###### Definition 1.6.
Let $`L`$ be a rationally framed link, with framing $`f_i=p_i/q_i`$ on the component $`x_i`$ (measured relative to the $`0`$ framing), and let $`L^0`$ be $`L`$ with all framings replaced by $`0`$. Set
$$Z(L):=Z(L^0)\underset{i}{}\chi \widehat{\mathrm{\Omega }}_{x_i}\left(\mathrm{\Omega }_{x_i}^1\mathrm{\Omega }_{x_i/q_i}\right)\mathrm{exp}\left(\frac{f_i}{2}{}_{x_i}{}^{}+\frac{S(f_i)}{48}\theta \right),$$
where in this equation:
* $`\mathrm{\Omega }^1`$ refers to inversion with respect to the disjoint union product of $`𝒜()`$.
* $`\mathrm{\Omega }_x^1`$ denotes $`\mathrm{\Omega }^1`$ with all of its univalent vertices (“legs”) colored $`x`$.
* $`\mathrm{\Omega }_{x/q}`$ denotes $`\mathrm{\Omega }`$ with all of its legs colored by $`x/q`$ (meaning that terms with $`k`$ legs get multiplied by $`1/q^k`$).
* $`\theta `$ denotes the trivalent diagram $`\begin{array}{c}\text{ }\text{}\text{ }\end{array}`$ in $`𝒜(\mathrm{})`$ and $`_{x_i}`$ denotes the “isolated chord” diagram in $`𝒜(_{x_i})`$.
Notice that if all the $`f_i`$’s are integers, then $`\mathrm{\Omega }_{x_i}^1\mathrm{\Omega }_{x_i/q_i}=1`$ and $`S(f_i)=0`$, and thus $`Z(L)=Z(L^0)_ie^{f_i{}_{x_i}{}^{}/2}`$, and hence the new definition of $`Z`$ extends the standard definition of the Kontsevich integral of an integrally framed link.
We can finally state our main theorem.
###### Theorem 1.
(Proof on page 4.1, alternative formulation in Theorem 6). Let $`L`$ be a rationally framed link and let $`M=S_L^3`$ be the rational homology 3-sphere obtained from $`S^3`$ by surgery along $`L`$. Then Equations (1)–(3) continue to hold for the computation of $`\widehat{Z}^{\mathrm{LMO}}(M)`$, only using Definition 1.6 for the definition of $`Z(L)`$.
###### Remark 1.7.
The LMO invariant can be generalized to be an invariant of 3-manifolds with an embedded link, and the surgery formula for the (thus generalized) LMO invariant holds even if the base manifold is not necessarily $`S^3`$. As our proof of Theorem 1 is completely local (see below), the theorem generalizes in the obvious manner to the case when the base manifold is arbitrary and some passive (non-surgery) link is also present.
### 1.2. Plan of the proof
It is well known (e.g. \[Ro, Section 9H\]) that rational surgery with parameter $`p/q`$ over a link component can be achieved by shackling that component with a framed Hopf chain and then performing integral surgery, in which the framings $`a_{1,\mathrm{},\mathrm{}}`$ are related to $`p/q`$ via a continued fraction expansion (cf. \[KM, Equation 0.5\]):
(6)
$$\begin{array}{c}\text{ }\text{}\text{ }\end{array}\frac{p}{q}=a_1,\mathrm{},a_{\mathrm{}}:=\frac{1}{a_1{\displaystyle \frac{1}{a_2\mathrm{}{\displaystyle \frac{1}{a_\mathrm{}1\frac{1}{a_{\mathrm{}}}}}}}}$$
Thus we compute the LMO invariant of $`S_L^3`$ by introducing an integrally framed shackled version $`L^s`$ of $`L`$, in which all of the original components of $`L`$ are shackled as above, and then by computing $`\widehat{Z}^{\mathrm{LMO}}(S_{L^s}^3)`$. We use a lemma, Lemma 4.1 below, that allows us to compute the integral in the definition of $`\widehat{Z}^{\mathrm{LMO}}(S_{L^s}^3)`$ in two steps: first we integrate over the Hopf chains, and then over the original components of $`L^0`$. The first step is computational in nature, rather complicated in the technical sense and takes up the bulk of the proof. When the output of the first step is appropriately normalized, it turns out to be $`Z(L)`$, which then needs to be fed into the second step, which is nothing but a re-run of the procedure in Equations (1)–(3). This proves Theorem 1.
### 1.3. Plan of the paper
In Section 2 we discuss some preliminaries: the Wheels and Wheeling theorems, formal Gaussian integration and assorted facts regarding continued fractions and matrices. In Section 3 we prove some necessary lemmas, and in Section 4 we carry out the computation mentioned above, of the Kontsevich integral and of surgery over Hopf chains, and thus finish the proof of Theorem 1.
Section 5 contains some further computations. Most importantly, we compute the LMO invariant of arbitrary lens spaces and of certain Seifert fibered spaces, and find that the LMO invariant does not separate lens spaces, is far from separating general Seifert fibered spaces, but it does separate Seifert fibered spaces which are integral homology spheres.
### 1.4. Acknowledgement
We wish to thank R. Kirby, S. Levy, J. Lieberum, L. Rozansky, H. Rubinstein, D. P. Thurston and B. Weiss for their remarks, suggestions and ideas. In particular we wish to thank D. P. Thurston for his help with the argument in Proposition 2.2. The first author’s research at MSRI was supported in part by NSF grant DMS-9701755. The second author’s research at the Hebrew University was supported in part by a Guastella Fellowship.
## 2. Preliminaries
### 2.1. Wheels and Wheeling
Let us start with the statements of the two fundamental theorems, Wheels and Wheeling.
###### Theorem 2 (Wheels, \[BGRT, BLT, Th\]).
The Kontsevich integral of the unknot, $`\nu =Z()`$, is the symmetrization of the Wheels element $`\mathrm{\Omega }`$ of Equation (4): $`\nu =\chi \mathrm{\Omega }`$. $`\mathrm{}`$
###### Theorem 3 (Wheeling, \[BGRT, Ko, Moc, BLT, Th\]).
The map $`\chi \widehat{\mathrm{\Omega }}_x:𝒜(_x)𝒜(_x)`$ is an algebra isomorphism, so that
(7)
$$\chi \widehat{\mathrm{\Omega }}_x(A\stackrel{\Gamma }{}B)=\chi \widehat{\mathrm{\Omega }}_x(A)\mathrm{\#}\chi \widehat{\mathrm{\Omega }}_x(B),\text{for }A,B𝒜(_x).$$
$`\mathrm{}`$
We will need a slightly more general form of the Wheeling theorem, also proven in \[BLT, Th\]:
Theorem 3’. Equation 7 holds even if $`A`$ and $`B`$ are allowed to have skeleton components and other univalent vertices beyond those labeled $`x`$, provided $`A`$ and $`B`$ are both “invariant with respect to $`x`$”, meaning that they satisfy the relation in Figure 1$`\mathrm{}`$
We note that Theorems 23’ have a common generalization as explained and proven in \[BLT, Th\] (see also Section 4.2 of this article):
###### Theorem 4.
Let $`{}_{}{}^{y}{}_{x}{}^{}`$ denote the open Hopf link with the open component labeled $`x`$ and the closed component labeled $`y`$. Let $`𝒜(_x_y)`$ denote the space of diagrams that have one directed interval skeleton component marked $`x`$ and (possibly) a number of univalent vertices marked $`y`$, modulo the STU, AS and IHX relations and modulo $`y`$-flavored link relations as in \[Å-II, Section 5.2\]. Then in $`𝒜(_x_y)`$,
$$\sigma _yZ({}_{x}{}^{y})=\mathrm{\Omega }_y\stackrel{\Gamma }{}\mathrm{exp}_\mathrm{\#}({}_{x}{}^{y})=(1+\frac{1}{48}yy+\mathrm{})(_x+{}_{x}{}^{y}+\frac{1}{2}{}_{x}{}^{y}\mathrm{\#}{}_{x}{}^{y}+\mathrm{})$$
### 2.2. Formal Gaussian integration
###### Definition 2.1.
Formal Gaussian integration is defined on formal Gaussian-like expressions (“perturbed $`\mathrm{\Lambda }`$-Gaussians”, for some invertible “covariance” $`X\times X`$ matrix $`\mathrm{\Lambda }=(l_{ij})`$) in $`𝒜(_XE)`$ ($`E`$ denotes some arbitrary extra skeleton components) by
(8)
$$^{FG}P\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)𝑑X:=\mathrm{exp}\left(\frac{1}{2}l^{ij}\begin{array}{c}_{x_i}_{x_j}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right),P_X=\mathrm{exp}\left(\frac{1}{2}l^{ij}\begin{array}{c}_{_{x_i}}_{_{x_j}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right)P|_{X0},$$
where $`\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}`$ ($`\begin{array}{c}_{_{x_i}}_{_{x_j}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}`$) denotes the $`x_ix_j`$-colored ($`_{x_i}_{x_j}`$-colored) strut, where $`P`$ is $`X`$-substantial (involves no struts both of whose ends are colored by colors in $`X`$, see \[Å-II, Definition 2.7\]), and where $`(l^{ij})`$ is the inverse matrix of $`\mathrm{\Lambda }=l_{ij}`$. (Here and below we employ the Einstein summation convention).
Below we will often need to compute the Kontsevich integral of links, and formal Gaussian integrals thereof. The Kontsevich integral of a link is valued in a quotient space $`𝒜(_X)`$ of $`𝒜(_X)`$, and thus it is useful to note that formal Gaussian integration is well defined for integrands in $`𝒜(_X)`$:
###### Proposition 2.2.
Formal Gaussian integration is well defined on $`𝒜(_XE)`$. Specifically, if $`Z𝒜(_XE)`$ can be written as a perturbed $`\mathrm{\Lambda }`$-Gaussian in two ways,
$$Z=P_1\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)=P_2\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right),$$
where $`P_{1,2}`$ are $`X`$-substantial and in $`𝒜(_XE)`$ but the equality holds in $`𝒜(_XE)`$, then
$$^{FG}𝑑XP_1\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)=^{FG}𝑑XP_2\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)$$
in $`𝒜(E)`$.
###### Proof.
We argue by induction on $`|X|`$. If $`X`$ is empty, so is the statement of the proposition. Otherwise there are two cases.
The lucky case: If $`l_{11}0`$ we can compute the $`x_1`$ integral first (\[Å-II, Proposition 2.13\]), and we need to show that
(9)
$$^{FG}d(X\backslash \{x_1\})^{FG}𝑑x_1(P_1P_2)\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)=0,$$
where we know that $`(P_1P_2)\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)`$ is link relation equivalent to $`0`$ via $`X`$-flavored link relations. Multiplication by $`\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_1}^{x_1}\hfill \end{array}`$ is well defined modulo link relations (exercise!), and thus
$$P:=(P_1P_2)\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)\mathrm{exp}\left(\frac{1}{2}l_{11}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_1}^{x_1}\hfill \end{array}\right)$$
is also link relation equivalent to $`0`$ via $`X`$-flavored link relations. By the definition of formal Gaussian integration the inner integral in Equation (9) is given by
$$P,\mathrm{exp}\left(\frac{1}{2}l_{11}^1\begin{array}{c}_{x_1}_{x_1}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right).$$
The map $`DD,\mathrm{exp}\left(\frac{1}{2}l_{11}^1\begin{array}{c}_{x_1}_{x_1}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right)`$ kills all $`x_1`$-flavored link relations (because the $`x_1`$-marked strut $`\begin{array}{c}_{x_1}_{x_1}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}`$ is $`x_1`$-invariant), and maps $`(X\backslash \{x_1\})`$-flavored link relations to $`(X\backslash \{x_1\})`$-flavored link relations. Therefore the inner integral in Equation (9) is link equivalent to $`0`$ via $`(X\backslash \{x_1\})`$-flavored link relations. By the induction hypothesis we now find that the outer integral in Equation (9) vanishes.
The ugly case: If $`l_{11}=0`$ we consider
$$I(ϵ):=^{FG}𝑑X(P_1P_2)\mathrm{exp}\left(\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)\mathrm{exp}ϵ\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_1}^{x_1}\hfill \end{array},$$
where $`ϵ`$ is an arbitrary scalar. Multiplication by $`\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_1}^{x_1}\hfill \end{array}`$ is well defined modulo link relations, and so the integrand here remains link equivalent to $`0`$. Thus by the lucky case, $`I(ϵ)`$ vanishes for all $`ϵ0`$. On the other hand, the coefficient of every diagram that appears in $`I(ϵ)`$ is a rational function in $`ϵ`$ which is non singular at $`ϵ=0`$ because $`\mathrm{\Lambda }`$ is regular. Thus it must be that $`I(0)=0`$. $`\mathrm{}`$
###### Remark 2.3.
The above proof gives us an opportunity to whine and complain about the state of our understanding of integration in spaces of diagrams. Two such integration theories exist. The LMO integration theory as in \[LMO, Le\], and the formal Gaussian theory of \[Å-I–III\]. Neither one of them is satisfactory:
* Formal Gaussian integration has a solid conceptual foundation; it is the diagrammatic analogue of an old time favorite, the theory of Gaussian integrals and Feynman diagrams. But it is defined only for a restricted kind of integrands, and thus it can only be used to define an invariant of a restricted class of 3-manifolds, namely rational homology spheres. And certain aspects of it, such as its being well-defined modulo link relations (as above), are somewhat tricky.
* The LMO integration theory is always defined and there’s no problem showing that it is well-defined modulo link relations, and thus it is superior to formal Gaussian integration, at least in some sense. But we (the authors) lack a conceptual understanding of what it really means. It involves the introduction out of thin air of some strange relations, and one needs to spend some time verifying that under these new relations the theory does not collapse to nothing. In \[LMO, Le\], these relations and the construction as a whole are not interpreted. In \[Å-III\] the relations are given a semi-satisfactory interpretation in terms of “negative dimensions”. But whereas formal Gaussian integration is clearly the diagrammatic counterpart of Gaussian integration over Lie algebras, the LMO integration theory is not fully understood as the diagrammatic counterpart of anything (be it integration or anything else). The situation clearly needs to be resolved. Is the LMO theory a diagrammatization of something known? What is it? If not, then it is a genuinely new piece of mathematics. Genuinely new mathematics is wonderful, but it is a rare commodity. Is the LMO theory really new?
In \[Å-III\] it is shown that the two integration theories agree whenever the weaker one is defined. Thus we could have deduced the previous proposition from the corresponding one for the LMO theory, which is easier to prove. But then we would have had to rely on a non trivial theory which is not yet properly understood.
### 2.3. Surgery, continued fractions and matrices
The continued fraction expansion in Equation (6) has a slightly refined version
(10)
$$\left(\begin{array}{c}p\\ q\end{array}\right)=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}a_1& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}a_2& 1\\ 1& 0\end{array}\right)\mathrm{}\left(\begin{array}{cc}a_{\mathrm{}}& 1\\ 1& 0\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right),$$
in which the signs of $`p`$ and $`q`$ are uniquely determined from $`a_{1,\mathrm{},\mathrm{}}`$. For convenience and without loss of generality, we assume that the signs of $`p`$ and $`q`$ are fixed so that (10) holds. This done, we define the integers $`u`$ and $`v`$ from the equality
$$\left(\begin{array}{cc}p& u\\ q& v\end{array}\right)=A(a_1,\mathrm{},a_{\mathrm{}}):=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}a_1& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}a_2& 1\\ 1& 0\end{array}\right)\mathrm{}\left(\begin{array}{cc}a_{\mathrm{}}& 1\\ 1& 0\end{array}\right).$$
(cf. \[KM, Lemma 1.9\]).
Let $`\mathrm{\Lambda }=(l_{ij})`$ be the (tri-diagonal) linking matrix of the Hopf chain of Equation (6) (cf. \[KM, Equation 0.6\]),
$$\mathrm{\Lambda }=\mathrm{\Lambda }(a_1,\mathrm{},a_{\mathrm{}}):=\left(\begin{array}{ccccc}a_1& 1& & & \\ 1& a_2& 1\\ & 1& \mathrm{}& & & \\ & & & \mathrm{}& 1\\ & & & 1& a_{\mathrm{}}\end{array}\right).$$
###### Proposition 2.4.
The four corners of the inverse matrix $`\mathrm{\Lambda }^1=(l^{ij})`$ of $`\mathrm{\Lambda }`$ are given by
(11)
$$l^{11}=\frac{p}{q},l^1\mathrm{}=l^\mathrm{}1=\frac{(1)^{\mathrm{}}}{q},l^{\mathrm{}\mathrm{}}=\frac{v}{q}.$$
###### Proof.
By induction on $`\mathrm{}`$ (and row expansion of the relevant determinants) one establishes the equality
$$A(a_1,\mathrm{},a_{\mathrm{}})=\left(\begin{array}{cc}det\mathrm{\Lambda }(a_2,\mathrm{},a_{\mathrm{}})& det\mathrm{\Lambda }(a_2,\mathrm{},a_\mathrm{}1)\\ det\mathrm{\Lambda }(a_1,\mathrm{},a_{\mathrm{}})& det\mathrm{\Lambda }(a_1,\mathrm{},a_\mathrm{}1)\end{array}\right).$$
After that, the theorem follows from simple observations regarding the determinant $`det\mathrm{\Lambda }`$ and the minors $`\mathrm{\Lambda }^{(ij)}`$ of the matrix $`\mathrm{\Lambda }`$. Namely, that $`det\mathrm{\Lambda }=det\mathrm{\Lambda }(a_1,\mathrm{},a_{\mathrm{}})`$, that $`\mathrm{\Lambda }^{(11)}`$ is $`det\mathrm{\Lambda }(a_2,\mathrm{},a_{\mathrm{}})`$, that $`\mathrm{\Lambda }^{(1\mathrm{})}`$ is triangular with ones on the diagonal, and that $`\mathrm{\Lambda }^{(\mathrm{}\mathrm{})}`$ is $`det\mathrm{\Lambda }(a_1,\mathrm{},a_\mathrm{}1)`$. $`\mathrm{}`$
We note that while the matrix $`\mathrm{\Lambda }`$ is not determined by $`p/q`$, a certain combination of its trace $`\tau `$, its signature $`\varsigma `$ and the numbers $`p/q`$ and $`v/q`$ appearing in the corners of its inverse matrix does depend only on $`p/q`$:
###### Theorem 5.
(cf. \[KM, Equation 0.8\]) The Dedekind symbol $`S(p/q)`$ of $`p/q`$ (see Table 1) is given by
$$S(p/q)=3\varsigma \tau l^{11}l^{\mathrm{}\mathrm{}}=3\varsigma \tau +\frac{p+v}{q}.$$
## 3. Some Lemmas and preliminary computations
### 3.1. Simple diagram manipulations
In the previous section we have used elements of $`𝒜(_X)`$ where $`X`$ is a set of labels. One may extend this notation in a multilinear way so as to allow labels which are elements of a formal vector space whose basis is indexed by $`X`$. Thinking of the elements of $`𝒜(_X)`$ as multinomials in the symbols $`X`$, whose coefficients are diagrams, one can also perform evaluations in which one or more of the symbols in $`X`$ are set to zero, which has the effect of selecting out those diagrams not mentioning these symbols.
###### Lemma 3.1.
For all $`A,B𝒜(_X)`$, $`_A(B)=A_{\overline{X}},B_{X+\overline{X}}_{\overline{X}}`$. Here $`\overline{X}`$ denotes a set in 1–1 correspondence to $`X`$, $`X+\overline{X}`$ is the set of corresponding sums, $`A_{\overline{X}}`$ denotes the same diagram as $`A`$ with all labels changed from $`X`$ labels to their corresponding $`\overline{X}`$-label, while the pairing takes place over the set of labels $`\overline{X}`$. $`\mathrm{}`$
###### Lemma 3.2.
Let $`E`$ be some skeleton specification (a list of colored $``$’s, $``$’s, $``$’s and $``$’s) not containing the colors $`t`$, $`x`$ and $`y`$. For any $`D𝒜(_t)`$ and $`A𝒜(_tE)`$,
$$\left[\widehat{D}_tA\right]_{tx+y}=\widehat{D}_x\left[A\right]_{tx+y}=\widehat{D}_y\left[A\right]_{tx+y}.$$
### 3.2. Properties of $`\mathrm{\Omega }`$
###### Proposition 3.3 (Pseudo-linearity of $`\mathrm{log}\mathrm{\Omega }`$).
For any $`x`$-invariant diagram $`D`$,
(12)
$$_D\mathrm{\Omega }_x=_D\mathrm{\Omega }_x|_{x=0}\mathrm{\Omega }_x=D,\mathrm{\Omega }_x_x\mathrm{\Omega }_x$$
(compare with standard calculus: if $`D`$ is any differential operator and $`f`$ is a linear function, then $`De^f=(Df)(0)e^f`$; we added the prefix “pseudo” above because Equation 12 does not hold for every $`D`$, but only for $`x`$-invariant $`D`$’s).
###### Proof.
In \[BLT, Th\] it is shown that $`\mathrm{\Omega }_{x+y}=\mathrm{\Omega }_x\stackrel{\Gamma }{}\mathrm{\Omega }_y`$ in $`𝒜(_x_y)`$, and thus using Lemma 3.1,
$$_D(\mathrm{\Omega })=D_y,\mathrm{\Omega }_{x+y}_y=D_y,\mathrm{\Omega }_x\stackrel{\Gamma }{}\mathrm{\Omega }_y_y=D_y,\mathrm{\Omega }_y_y\mathrm{\Omega }_x=D_x,\mathrm{\Omega }_x_x\mathrm{\Omega }_x.$$
The $`x`$-invariance of $`D`$ is used in asserting the equality between the contractions $`D_y,\mathrm{\Omega }_{x+y}_y`$ and $`D_y,\mathrm{\Omega }_x\stackrel{\Gamma }{}\mathrm{\Omega }_y_y`$. If $`x`$-invariance is not assumed, the contraction map $`D_y,_y:𝒜(_x_y)𝒜(_x)`$ may not descend to a map $`𝒜(_x_y)𝒜(_x)`$. $`\mathrm{}`$
###### Corollary 3.4.
For any $`D𝒜()`$ we have $`e^_D\mathrm{\Omega }=e^{D,\mathrm{\Omega }}\mathrm{\Omega }`$. In particular, for any $`\alpha `$,
$$\mathrm{exp}\left(\frac{\alpha }{2}\begin{array}{c}__x__x\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right)\mathrm{\Omega }_x=\mathrm{exp}\left(\frac{\alpha \theta }{48}\right)\mathrm{\Omega }_x$$
###### Proof.
The first assertion follows immediately from Proposition 3.3. The second assertion follows from the first and from the equality $`\begin{array}{c}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array},\mathrm{\Omega }=\theta /24`$, which follows from the fact that the two-legged part of $`\mathrm{\Omega }`$ is $`/48`$. $`\mathrm{}`$
###### Corollary 3.5.
For any $`A,B𝒜()`$,
$$A\stackrel{\Gamma }{}B,\mathrm{\Omega }=A,\mathrm{\Omega }B,\mathrm{\Omega }.$$
In particular, if $`A`$ is invertible, then $`A^1,\mathrm{\Omega }=A,\mathrm{\Omega }^1`$.
###### Proof.
Clearly, $`_{A\stackrel{\Gamma }{}B}=_A_B`$. And so, using Proposition 3.3 twice and the fact that $`\mathrm{\Omega }|_{x=0}=1`$, we get
$`A\stackrel{\Gamma }{}B,\mathrm{\Omega }`$ $`=_{A\stackrel{\Gamma }{}B}\mathrm{\Omega }|_{x=0}=_A(_B(\mathrm{\Omega }))|_{x=0}`$
$`=_A(\mathrm{\Omega })|_{x=0}B,\mathrm{\Omega }=\mathrm{\Omega }|_{x=0}A,\mathrm{\Omega }B,\mathrm{\Omega }=A,\mathrm{\Omega }B,\mathrm{\Omega }.`$
###### Lemma 3.6.
The following two formulas hold when observed<sup>2</sup><sup>2</sup>2See \[Me\].:
(17) $`\widehat{\mathrm{\Omega }}\mathrm{exp}\left({\displaystyle \frac{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ \hfill \end{array}}{2}}\right)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\right)\mathrm{\Omega }\mathrm{exp}\left({\displaystyle \frac{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ \hfill \end{array}}{2}}\right)`$
(22) $`\widehat{\mathrm{\Omega }}_x\mathrm{exp}\left({\displaystyle \frac{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x+y}^{x+y}\hfill \end{array}}{2}}\right)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\right)\mathrm{\Omega }_{x+y}\mathrm{exp}\left({\displaystyle \frac{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x+y}^{x+y}\hfill \end{array}}{2}}\right)`$
###### Proof.
Equation (17) follows from the following computation:
$`\widehat{\mathrm{\Omega }}_x\mathrm{exp}\left({\displaystyle \frac{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}}{2}}\right)`$ $`=\mathrm{exp}\left({\displaystyle \frac{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x+y}^{x+y}\hfill \end{array}}{2}}\right),\mathrm{\Omega }_y_y`$ by Lemma 3.1
$`=e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}/2}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^y^y\hfill \end{array}/2},\mathrm{\Omega }_y_y=e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}/2}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^y^y\hfill \end{array}/2},\mathrm{\Omega }_y_y`$
$`=e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}/2}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}},e^{\begin{array}{c}__y__y\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}/2}\mathrm{\Omega }_y_y`$ see Remark 1.4
$`=\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\right)e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}/2}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}},\mathrm{\Omega }_y_y`$ by Corollary 3.4,
$`=\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\right)e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}/2}\mathrm{\Omega }_x`$
The operator $`\widehat{\mathrm{\Omega }}`$ commutes with translation by $`y`$ (the map $`xx+y`$), because $`\widehat{\mathrm{\Omega }}`$ has “constant coefficients”. Thus Equation (22) is a consequence of Equation (17). Alternatively, it can be proven directly along the same lines. $`\mathrm{}`$
### 3.3. Properties of $`Z`$
The Wheeling Theorem allows us to introduce a variant $`Z`$ of the Kontsevich integral $`Z`$ which has some nice multiplicative properties:
###### Definition 3.7.
If $`L`$ is an $`X`$-marked link (or tangle), set $`Z:=\left(_i\widehat{\mathrm{\Omega }}_{x_i}^1\sigma _{x_i}\right)Z`$ to be the “wheeled Kontsevich integral”. As $`\sigma `$ and $`\widehat{\mathrm{\Omega }}^1`$ are invertible, $`Z`$ carries just as much information as the original Kontsevich integral. In a completely parallel manner, set $`\stackrel{ˇ}{Z}:=\left(_i\widehat{\mathrm{\Omega }}_{x_i}^1\sigma _{x_i}\right)\stackrel{ˇ}{Z}`$.
A first example of a nice multiplicative property of $`Z`$ is the following lemma:
###### Lemma 3.8.
If $`L`$ is a (possibly rationally) framed link and $`L^0`$ is the corresponding $`0`$-framed link, then
$$Z(L)=Z(L^0)\underset{i}{}\mathrm{exp}\left((S(f_i)f_i)\frac{\theta }{48}\right)\mathrm{\Omega }_{x_i}^1\mathrm{\Omega }_{x_i/q_i}\mathrm{exp}\frac{f_i}{2}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_i}\hfill \end{array}.$$
In particular, if $`L`$ is integrally framed, then simply
$$Z(L)=Z(L^0)\underset{i}{}\mathrm{exp}\left(\frac{f_i\theta }{48}\right)\mathrm{exp}\frac{f_i}{2}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_i}\hfill \end{array}.$$
###### Proof.
Indeed,
$`Z(L)`$ $`=\left({\displaystyle \underset{i}{}}\widehat{\mathrm{\Omega }}_{x_i}^1\sigma _{x_i}\right)Z(L)`$
$`=Z(L^0){\displaystyle \underset{i}{}}\mathrm{\Omega }_{x_i}^1\mathrm{\Omega }_{x_i/q_i}\mathrm{exp}\left({\displaystyle \frac{f_i}{2}}\widehat{\mathrm{\Omega }}_{x_i}^1\sigma _{x_i}{}_{x_i}{}^{}+{\displaystyle \frac{S(f_i)}{48}}\theta \right)`$
$`=Z(L^0){\displaystyle \underset{i}{}}\mathrm{\Omega }_{x_i}^1\mathrm{\Omega }_{x_i/q_i}\mathrm{exp}\left({\displaystyle \frac{f_i}{2}}\left(\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_i}\hfill \end{array}{\displaystyle \frac{\theta }{24}}\right)+{\displaystyle \frac{S(f_i)}{48}}\theta \right)`$
$`\mathrm{}`$
Another nice multiplicative property of $`Z`$ is its behavior under the operation of taking the connected sum of links:
###### Lemma 3.9.
Let $`L_i`$ be an $`X_i`$-marked link and let $`x_i`$ be some specific colors in $`X_i`$, for $`i=1,2`$, with $`X_i`$ disjoint sets of colors. Let $`t`$ be a color not in $`X_1X_2`$, and let $`L`$ be the connected sum of $`L_1`$ and $`L_2`$ along the components $`x_1`$ and $`x_2`$, with the “merged” component marked $`t`$. Then
(23) $`Z(L)`$ $`=`$ $`Z(_t)^1\left(Z(L_1)/x_1t\right)\left(Z(L_2)/x_2t\right)`$
(24) $`=`$ $`\mathrm{\Omega },\mathrm{\Omega }\mathrm{\Omega }_t^1\left(Z(L_1)/x_1t\right)\left(Z(L_2)/x_2t\right).`$
###### Proof.
Equation (23) follows immediately from the generalized form of the Wheeling Theorem, Theorem 3’, and the known multiplicative property of the Kontsevich integral $`Z`$. Equation (24) follows from Equation (23) using the following lemma, which is of independent interest. $`\mathrm{}`$
###### Lemma 3.10.
The wheeled Kontsevich integral of the 0-framed unknot is given by
(25)
$$Z()=\widehat{\mathrm{\Omega }}^1\sigma \nu =\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{\Omega }.$$
###### Proof.
$`Z()=\widehat{\mathrm{\Omega }}^1\sigma \nu `$ $`=\widehat{\mathrm{\Omega }}^1\mathrm{\Omega }`$ by Wheels
$`=\mathrm{\Omega }^1,\mathrm{\Omega }\mathrm{\Omega },`$ by Proposition 3.3,
$`=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{\Omega },`$ by Corollary 3.5.
###### Corollary 3.11.
If $`L`$ is an $`X`$-marked link, then
$$\stackrel{ˇ}{Z}(L)=\mathrm{\Omega },\mathrm{\Omega }^{|X|}\left(\underset{xX}{}\mathrm{\Omega }_x\right)Z(L).$$
Recall that the integration by parts formula of \[Å-II, Proposition 2.15\] implies that for any $`D,Z𝒜(_x)`$ one has $`\widehat{D}Z𝑑x=Z𝑑x`$ ($`\widehat{D}`$ is divergence-free). In particular, $`\stackrel{ˇ}{Z}𝑑X=\sigma \stackrel{ˇ}{Z}𝑑X`$ and thus:
###### Corollary 3.12.
If $`L`$ is an $`X`$-marked link (possibly rationally framed), then
(26)
$$\text{Å}_0(L)=\mathrm{\Omega },\mathrm{\Omega }^{|X|}\left(\underset{xX}{}\mathrm{\Omega }_x\right)Z(L)𝑑X.$$
$`\mathrm{}`$
###### Exercise 3.13.
Determine the behavior of $`Z`$ under
* doubling a component (see \[LM2\] and Proposition 4.2 below).
* dropping a component.
### 3.4. One specific integral
At several points below we will need the value of a certain specific integral, which we hereby evaluate:
###### Lemma 3.14.
For any scalars $`\alpha `$, $`\beta `$ and $`\gamma `$ the following equality holds:
$$I(\alpha ,\beta ;\gamma ):=\mathrm{\Omega }_{\alpha x}\mathrm{\Omega }_{\beta x}\mathrm{exp}\left(\frac{\gamma }{2}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}\right)𝑑x=\mathrm{exp}\left(\frac{\alpha ^2+\beta ^2}{48\gamma }\theta \right)\mathrm{\Omega }_x,\mathrm{\Omega }_{\alpha \beta x/\gamma }_x$$
###### Proof.
$`I(\alpha ,\beta ;\gamma )`$ $`=\mathrm{exp}\left({\displaystyle \frac{1}{2\gamma }}\begin{array}{c}__u__u\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right)\mathrm{\Omega }_{\alpha u}\mathrm{\Omega }_{\beta u}|_{u=0}`$ by (8), with a dummy $`u`$
$`=\mathrm{exp}\left({\displaystyle \frac{1}{2\gamma }}\begin{array}{c}__u__u\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right)\mathrm{\Omega }_{\alpha (u+v)}\mathrm{\Omega }_{\beta (uv)}|_{u=v=0}`$ introducing a spurious $`v`$
$`=\mathrm{exp}\left({\displaystyle \frac{1}{2\gamma }}\begin{array}{c}_{(\alpha _x+\beta _y)}_{(\alpha _x+\beta _y)}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right)\mathrm{\Omega }_x\mathrm{\Omega }_y|_{x=y=0}`$
$`=e^{\frac{\alpha \beta \begin{array}{c}__x__y\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{\gamma }}e^{\frac{\alpha ^2\begin{array}{c}__x__x\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{2\gamma }}e^{\frac{\beta ^2\begin{array}{c}__y__y\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{2\gamma }}\mathrm{\Omega }_x\mathrm{\Omega }_y|_{x=y=0}`$
$`=e^{\frac{\alpha ^2\begin{array}{c}_x_x\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{2\gamma }},\mathrm{\Omega }_x_xe^{\frac{\beta ^2\begin{array}{c}_y_y\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{2\gamma }},\mathrm{\Omega }_y_ye^{\frac{\alpha \beta \begin{array}{c}__x__y\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{\gamma }}\mathrm{\Omega }_x\mathrm{\Omega }_y|_{x=y=0}`$
$`=\mathrm{exp}\left({\displaystyle \frac{\alpha ^2\theta }{48\gamma }}\right)\mathrm{exp}\left({\displaystyle \frac{\beta ^2\theta }{48\gamma }}\right)\mathrm{\Omega }_x,\mathrm{\Omega }_{\alpha \beta x/\gamma }_x`$ by Corollary 3.4
### 3.5. The framed unknot and an alternative formula for $`\widehat{Z}^{\mathrm{LMO}}`$
We now know enough to write a cleaner formula for $`\widehat{Z}^{\mathrm{LMO}}`$. We start by computing the Kontsevich integral of the $`p`$-framed unknot (for an integer $`p`$). Using Lemmas 3.8 and 3.10 we get
(27)
$$Z(^p)=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\left(\frac{p\theta }{48}\right)\mathrm{exp}\left(\frac{p}{2}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ \hfill \end{array}\right)\mathrm{\Omega }.$$
We can now compute $`\text{Å}_0(^p)`$:
$`\text{Å}_0(^p)`$ $`=\mathrm{\Omega },\mathrm{\Omega }^1{\displaystyle \mathrm{\Omega }_xZ(_x^p)𝑑x}`$ by Corollary 3.12
(30) $`=\mathrm{\Omega },\mathrm{\Omega }^2\mathrm{exp}\left({\displaystyle \frac{p\theta }{48}}\right){\displaystyle \mathrm{\Omega }_x^2\mathrm{exp}\left(\frac{p}{2}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}\right)𝑑x}`$ by Equation (27),
$`=\mathrm{\Omega },\mathrm{\Omega }^2\mathrm{exp}\left({\displaystyle \frac{p\theta }{48}}\right)\mathrm{exp}\left({\displaystyle \frac{\theta }{24p}}\right)\mathrm{\Omega }_x,\mathrm{\Omega }_{x/p}_x`$
(31) $`=\mathrm{\Omega }_x,\mathrm{\Omega }_x^2\mathrm{\Omega }_{x/p}_x\mathrm{exp}\left[\left({\displaystyle \frac{p}{48}}+{\displaystyle \frac{1}{24p}}\right)\theta \right].`$
In particular, the normalization factors in the definition of the Århus integral are
(32)
$$Z_\pm =\text{Å}_0(^{\pm 1})=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\frac{\theta }{16}$$
We can now write a rather clean formula for $`\widehat{Z}^{\mathrm{LMO}}`$ in terms of $`Z`$:
###### Theorem 6.
Let $`L`$ be a framed link, let $`M=S_L^3`$ be the result of surgery on $`L`$, and let $`\varsigma (L)`$ be the signature of its linking matrix.
* If $`L`$ is integrally framed, then
(33)
$$\widehat{Z}^{\mathrm{LMO}}(M)=\mathrm{exp}\left(\frac{\theta }{16}\varsigma (L)\right)Z(L)\underset{i}{}\mathrm{\Omega }_{x_i}dx_i.$$
* If $`L`$ is rationally framed with framing $`f_i=p_i/q_i`$ on the component $`x_i`$ (relative to the $`0`$ framing) and if $`L^0`$ is the zero framed version of $`L`$, then the following surgery formula is equivalent to Theorem 1 (an in particular, as we shall later see, it holds):
(34)
$$\widehat{Z}^{\mathrm{LMO}}(M)=\mathrm{exp}\left(\frac{\theta }{48}\left(3\varsigma (L)+\underset{i}{}S(f_i)f_i\right)\right)Z(L^0)\underset{i}{}e^{f_i\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_i}\hfill \end{array}/2}\mathrm{\Omega }_{x_i/q_i}dx_i.$$
###### Proof.
The first assertion, Equation (33) is a simple assembly of Equations (3), (26) and (32). The second assertion follows by applying $`\widehat{\mathrm{\Omega }}^1\sigma `$ to the definition of $`Z(L)`$ (Definition 1.6), using Wheeling and substituting the result into Equation (33). $`\mathrm{}`$
## 4. The Rational Surgery Formula
### 4.1. Proof modulo computations
Let $`L`$ be a rationally framed link, with framing $`f_i=p_i/q_i`$ on the component $`x_i`$ (measured relative to the $`0`$ framing), and let $`L^0`$ be $`L`$ with all framing replaced by $`0`$. Let $`L^s`$ be a “shackled” version of $`L`$ — an integrally framed link surgery upon which is equivalent to surgery upon $`L`$, as in Section 2.3 and as in Figure 2. Let $`X=\{x_1,\mathrm{},x_n\}`$ be the original set of labels for the components of $`L`$, let $`X^e=x_1^1,\mathrm{},x_1^\mathrm{}_1,\mathrm{},x_n^1,\mathrm{},x_n^\mathrm{}_n`$ be the set of ‘extra’ labels in $`L^s`$, and let $`X^s=X\stackrel{\Gamma }{}X^e`$ be the full set of labels of $`L^s`$. Let $`L^e`$ be the part of $`L^s`$ colored by $`X^e`$ (the Hopf chains). Let $`M=S_L^3=S_{L^s}^3`$.
In order to prove Theorem 1, we compute $`\widehat{Z}^{\mathrm{LMO}}(M)`$ in two steps:
1. We first run the procedure of Equations (1)–(3) only on the ‘extra’ components of $`L^s`$. Namely, we set
(35) $`\stackrel{ˇ}{Z}_1`$ $`:=\nu ^{X^e}Z(L^s),`$ $`\text{Å}_1`$ $`:={\displaystyle \sigma _{X^e}\stackrel{ˇ}{Z}_1𝑑X^e},`$ $`\widehat{Z}_1^{\mathrm{LMO}}`$ $`:=Z_+^{\varsigma _{1+}}Z_{}^{\varsigma _1}\text{Å}_1`$
(here $`\varsigma _{1\pm }`$ are the numbers of positive/negative eigenvalues of the covariance matrix of $`Z(L^s)`$ with respect to only the variables in $`X^e`$).
2. We then run the procedure of Equations (1)–(3) over the remaining components, taking as our input the result of the first step. Namely, we set
$`\stackrel{ˇ}{Z}_2`$ $`:=\nu ^X\widehat{Z}_1^{\mathrm{LMO}},`$ $`\text{Å}_2`$ $`:={\displaystyle \sigma _X\stackrel{ˇ}{Z}_2𝑑X},`$ $`\widehat{Z}_2^{\mathrm{LMO}}`$ $`:=Z_+^{\varsigma _{2+}}Z_{}^{\varsigma _2}\text{Å}_2`$
(here $`\varsigma _{1\pm }`$ are the numbers of positive/negative eigenvalues of the covariance matrix of $`\widehat{Z}_1^{\mathrm{LMO}}`$).
###### Proof of Theorem 1.
Theorem 1 follows immediately from Lemma 4.1 (right below) and Proposition 4.5 (on page 4.5), which assert that $`\widehat{Z}_2^{\mathrm{LMO}}=\widehat{Z}^{\mathrm{LMO}}(M)`$ and that $`\widehat{Z}_1^{\mathrm{LMO}}=Z(L)`$ respectively. $`\mathrm{}`$
###### Lemma 4.1.
The above two step procedure really does compute $`\widehat{Z}^{\mathrm{LMO}}(M)`$. Namely, $`\widehat{Z}_2^{\mathrm{LMO}}=\widehat{Z}^{\mathrm{LMO}}(M)`$.
###### Proof.
This is Theorem 6.6 of \[LMO\]. It also follows from the formalism of \[Å-I–III\] noting that
* Formal Gaussian integration behaves correctly under iteration (\[Å-II, Proposition 2.13\]).
* The covariance matrix of $`\widehat{Z}_1^{\mathrm{LMO}}`$ is the linking matrix of $`L`$ and if $`\varsigma _\pm `$ denotes the numbers of positive/negative eigenvalues of the covariance matrix of $`Z(L^s)`$, then $`\varsigma _\pm =\varsigma _{1\pm }+\varsigma _{2\pm }`$.
$`\mathrm{}`$
### 4.2. The Hopf link and the open Hopf link
As a first step towards understanding the Kontsevich integral of shackled links, we compute $`Z(H(0,0))`$, the wheeled Kontsevich integral of the 0-framed Hopf link:
###### Proposition 4.2.
The wheeled Kontsevich integral of the $`(x,y)`$-marked 0-framed Hopf link $`H_{x,y}(0,0)`$ is given by
$$Z(H_{x,y}(0,0))=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}.$$
###### Proof.
The $`(1,1)`$-framed Hopf link $`H_{x,y}(1,1)`$ is the $`(x,y)`$-marked double $`\mathrm{\Delta }_{x,y}^t_t^1`$ of the 1-framed unknot $`_t^1`$, and hence
$`Z(H_{x,y}(1,1))`$ $`=\widehat{\mathrm{\Omega }}_x^1\widehat{\mathrm{\Omega }}_y^1\sigma _x\sigma _yZ(\mathrm{\Delta }_{x,y}^t_t^1)`$
$`=\widehat{\mathrm{\Omega }}_x^1\widehat{\mathrm{\Omega }}_y^1\left[\sigma _tZ(_t^1)\right]_{tx+y}`$ by \[LM2\],
$`=\widehat{\mathrm{\Omega }}_x^1\left[\widehat{\mathrm{\Omega }}_t^1\sigma _tZ(_t^1)\right]_{tx+y}`$ by Lemma 3.2,
$`=\widehat{\mathrm{\Omega }}_x^1\left[Z(_t^1)\right]_{tx+y}`$
$`=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\right)\widehat{\mathrm{\Omega }}_x^1\left[\mathrm{exp}\left({\displaystyle \frac{1}{2}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x+y}^{x+y}\hfill \end{array}\right)\mathrm{\Omega }_{x+y}\right]`$ by (27),
$`=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\left({\displaystyle \frac{\theta }{24}}\right)\mathrm{exp}\left({\displaystyle \frac{1}{2}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x+y}^{x+y}\hfill \end{array}\right)`$ by Lemma 3.6, (22).
It remains to undo the $`(1,1)`$ framing by using Lemma 3.8 on each component:
$`Z(H_{x,y}(0,0))`$ $`=\left[\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\right)\right]^2\mathrm{exp}\left({\displaystyle \frac{1}{2}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}\right)\mathrm{exp}\left({\displaystyle \frac{1}{2}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^y^y\hfill \end{array}\right)Z(H_{x,y}(1,1))`$
$`=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}.`$
For the sake of completeness, we can now re-prove Theorem 4:
###### Proof of Theorem 4.
To get $`\sigma _x\sigma _yZ({}_{x}{}^{y})`$ from $`Z(H_{x,y}(0,0))`$, we need to apply $`\widehat{\mathrm{\Omega }}_x`$ and $`\widehat{\mathrm{\Omega }}_y`$, and to “open up” the $`x`$-component, which amounts to multiplication (using the product $`\mathrm{\#}`$) by $`Z(_x)^1`$ (recall that the Kontsevich integral of an open unknot is trivial). The latter operation can more easily be performed first, by $`\stackrel{\Gamma }{}`$-multiplying $`Z(H_{x,y}(0,0))`$ by $`\mathrm{\Omega },\mathrm{\Omega }\mathrm{\Omega }_x^1`$, as in (25). Thus
(36)
$$Z({}_{x}{}^{y})=\mathrm{\Omega },\mathrm{\Omega }\mathrm{\Omega }_x^1Z(H_{x,y}(0,0))=\mathrm{\Omega }_x^1\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}$$
and hence
$`\sigma _x\sigma _yZ({}_{x}{}^{y})`$ $`=\widehat{\mathrm{\Omega }}_x\widehat{\mathrm{\Omega }}_yZ({}_{x}{}^{y})`$
$`=\widehat{\mathrm{\Omega }}_x\widehat{\mathrm{\Omega }}_y\left(\mathrm{\Omega }_x^1\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}\right)`$ by Equation (36),
$`=\widehat{\mathrm{\Omega }}_x\left(\mathrm{\Omega }_x^1\widehat{\mathrm{\Omega }}_y\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}\right)=\widehat{\mathrm{\Omega }}_x\left(\mathrm{\Omega }_x^1\mathrm{\Omega }_x\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}\right)`$
$`=\widehat{\mathrm{\Omega }}_x\left(\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}\right)=\mathrm{\Omega }_y\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}.`$
It only remains to note that $`\sigma _yZ({}_{x}{}^{y})=\chi _x\sigma _x\sigma _yZ({}_{x}{}^{y})=\mathrm{\Omega }_y\chi _x\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^y\hfill \end{array}=\mathrm{\Omega }_y\mathrm{exp}_\mathrm{\#}({}_{x}{}^{y})`$. $`\mathrm{}`$
### 4.3. The Hopf chain and the shackling element
We can use the connect-sum lemma (Lemma 3.9) repeatedly in order to compute the wheeled Kontsevich integral of Hopf chains (see Figure 3). The result is
$$Z(H_{x_1,\mathrm{},x_{\mathrm{}}}(0,\mathrm{},0))=\mathrm{\Omega },\mathrm{\Omega }^1\underset{i=2}{\overset{\mathrm{}1}{}}\mathrm{\Omega }_{x_i}^1\underset{i=1}{\overset{\mathrm{}1}{}}\mathrm{exp}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_{i+1}}\hfill \end{array}.$$
Using Lemma 3.8 and some repackaging, we get:
###### Proposition 4.3.
Let $`(l_{ij})`$ be the linking matrix of $`H_{x_1,\mathrm{},x_{\mathrm{}}}(a_1,\mathrm{},a_{\mathrm{}})`$. Then
$$Z(H_{x_1,\mathrm{},x_{\mathrm{}}}(a_1,\mathrm{},a_{\mathrm{}}))=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\left(\frac{\theta a_i}{48}\right)\left(\underset{i=2}{\overset{\mathrm{}1}{}}\mathrm{\Omega }_{x_i}^1\right)\mathrm{exp}\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}.$$
We can now use Proposition 4.3, Equation (36) and the connect-sum lemma (Lemma 3.9) to compute the wheeled Kontsevich integral of the shackling element (see Figure 4):
###### Proposition 4.4.
Let $`(l_{ij})`$ be the $`\mathrm{}\times \mathrm{}`$ linking matrix of the closed components of the shackling element $`S_{x;x_{1..\mathrm{}}}(a_{1..\mathrm{}})`$. Then
$$Z(S_{x;x_{1..\mathrm{}}}(a_{1..\mathrm{}}))=\mathrm{\Omega }_x^1\left(\underset{i=1}{\overset{\mathrm{}1}{}}\mathrm{\Omega }_{x_i}^1\right)\mathrm{exp}\left(\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^{x_1}\hfill \end{array}+\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\frac{\theta a_i}{48}\right).$$
### 4.4. The Kontsevich integral of rationally framed links
###### Proposition 4.5.
The output $`\widehat{Z}_1^{\mathrm{LMO}}`$ of the first step of the computation procedure of Section 4.1 is equal to the Kontsevich integral for rationally framed links $`Z(L)`$ of Definition 1.6: $`\widehat{Z}_1^{\mathrm{LMO}}=Z(L)`$.
###### Proof.
Both sides of the required equality are clearly made of local contributions, one per each shackling element or framing fraction $`p_i/q_i`$. Thus it is enough to prove the proposition at the locale where all the actors act. Ergo we may as well assume that the link $`L`$ in question is a straight line $`_x^f`$ marked $`f=p/q`$, and then, after choosing $`\mathrm{}`$ and $`a_1,\mathrm{},a_{\mathrm{}}`$ as in Equation (6) (more precisely, as in Equation (10)), the shackled $`L`$ becomes the shackling element $`S_{x;x_{1..\mathrm{}}}(a_{1..\mathrm{}})`$ (so $`X^e`$, the set of “extra” labels, is $`\{x_1,\mathrm{},x_{\mathrm{}}\}`$). We just need to compute $`\widehat{Z}_1^{\mathrm{LMO}}`$ as in Section 4.1, starting from $`L^s=S_{x;x_{1..\mathrm{}}}(a_{1..\mathrm{}})`$. Set $`C=\mathrm{\Omega },\mathrm{\Omega }^{\mathrm{}}\mathrm{exp}\left(\frac{\theta {\scriptscriptstyle a_i}}{48}\right)\mathrm{\Omega }_x^1`$ and start crunching:
$`\widehat{\mathrm{\Omega }}_x^1\sigma _x\text{Å}_1`$ $`=\mathrm{\Omega },\mathrm{\Omega }^{\mathrm{}}{\displaystyle Z(S_{x;x_{1..\mathrm{}}}(a_{1..\mathrm{}}))\underset{i}{}\mathrm{\Omega }_{x_i}dx_i}`$
(41) $`=C{\displaystyle \mathrm{\Omega }_x_{\mathrm{}}\mathrm{exp}\left(\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^{x_1}\hfill \end{array}+\frac{1}{2}l_{ij}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^{x_i}^{x_j}\hfill \end{array}\right)\underset{i}{}dx_i}`$ by Proposition 4.4
(46) $`=C\mathrm{exp}\left({\displaystyle \frac{1}{2}}l^{ij}\begin{array}{c}_{x_i}_{x_j}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}\right),\mathrm{\Omega }_x_{\mathrm{}}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^{x_1}\hfill \end{array}}_{X^e}`$ by Definition 2.1
(55) $`=Ce^{\frac{1}{2}l^{11}\begin{array}{c}_{x_1}_{x_1}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}e^{l^1\mathrm{}\begin{array}{c}_{x_1}_x_{\mathrm{}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}e^{\frac{1}{2}l^{\mathrm{}\mathrm{}}\begin{array}{c}_x_{\mathrm{}}_x_{\mathrm{}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}},\mathrm{\Omega }_x_{\mathrm{}}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^{x_1}\hfill \end{array}}_{x_1,x_{\mathrm{}}}`$
(64) $`=Ce^{\frac{p}{2q}\begin{array}{c}_{x_1}_{x_1}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}e^{\frac{(1)^{\mathrm{}}}{q}\begin{array}{c}_{x_1}_x_{\mathrm{}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}e^{\frac{v}{2q}\begin{array}{c}_x_{\mathrm{}}_x_{\mathrm{}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}},\mathrm{\Omega }_x_{\mathrm{}}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^{x_1}\hfill \end{array}}_{x_1,x_{\mathrm{}}}`$ by Equation (11)
(73) $`=Ce^{\frac{p}{2q}\begin{array}{c}_{x_1}_{x_1}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}e^{\frac{(1)^{\mathrm{}}}{q}\begin{array}{c}_{x_1}_x_{\mathrm{}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}},e^{\frac{v}{2q}\begin{array}{c}__x_{\mathrm{}}__x_{\mathrm{}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}\mathrm{\Omega }_x_{\mathrm{}}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^{x_1}\hfill \end{array}}_{x_1,x_{\mathrm{}}}`$
(80) $`=C\mathrm{exp}\left({\displaystyle \frac{v\theta }{48q}}\right)e^{\frac{p}{2q}\begin{array}{c}_{x_1}_{x_1}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}e^{\frac{(1)^{\mathrm{}}}{q}\begin{array}{c}_{x_1}_x_{\mathrm{}}\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}},\mathrm{\Omega }_x_{\mathrm{}}e^{\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^{x_1}\hfill \end{array}}_{x_1,x_{\mathrm{}}}`$ by Corollary 3.4
(83) $`=C\mathrm{exp}\left({\displaystyle \frac{v\theta }{48q}}\right)\mathrm{\Omega }_{x/q}\mathrm{exp}\left({\displaystyle \frac{p}{2q}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}\right)`$
(86) $`=\mathrm{\Omega },\mathrm{\Omega }^{\mathrm{}}\mathrm{exp}\left[\left({\displaystyle \frac{v}{q}}{\displaystyle a_i}\right){\displaystyle \frac{\theta }{48}}\right]\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/q}\mathrm{exp}\left({\displaystyle \frac{p}{2q}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}\right)`$
$`\widehat{\mathrm{\Omega }}^1\sigma \widehat{Z}_1^{\mathrm{LMO}}`$ $`=Z_+^{\varsigma _+}Z_{}^\varsigma _{}\text{Å}_1`$
(89) $`=\mathrm{exp}\left[\left({\displaystyle \frac{v}{q}}+3\varsigma \tau \right){\displaystyle \frac{\theta }{48}}\right]\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/q}\mathrm{exp}\left({\displaystyle \frac{p}{2q}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}\right)`$ by Equation (32)
(92) $`=\mathrm{exp}\left[\left(S(p/q)p/q\right){\displaystyle \frac{\theta }{48}}\right]\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/q}\mathrm{exp}\left({\displaystyle \frac{p}{2q}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}\right)`$ by Theorem 5
(95) $`\widehat{Z}_1^{\mathrm{LMO}}`$ $`=\mathrm{exp}\left[\left(S(p/q)p/q\right){\displaystyle \frac{\theta }{48}}\right]\chi \widehat{\mathrm{\Omega }}\left(\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/q}\right)\mathrm{exp}\left({\displaystyle \frac{p}{2q}}\chi \widehat{\mathrm{\Omega }}(\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array})\right)`$ by Wheeling
$`=\chi \widehat{\mathrm{\Omega }}\left(\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/q}\right)\mathrm{exp}\left({\displaystyle \frac{f}{2}}{}_{x}{}^{}+{\displaystyle \frac{S(f)}{48}}\theta \right)`$
as required. $`\mathrm{}`$
###### Remark 4.6.
Strictly speaking, our arguments in this section assumed that $`\mathrm{}2`$. This is not a serious limitation. One may either read the arguments again and figure what the appropriate defaults are for $`\mathrm{}<2`$, or simply use a longer continued fraction expansion for $`p/q`$. Notice that we made no restrictive assumptions on the integers $`a_i`$, and so every fraction $`p/q`$ has arbitrarily long continued fraction expansions.
## 5. Some Computations
### 5.1. General $`(p,q)`$ lens spaces
As a first application of the rational surgery formula we compute the LMO invariant of general $`(p,q)`$ lens spaces.
###### Proposition 5.1.
The LMO invariant of arbitrary $`(p,q)`$ lens spaces is given by
$$\widehat{Z}^{\mathrm{LMO}}(L_{p,q})=\mathrm{\Omega }_x,\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/p}_x\mathrm{exp}\frac{S(q/p)}{48}\theta $$
###### Proof.
Recall that the $`(p,q)`$ lens space $`L_{p,q}`$ is obtained from $`S^3`$ by surgery over the $`p/q`$-framed unknot $`^{p/q}`$. Thus we can use Equation (34) to compute $`\widehat{Z}^{\mathrm{LMO}}`$, remembering also that by Equation (25), $`Z()=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{\Omega }`$:
$`\widehat{Z}^{\mathrm{LMO}}`$ $`(L_{p,q})=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}(3\mathrm{sign}(p/q)+S(p/q)p/q)\right){\displaystyle e^{p\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}/2q}\mathrm{\Omega }_x\mathrm{\Omega }_{x/q}𝑑x}`$
$`=\mathrm{\Omega }_x,\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/p}_x\mathrm{exp}\left(S\left({\displaystyle \frac{p}{q}}\right){\displaystyle \frac{p}{q}}{\displaystyle \frac{q}{p}}{\displaystyle \frac{1}{pq}}+3\mathrm{sign}(pq)\right){\displaystyle \frac{\theta }{48}}`$
$`=\mathrm{\Omega }_x,\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/p}_x\mathrm{exp}{\displaystyle \frac{S(q/p)}{48}}\theta `$ by Definition 1.5.
$`\mathrm{}`$
###### Corollary 5.2.
The LMO invariant does not separate lens spaces. (Though see Section 5.4 for some more encouraging news about the LMO invariant).
###### Proof.
As noted in \[KM, pp. 247\], the lens spaces $`L_{25,4}`$ and $`L_{25,9}`$ are not homeomorphic but the Dedekind symbols $`S(4/25)=S(9/25)=48/25`$ are equal (see Table 1), and thus their LMO invariants are equal. $`\mathrm{}`$
###### Exercise 5.3.
Recompute $`\widehat{Z}^{\mathrm{LMO}}(L_{p,q})`$ directly from the Kontsevich integral of Hopf chains (Proposition 4.3), using the fact that surgery on a the Hopf chain $`H_{x_1,\mathrm{},x_{\mathrm{}}}(a_1,\mathrm{},a_{\mathrm{}})`$ gives $`L_{q,p}`$.
### 5.2. Seifert fibered spaces with a spherical base
Let $`M=S^3(b;p_1/q_1,\mathrm{},p_n/q_n)`$ be the Seifert fibered space with base space $`S^2`$, with $`n`$ exceptional fibers with orbit invariants $`(p_i,q_i)`$ ($`p_i0`$, $`0<q_i<p_i`$, $`p_i`$ and $`q_i`$ relatively prime), and with bundle invariant $`b`$. The following is easily established (see \[Mon, Sc\]):
* If $`e_0:=b+_iq_i/p_i0`$, then $`M`$ is a rational homology sphere and $`\left|H_1(M)\right|=|e_0|_ip_i`$.
* $`M`$ is given by surgery on the following (rationally framed) “key chain” link in $`S^3`$:
$$L=L(b;p_1/q_1,\mathrm{},p_n/q_n):=\begin{array}{c}\text{ }\text{}\text{ }\end{array}$$
We could use the rational surgery formula of Theorem 1 to compute the LMO invariant of $`M`$, but it is somewhat easier to backtrack a bit, and use the main points of the proof of Theorem 1, applied in a slightly different manner. So we shackle the components $`x_1,\mathrm{},x_n`$ of $`L`$ as in Figure 2, and observe that the result $`L^s`$ is a $`(b)`$-framed unknot marked $`x`$ with $`n`$ one-step-longer Hopf chains hanging from it. We then run the “first step” of Section 4.1 on these longer Hopf chains. As there the result is $`Z(_x^b)`$ corrected by $`n`$ insertions of factors corresponding to the $`n`$ one-step-longer Hopf chains. So noting that $`0,a_i^1,\mathrm{},a_i^{\mathrm{}}=1/a_i^1,\mathrm{},a_i^{\mathrm{}}`$ and using Equation (92) we find that the result of the (revised) first step of the computation is
$`\widehat{\mathrm{\Omega }}^1\sigma Z_1^{}`$ $`=Z(_x^b){\displaystyle \underset{i}{}}\mathrm{exp}\left[\left(S(q_i/p_i)+q_i/p_i\right){\displaystyle \frac{\theta }{48}}\right]\mathrm{\Omega }_x^1\mathrm{\Omega }_{x/p_i}\mathrm{exp}\left({\displaystyle \frac{q_i}{2p_i}}\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}\right)`$
$`=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\left(e_0{\displaystyle \underset{i}{}}S\left({\displaystyle \frac{q_i}{p_i}}\right)\right)\right)\mathrm{\Omega }_x^{1n}\left({\displaystyle \underset{i}{}}\mathrm{\Omega }_{x/p_i}\right)\mathrm{exp}{\displaystyle \frac{e_0\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}}{2}}`$
This result serves as the input into the second step of the computation in the usual manner:
(98) $`\widehat{\mathrm{\Omega }}^1\sigma \stackrel{ˇ}{Z}_1^{}`$ $`=\mathrm{\Omega },\mathrm{\Omega }^2\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\left(e_0{\displaystyle \underset{i}{}}S\left({\displaystyle \frac{q_i}{p_i}}\right)\right)\right)\mathrm{\Omega }_x^{2n}\left({\displaystyle \underset{i}{}}\mathrm{\Omega }_{x/p_i}\right)\mathrm{exp}{\displaystyle \frac{e_0\begin{array}{c}\mathrm{}\text{ }\mathrm{}\hfill \\ ^x^x\hfill \end{array}}{2}}`$
$`\text{Å}_2^{}`$ $`={\displaystyle \widehat{\mathrm{\Omega }}^1\sigma \stackrel{ˇ}{Z}_1^{}𝑑x}`$
(101) $`=\mathrm{\Omega },\mathrm{\Omega }^2\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}\left(e_0{\displaystyle \underset{i}{}}S\left({\displaystyle \frac{q_i}{p_i}}\right)\right)\right)\mathrm{exp}{\displaystyle \frac{\begin{array}{c}_x_x\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{2e_0}},\mathrm{\Omega }_x^{2n}{\displaystyle \underset{i}{}}\mathrm{\Omega }_{x/p_i}`$
(104) $`\widehat{Z}^{\mathrm{LMO}}(M)`$ $`=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\left({\displaystyle \frac{\theta }{48}}(e_03\mathrm{sign}(e_0){\displaystyle \underset{i}{}}S\left({\displaystyle \frac{q_i}{p_i}}\right))\right)\mathrm{exp}{\displaystyle \frac{\begin{array}{c}_x_x\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{2e_0}},\mathrm{\Omega }_x^{2n}{\displaystyle \underset{i}{}}\mathrm{\Omega }_{x/p_i}`$
###### Remark 5.4.
According to \[LMO, Proposition 5.3\] (see also \[LMMO\]), for any rational homology sphere $`M`$ the coefficient of $`\theta `$ in $`\widehat{Z}^{\mathrm{LMO}}(M)`$ is $`\lambda _w(M)/4`$, where $`\lambda _w(M)`$ denotes the Casson-Walker invariant of $`M`$ (see \[Wa\]). Thus in our case, extracting the coefficient of $`\theta `$ in $`\widehat{Z}^{\mathrm{LMO}}(M)`$ from Equation (104), we find that
(105)
$$\lambda _w(S^3(b;p_1/q_1,\mathrm{},p_n/q_n))=\frac{1}{12}\left(e_0+\frac{2n}{e_0}3\mathrm{sign}(e_0)+\underset{i}{}\frac{1}{e_0p_i^2}S\left(\frac{q_i}{p_i}\right)\right).$$
This agrees with Lescop’s evaluation of the same quantity at \[Les, Proposition 6.1.1\] (though notice the different normalization, \[Les, Section 1.5, T5.0\]).
###### Remark 5.5.
We do not know how to simplify the pairing in Equation (104) any further, except in the case when $`n2`$ (using Lemma 3.14). But in this case $`M`$ is a lens space, so we learn nothing new. At any rate, Remark 5.4 allows us to rewrite Equation (104) in an alternate way:
(106)
$$\widehat{Z}^{\mathrm{LMO}}(M)=\mathrm{\Omega },\mathrm{\Omega }^1\mathrm{exp}\frac{\lambda _w(M)\theta }{4}\mathrm{exp}\left(\frac{\theta }{48e_0}(n2\underset{i}{}\frac{1}{p_i^2})\right)\mathrm{exp}\frac{\begin{array}{c}_x_x\hfill \\ \mathrm{}\text{ }\mathrm{}\hfill \end{array}}{2e_0},\mathrm{\Omega }_x^{2n}\underset{i}{}\mathrm{\Omega }_{x/p_i}$$
###### Remark 5.6.
Equation (106) shows that the LMO invariant of Seifert fibered spaces with a spherical base factors through at most $`n+2`$ integer parameters: the $`p_i`$’s, $`e_0`$ and the Casson-Walker invariant of $`M`$. As $`S^3(b;p_1/q_1,\mathrm{},p_n/q_n)`$ depends on $`2n`$ integers (albeit with some range restrictions), it is clear that the LMO invariant grossly fails to separate rational homology Seifert fibered spaces. Some more encouraging news about the LMO invariant are in Section 5.4.
### 5.3. Some $`sl(2)`$ computations
Recall that diagram spaces such as $`𝒜()`$, $`𝒜()`$ and $`𝒜(\mathrm{})`$ are in some sense “universal” versions of certain tensor spaces associated with Lie algebras (see \[B-N\]), and that the Kontsevich integral and the LMO invariant are in this sense universal versions of the Reshetikhin-Turaev invariants of links and of manifolds (see \[Kas, LM1, Le, Å-I\]). Thus we wish to explicitly compare our results with some known answers for the Reshetikhin-Turaev invariants in the case of the Lie algebra $`sl(2)`$.
In the particular case of $`sl(2)`$, we may divide all diagram spaces $`𝒜(\mathrm{})`$ by the additional “$`A_1`$ relations” (with $`\mathrm{}`$ a formal parameter),
$`=3`$ $`\text{and}=2\mathrm{}().`$
(The first relations says that $`sl(2)`$ is three dimensional, and the second becomes the identity $`A\times (B\times C)=(AB)C(AC)B`$ in the vector calculus guise of $`sl(2)`$. See e.g. \[CV\].) Calling the corresponding quotient spaces $`𝒜^1(\mathrm{})`$, we easily find by repeated applications of the $`A_1`$ relations that
$`𝒜^1(\mathrm{})=[[\mathrm{}]]`$ $`\text{and}𝒜^1()=[[\mathrm{},s]],`$
(here $`s=()`$ is a single strut). One easily computes that in $`𝒜^1(\mathrm{})`$, $`\theta =12\mathrm{}`$. Furthermore, simple induction (see e.g. \[Th\]) shows that in $`𝒜^1(\mathrm{})`$,
(107) $`\omega _{2m}=2(2\mathrm{}s)^m`$ $`\text{and}s^m,s^m=(2m+1)!`$
We can now find in $`𝒜^1(\mathrm{})`$:
$`\mathrm{\Omega }`$ $`=\mathrm{exp}_\stackrel{\Gamma }{}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}b_{2m}\omega _{2m}`$ by Equation (4)
$`=\mathrm{exp}2{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}b_{2m}(2\mathrm{}s)^m`$ by Equation (107)
$`={\displaystyle \frac{\mathrm{sinh}\sqrt{\mathrm{}s/2}}{\sqrt{\mathrm{}s/2}}}`$ by Equation (5)
$`={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{}^m}{2^m(2m+1)!}}s^m.`$
Hence by Equation (107),
$$\mathrm{\Omega },\mathrm{\Omega }=\underset{m=0}{\overset{\mathrm{}}{}}\frac{\mathrm{}^{2m}}{2^{2m}(2m+1)!}=\frac{\mathrm{sinh}\mathrm{}/2}{\mathrm{}/2}.$$
Assembling everything, we find from Equation (106) that in $`𝒜^1(\mathrm{})`$ the LMO invariant of $`M=S^3(b;p_1/q_1,\mathrm{},p_n/q_n)`$ is given as a product of an easy factor
$$\frac{\mathrm{}/2}{\mathrm{sinh}\mathrm{}/2}\mathrm{exp}3\mathrm{}\lambda _w(M),$$
and a difficult ‘rest’:
(108)
$$Z^{\text{rest}}(M):=\mathrm{exp}\frac{s}{2e_0},\frac{2_ip_i\mathrm{sinh}\sqrt{\mathrm{}s/2p_i^2}}{\mathrm{}s\left(\mathrm{sinh}\sqrt{\mathrm{}s/2}\right)^{n2}}\mathrm{exp}\frac{\mathrm{}}{4e_0}\left(n2\underset{i}{}\frac{1}{p_i^2}\right).$$
Next, we notice two equalities; the first is an immediate corollary of Equation (107), and the second is an exercise in Gaussian integration (for simplicity we restrict to $`e_0>0`$; otherwise we need to use a different contour for the integration):
$$\mathrm{exp}\frac{s}{2e_0},\left(\frac{\mathrm{}s}{2}\right)^m=\frac{(2m+1)!\mathrm{}^m}{(4e_0)^mm!}=\sqrt{\frac{\mathrm{}^3e_0^3}{16\pi }}_{}𝑑\beta e^{\mathrm{}e_0\beta ^2/4}\beta ^2(\mathrm{}\beta /2)^{2m}.$$
This known for every monomial in $`\mathrm{}s/2`$, we can now replace the bracket in Equation (108) by a Gaussian integral, substituting $`(\mathrm{}\beta /2)^2`$ for every occurrence of $`\mathrm{}s/2`$ in (108). Thus,
(109)
$$Z^{\text{rest}}(M)=\sqrt{\frac{e_0^3}{\pi \mathrm{}}}\mathrm{exp}\frac{\mathrm{}}{4e_0}\left(n2\underset{i}{}\frac{1}{p_i^2}\right)_{}𝑑\beta e^{\mathrm{}e_0\beta ^2/4}\frac{_ip_i\mathrm{sinh}\mathrm{}\beta /2p_i}{\left(\mathrm{sinh}\mathrm{}\beta /2\right)^{n2}}.$$
Disregarding some thorny normalization issues, this result is in agreement with \[LR, Section 4.5\].
### 5.4. The LMO invariant of integral homology Seifert fibered spaces
Corollary 5.2 and Remark 5.6 seem like bad news for finite type invariants of 3-manifolds. But things aren’t as bad as they seem. The LMO invariant is a universal finite type invariant merely over the rationals, and it may well be that rational homology spheres, whose homology groups may contain torsion, are separated by finite-group-valued finite type invariants. At any rate, it is nice that we can complement our non-separation results for the LMO invariant of rational homology spheres with a separation result for (some) integral homology spheres:
###### Theorem 7.
The LMO invariant separates integral homology Seifert fibered spaces.
###### Proof.
To be an integral homology sphere, a Seifert fibered space must be of the form $`M=S^3(b;p_1/q_1,\mathrm{},p_n/q_n)`$ discussed above. Furthermore, it must have $`\left|H_1(M)\right|=|e_0|_ip_i=|b+_iq_i/p_i|_ip_i=1`$. For this to happen, $`e_0`$ must equal $`\pm 1/_ip_i`$, the $`p_i`$’s must be pairwise relatively prime, and then the $`q_i`$’s and $`b`$ are uniquely determined up to an overall sign by the $`p_i`$’s, using the Chinese remainder theorem. The overall sign of $`b`$ and the $`q_i`$’s (and therefore also of $`e_0`$) can be read from the coefficient of $`\theta `$ in $`\widehat{Z}^{\mathrm{LMO}}(M)`$ (see Equation (105)). It remains to check to what extent does $`Z^{\text{rest}}(M)`$ determine the $`p_i`$’s. Thus we regard $`Z^{\text{rest}}(M)`$ as ‘known’, and try to read out $`n`$ and $`p_1,\mathrm{},p_n`$.
For the sake of simplicity we assume that $`e_0>0`$; that is, that $`e_0=+1/_ip_i`$. This done, Equation (109) shows that $`Z^{\text{rest}}(M)`$ makes sense as an honest analytic function of $`\mathrm{}>0`$, and not merely as a formal power series. Assume for the moment that $`n>2`$. For large values of $`\mathrm{}`$ it is easy to bound the integral in Equation (109) above and below by rational functions of $`\mathrm{}`$, and thus the growth rate of $`Z^{\text{rest}}(M)`$ as a function of $`\mathrm{}`$ is determined by the exponential prefactor. Thus by observing the growth rate of $`Z^{\text{rest}}(M)`$ we can determine $`(n2_i1/p_i^2)/4e_0`$, factor that term out, and therefore regard the integral in Equation (109),
$$\sqrt{\frac{e_0^3}{\mathrm{}}}_{}𝑑\beta e^{\mathrm{}e_0\beta ^2/4}\frac{_ip_i\mathrm{sinh}\mathrm{}\beta /2p_i}{\left(\mathrm{sinh}\mathrm{}\beta /2\right)^{n2}}=\frac{1}{\mathrm{}^{3/2}}_{}𝑑\beta e^{\beta ^2/4\mathrm{}}\frac{_i\mathrm{sinh}(\beta /2p_i\sqrt{e_0})}{\left(\mathrm{sinh}(\beta /2\sqrt{e_0})\right)^{n2}}$$
as known. As $`\mathrm{}`$ varies, this is the integral of a fixed function against all possible Gaussians. Knowing all of these integrals we know what the function is. Thus we know all the quantities $`p_i\sqrt{e_0}`$ and the value of $`\sqrt{e_0}`$ (recall that the $`p_i`$’s are all distinct and greater than $`1`$, so no accidental cancellations can occur). This finishes the case of $`n>2`$.
If $`n2`$, then $`M`$ is a sphere and therefore $`\widehat{Z}^{\mathrm{LMO}}(M)=1`$. This leads to $`Z^{\text{rest}}(M)=\frac{\mathrm{sinh}\mathrm{}/2}{\mathrm{}/2}`$, which grows like $`e^{\mathrm{}/2}`$. This growth rate is much slower than the minimal possible value of $`(n2_i1/p_i^2)/4e_0`$ for $`n>2`$, which is attained when $`n=3`$ and $`(p_i)=(2,3,5)`$. Thus the case of $`n2`$ is easily separated from the case of $`n>2`$. $`\mathrm{}`$
###### Exercise 5.7.
Verify directly from Equation (109) that in the $`sl(2)`$ case $`Z^{\text{rest}}(M)=\frac{\mathrm{sinh}\mathrm{}/2}{\mathrm{}/2}`$ for $`n2`$.
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# Contents
## 1 Introduction
There is an obvious interest in having practical algorithms for predicting the future, and there is a correspondingly large literature on the problem of time series extrapolation.<sup>1</sup><sup>1</sup>1The classic papers are by Kolmogoroff (1939, 1941) and Wiener (1949), who essentially solved all the extrapolation problems that could be solved by linear methods. Our understanding of predictability was changed by developments in dynamical systems, which showed that apparently random (chaotic) time series could arise from simple deterministic rules, and this led to vigorous exploration of nonlinear extrapolation algorithms (Abarbanel et al. 1993). For a review comparing different approaches, see the conference proceedings edited by Weigend and Gershenfeld (1994). But prediction is both more and less than extrapolation: we might be able to predict, for example, the chance of rain in the coming week even if we cannot extrapolate the trajectory of temperature fluctuations. In the spirit of its thermodynamic origins, information theory (Shannon 1948) characterizes the potentialities and limitations of all possible prediction algorithms, as well as unifying the analysis of extrapolation with the more general notion of predictability. Specifically, we can define a quantity—the predictive information—that measures how much our observations of the past can tell us about the future. The predictive information characterizes the world we are observing, and we shall see that this characterization is close to our intuition about the complexity of the underlying dynamics.
Prediction is one of the fundamental problems in neural computation. Much of what we admire in expert human performance is predictive in character—the point guard who passes the basketball to a place where his teammate will arrive in a split second, the chess master who knows how moves made now will influence the end game two hours hence, the investor who buys a stock in anticipation that it will grow in the year to come. More generally, we gather sensory information not for its own sake but in the hope that this information will guide our actions (including our verbal actions). But acting takes time, and sense data can guide us only to the extent that those data inform us about the state of the world at the time of our actions, so the only components of the incoming data that have a chance of being useful are those that are predictive. Put bluntly, nonpredictive information is useless to the organism, and it therefore makes sense to isolate the predictive information. It will turn out that most of the information we collect over a long period of time is nonpredictive, so that isolating the predictive information must go a long way toward separating out those features of the sensory world that are relevant for behavior.
One of the most important examples of prediction is the phenomenon of generalization in learning. Learning is formalized as finding a model that explains or describes a set of observations, but again this is useful only because we expect this model will continue to be valid: in the language of learning theory \[see, for example, Vapnik (1998)\] an animal can gain selective advantage not from its performance on the training data but only from its performance at generalization. Generalizing—and not “overfitting” the training data—is precisely the problem of isolating those features of the data that have predictive value (see also Bialek and Tishby, in preparation). Further, we know that the success of generalization hinges on controlling the complexity of the models that we are willing to consider as possibilities. Finally, learning a model to describe a data set can be seen as an encoding of those data, as emphasized by Rissanen (1989), and the quality of this encoding can be measured using the ideas of information theory. Thus the exploration of learning problems should provide us with explicit links among the concepts of entropy, predictability, and complexity.
The notion of complexity arises not only in learning theory, but also in several other contexts. Some physical systems exhibit more complex dynamics than others (turbulent vs. laminar flows in fluids), and some systems evolve toward more complex states than others (spin glasses vs. ferromagnets). The problem of characterizing complexity in physical systems has a substantial literature of its own; for an overview see Bennett (1990). In this context several authors have considered complexity measures based on entropy or mutual information, although as far as we know no clear connections have been drawn among the measures of complexity that arise in learning theory and those that arise in dynamical systems and statistical mechanics.
An essential difficulty in quantifying complexity is to distinguish complexity from randomness. A true random string cannot be compressed and hence requires a long description; it thus is complex in the sense defined by Kolmogorov (1965, Li and Vitányi 1993, Vitányi and Li 2000), yet the physical process that generates this string may have a very simple description. Both in statistical mechanics and in learning theory our intuitive notions of complexity correspond to the statements about complexity of the underlying process, and not directly to the description length or Kolmogorov complexity.
Our central result is that the predictive information provides a general measure of complexity which includes as special cases the relevant concepts from learning theory and from dynamical systems. While work on complexity in learning theory rests specifically on the idea that one is trying to infer a model from data, the predictive information is a property of the data (or, more precisely, of an ensemble of data) itself without reference to a specific class of underlying models. If the data are generated by a process in a known class but with unknown parameters, then we can calculate the predictive information explicitly and show that this information diverges logarithmically with the size of the data set we have observed; the coefficient of this divergence counts the number of parameters in the model, or more precisely the effective dimension of the model class, and this provides a link to known results of Rissanen and others. We also can quantify the complexity of processes that fall outside the conventional finite dimensional models, and we show that these more complex processes are characterized by a power–law rather than a logarithmic divergence of the predictive information.
By analogy with the analysis of critical phenomena in statistical physics, the separation of logarithmic from power–law divergences, together with the measurement of coefficients and exponents for these divergences, allows us to define “universality classes” for the complexity of data streams. The power–law or nonparametric class of processes may be crucial in real world learning tasks, where the effective number of parameters becomes so large that asymptotic results for finitely parameterizable models are inaccessible in practice. There is empirical evidence that simple physical systems can generate dynamics in this complexity class, and there are hints that language also may fall in this class.
Finally, we argue that the divergent components of the predictive information provide a unique measure of complexity that is consistent with certain simple requirements. This argument is in the spirit of Shannon’s original derivation of entropy as the unique measure of available information. We believe that this uniqueness argument provides a conclusive answer to the question of how one should quantify the complexity of a process generating a time series.
With the evident cost of lengthening our discussion, we have tried to give a self–contained presentation that develops our point of view, uses simple examples to connect with known results, and then generalizes and goes beyond these results.<sup>2</sup><sup>2</sup>2Some of the basic ideas presented here, together with some connections to earlier work, can be found in brief preliminary reports (Bialek 1995; Bialek and Tishby 1999). The central results of the present work, however, were at best conjectures in these preliminary accounts. Even in cases where at least the qualitative form of our results is known from previous work, we believe that our point of view elucidates some issues that may have been less the focus of earlier studies. Last but not least, we explore the possibilities for connecting our theoretical discussion with the experimental characterization of learning and complexity in neural systems.
## 2 A curious observation
Before starting the systematic analysis of the problem, we want to motivate our discussion further by presenting results of some simple numerical experiments. Since most of the paper draws examples from learning, here we consider examples from equilibrium statistical mechanics. Suppose that we have a one dimensional chain of Ising spins with the Hamiltonian given by
$$H=\underset{\mathrm{i},\mathrm{j}}{}J_{\mathrm{ij}}\sigma _\mathrm{i}\sigma _\mathrm{j},$$
(1)
where the matrix of interactions $`J_{\mathrm{ij}}`$ is not restricted to nearest neighbors—long range interactions are also allowed. One may identify spins pointing upwards with $`1`$ and downwards with $`0`$, and then a spin chain is equivalent to some sequence of binary digits. This sequence consists of (overlapping) words of $`N`$ digits each, $`W_\mathrm{k}`$, $`k=0,1\mathrm{}2^N1`$. There are $`2^N`$ such words total, and they appear with very different frequencies $`n(W_k)`$ in the spin chain \[see Fig. (1) for details\].
If the number of spins is large, then counting these frequencies provides a good empirical estimate of $`P_N(W_k)`$, the probability distribution of different words of length $`N`$. Then one can calculate the entropy $`S(N)`$ of this probability distribution by the usual formula
$$S(N)=\underset{k=0}{\overset{2^N1}{}}P_N(W_k)\mathrm{log}_2P_N(W_k)(\mathrm{bits}).$$
(2)
Note that this is not the entropy of a finite chain with length $`N`$; instead it is the entropy of words or strings with length $`N`$ drawn from a much longer chain. Nonetheless, since entropy is an extensive property, $`S(N)`$ is proportional asymptotically to $`N`$ for any spin chain, that is $`S(N)𝒮_0N`$. The usual goal in statistical mechanics is to understand this “thermodynamic limit” $`N\mathrm{}`$, and hence to calculate the entropy density $`𝒮_0`$. Different sets of interactions $`J_{\mathrm{ij}}`$ result in different values of $`𝒮_0`$, but the qualitative result $`S(N)N`$ is true for all reasonable $`\{J_{\mathrm{ij}}\}`$.
We investigated three different spin chains of one billion spins each. As usual in statistical mechanics the probability of any configuration of spins $`\{\sigma _\mathrm{i}\}`$ is given by the Boltzmann distribution,
$$P[\{\sigma _\mathrm{i}\}]\mathrm{exp}(H/k_BT),$$
(3)
where to normalize the scale of the $`J_{\mathrm{ij}}`$ we set $`k_BT=1`$. For the first chain, only $`J_{\mathrm{i},\mathrm{i}+1}=1`$ was nonzero, and its value was the same for all $`\mathrm{i}`$. The second chain was also generated using the nearest neighbor interactions, but the value of the coupling was reset every 400,000 spins by taking a random number from a Gaussian distribution with zero mean and unit variance. In the third case, we again reset the interactions at the same frequency, but now interactions were long–ranged; the variances of coupling constants decreased with the distance between the spins as $`J_{\mathrm{ij}}^2=1/(\mathrm{i}\mathrm{j})^2`$. We plot $`S(N)`$ for all these cases in Fig. (2),
and, of course, the asymptotically linear behavior is evident—the extensive entropy shows no qualitative distinction among the three cases we consider.
However, the situation changes drastically if we remove the asymptotic linear contribution and focus on the corrections to extensive behavior. Specifically, we write $`S(N)=𝒮_0N+S_1(N)`$, and plot only the sublinear component $`S_1(N)`$ of the entropy.
As we see in Fig. (3), the three chains then exhibit qualitatively different features: for the first one, $`S_1`$ is constant; for the second one, it is logarithmic; and, for the third one, it clearly shows a power–law behavior.
What is the significance of these observations? Of course, the differences in the behavior of $`S_1(N)`$ must be related to the ways we chose $`J`$’s for the simulations. In the first case, $`J`$ is fixed, and if we see $`N`$ spins and try to predict the state of the $`N+1^{\mathrm{st}}`$ spin, all that really matters is the state of the spin $`\sigma _\mathrm{N}`$—there is nothing to “learn” from observations on longer segments of the chain. For the second chain, $`J`$ changes, and the statistics of the spin–words are different in different parts of the sequence. By looking at these statistics, one can “learn” the coupling at the current position; this estimate improves the more spins (longer words) we observe. Finally, in the third case there are many coupling constants that can be learned—as $`N`$ increases one becomes sensitive to weaker correlations caused by interactions over larger and larger distances. So, intuitively, the qualitatively different behaviors of $`S_1(N)`$ in the three plotted cases are correlated with differences in the problem of learning the underlying dynamics of the spin chains from observations on samples of the spins themselves. Much of this paper can be seen as expanding on and quantifying this intuitive observation.
## 3 Fundamentals
The problem of prediction comes in various forms, as noted above. Information theory allows us to treat the different notions of prediction on the same footing. The first step is to recognize that all predictions are probabilistic—even if we can predict the temperature at noon tomorrow, we should provide error bars or confidence limits on our prediction. The next step is to remember that, even before we look at the data, we know that certain futures are more likely than others, and we can summarize this knowledge by a prior probability distribution for the future. Our observations on the past lead us to a new, more tightly concentrated distribution, the distribution of futures conditional on the past data. Different kinds of predictions are different slices through or averages over this conditional distribution, but information theory quantifies the “concentration” of the distribution without making any commitment as to which averages will be most interesting.
Imagine that we observe a stream of data $`x(t)`$ over a time interval $`T<t<0`$; let all of these past data be denoted by the shorthand $`x_{\mathrm{past}}`$. We are interested in saying something about the future, so we want to know about the data $`x(t)`$ that will be observed in the time interval $`0<t<T^{}`$; let these future data be called $`x_{\mathrm{future}}`$. In the absence of any other knowledge, futures are drawn from the probability distribution $`P(x_{\mathrm{future}})`$, while observations of particular past data $`x_{\mathrm{past}}`$ tell us that futures will be drawn from the conditional distribution $`P(x_{\mathrm{future}}|x_{\mathrm{past}})`$. The greater concentration of the conditional distribution can be quantified by the fact that it has smaller entropy than the prior distribution, and this reduction in entropy is Shannon’s definition of the information that the past provides about the future. We can write the average of this predictive information as
$`_{\mathrm{pred}}(T,T^{})`$ $`=`$ $`\mathrm{log}_2\left[{\displaystyle \frac{P(x_{\mathrm{future}}|x_{\mathrm{past}})}{P(x_{\mathrm{future}})}}\right]`$ (5)
$`=`$ $`\mathrm{log}_2P(x_{\mathrm{future}})\mathrm{log}_2P(x_{\mathrm{past}})`$
$`\left[\mathrm{log}_2P(x_{\mathrm{future}},x_{\mathrm{past}})\right],`$
where $`\mathrm{}`$ denotes an average over the joint distribution of the past and the future, $`P(x_{\mathrm{future}},x_{\mathrm{past}})`$.
Each of the terms in Eq. (5) is an entropy. Since we are interested in predictability or generalization, which are associated with some features of the signal persisting forever, we may assume stationarity or invariance under time translations. Then the entropy of the past data depends only on the duration of our observations, so we can write $`\mathrm{log}_2P(x_{\mathrm{past}})=S(T)`$, and by the same argument $`\mathrm{log}_2P(x_{\mathrm{future}})=S(T^{})`$. Finally, the entropy of the past and the future taken together is the entropy of observations on a window of duration $`T+T^{}`$, so that $`\mathrm{log}_2P(x_{\mathrm{future}},x_{\mathrm{past}})=S(T+T^{})`$. Putting these equations together, we obtain
$$_{\mathrm{pred}}(T,T^{})=S(T)+S(T^{})S(T+T^{}).$$
(6)
It is important to recall that mutual information is a symmetric quantity. Thus we can view $`_{\mathrm{pred}}(T,T^{})`$ either as the information that a data segment of duration $`T`$ provides about the future of length $`T^{}`$, or as the information that a data segment of duration $`T^{}`$ provides about the immediate past of duration $`T`$. This is a direct application of the definitions of information, but seems counterintuitive: shouldn’t it be more difficult to predict than to postdict? One can perhaps recover the correct intuition by thinking about a large ensemble of question/answer pairs. Prediction corresponds to generating the answer to a given question, while postdiction corresponds to generating the question that goes with a given answer. We know that guessing questions given answers is also a hard problem,<sup>3</sup><sup>3</sup>3This is the basis of the popular American television game show Jeopardy!, widely viewed as the most ‘intellectual’ of its genre. and can be just as challenging as the more conventional problem of answering the questions themselves. Our focus here is on prediction because we want to make connections with the phenomenon of generalization in learning, but it is clear that generating a consistent interpretation of observed data may involve elements of both prediction and postdiction \[see, for example, Eagleman and Sejnowski (2000)\]; it is attractive that the information theoretic formulation treats these problems symmetrically.
In the same way that the entropy of a gas at fixed density is proportional to the volume, the entropy of a time series (asymptotically) is proportional to its duration, so that $`lim_T\mathrm{}S(T)/T=𝒮_0`$; entropy is an extensive quantity. But from Eq. (6) any extensive component of the entropy cancels in the computation of the predictive information: predictability is a deviation from extensivity. If we write $`S(T)=𝒮_0T+S_1(T)`$, then Eq. (6) tells us that the predictive information is related only to the nonextensive term $`S_1(T)`$. Note that if we are observing a deterministic system, then $`𝒮_0=0`$, but this is independent of questions about the structure of the subextensive term $`S_1(T)`$. It is attractive that information theory gives us a unified discussion of prediction in deterministic and probabilistic cases.
We know two general facts about the behavior of $`S_1(T)`$. First, the corrections to extensive behavior are positive, $`S_1(T)0`$. Second, the statement that entropy is extensive is the statement that the limit
$`\underset{T\mathrm{}}{lim}{\displaystyle \frac{S(T)}{T}}=𝒮_0`$ (7)
exists, and for this to be true we must also have
$`\underset{T\mathrm{}}{lim}{\displaystyle \frac{S_1(T)}{T}}=0.`$ (8)
Thus the nonextensive terms in the entropy must be subextensive, that is they must grow with $`T`$ less rapidly than a linear function. Taken together, these facts guarantee that the predictive information is positive and subextensive. Further, if we let the future extend forward for a very long time, $`T^{}\mathrm{}`$, then we can measure the information that our sample provides about the entire future,
$$I_{\mathrm{pred}}(T)=\underset{T^{}\mathrm{}}{lim}_{\mathrm{pred}}(T,T^{})=S_1(T).$$
(9)
Similarly, instead of increasing the duration of the future to infinity we could have considered the mutual information between a sample of length $`T`$ and all of the infinite past. Then the postdictive information also is equal to $`S_1(T)`$, and the symmetry between prediction and postdiction is even more profound: not only is there symmetry between questions and answers, but observations on a given period of time provide the same amount of information about the historical path that led to our observations as about the future that will unfold from them. In some cases this statement becomes even stronger. For example, if the subextensive entropy of a long discontinuous observation of a total length $`T`$ with a gap of a duration $`\delta TT`$ is equal to $`S_1(T)+O(\frac{\delta T}{T})`$, then the subextensive entropy of the present is not only its information about the past or the future, but also the information about the past and the future.
If we have been observing a time series for a (long) time $`T`$, then the total amount of data we have collected in is measured by the entropy $`S(T)`$, and at large $`T`$ this is given approximately by $`𝒮_0T`$. But the predictive information that we have gathered cannot grow linearly with time, even if we are making predictions about a future which stretches out to infinity. As a result, of the total information we have taken in by observing $`x_{\mathrm{past}}`$, only a vanishing fraction is of relevance to the prediction:
$$\underset{T\mathrm{}}{lim}\frac{\mathrm{Predictive}\mathrm{Information}}{\mathrm{Total}\mathrm{Information}}=\frac{I_{\mathrm{pred}}(T)}{S(T)}0.$$
(10)
In this precise sense, most of what we observe is irrelevant to the problem of predicting the future.<sup>4</sup><sup>4</sup>4We can think of Eq. (10) as a law of diminishing returns: although we collect data in proportion to our observation time $`T`$, a smaller and smaller fraction of this information is useful in the problem of prediction. These diminishing returns are not due to a limited lifetime, since we calculate the predictive information assuming that we have a future extending forward to infinity. A senior colleague points out that this is an argument for changing fields before becoming too expert.
Consider the case where time is measured in discrete steps, so that we have seen $`N`$ time points $`x_1,x_2,\mathrm{},x_N`$. How much have we learned about the underlying pattern in these data? The more we know, the more effectively we can predict the next data point $`x_{N+1}`$ and hence the fewer bits we will need to describe the deviation of this data point from our prediction: our accumulated knowledge about the time series is measured by the degree to which we can compress the description of new observations. On average, the length of the code word required to describe the point $`x_{N+1}`$, given that we have seen the previous $`N`$ points, is given by
$$\mathrm{}(N)=\mathrm{log}_2P(x_{N+1}|x_1,x_2,\mathrm{},x_N)\mathrm{bits},$$
(11)
where the expectation value is taken over the joint distribution of all the $`N+1`$ points, $`P(x_1,x_2,\mathrm{},x_N,x_{N+1})`$. It is easy to see that
$`\mathrm{}(N)=S(N+1)S(N){\displaystyle \frac{S(N)}{N}}.`$ (12)
As we observe for longer times, we learn more and this word length decreases. It is natural to define a learning curve that measures this improvement. Usually we define learning curves by measuring the frequency or costs of errors; here the cost is that our encoding of the point $`x_{N+1}`$ is longer than it could be if we had perfect knowledge. This ideal encoding has a length which we can find by imagining that we observe the time series for an infinitely long time, $`\mathrm{}_{\mathrm{ideal}}=lim_N\mathrm{}\mathrm{}(N)`$, but this is just another way of defining the extensive component of the entropy $`𝒮_0`$. Thus we can define a learning curve
$`\mathrm{\Lambda }(N)`$ $``$ $`\mathrm{}(N)\mathrm{}_{\mathrm{ideal}}`$
$`=`$ $`S(N+1)S(N)𝒮_0`$
$`=`$ $`S_1(N+1)S_1(N)`$
$``$ $`{\displaystyle \frac{S_1(N)}{N}}={\displaystyle \frac{I_{\mathrm{pred}}(N)}{N}},`$ (14)
and we see once again that the extensive component of the entropy cancels.
It is well known that the problems of prediction and compression are related, and what we have done here is to illustrate one aspect of this connection. Specifically, if we ask how much one segment of a time series can tell us about the future, the answer is contained in the subextensive behavior of the entropy. If we ask how much we are learning about the structure of the time series, then the natural and universally defined learning curve is related again to the subextensive entropy: the learning curve is the derivative of the predictive information.
This universal learning curve is connected to the more conventional learning curves in specific contexts. As an example (cf. Section 4.1), consider fitting a set of data points $`\{x_\mathrm{n},y_\mathrm{n}\}`$ with some class of functions $`y=f(x;𝜶),`$ where the $`𝜶`$ are unknown parameters that need to be learned; we also allow for some Gaussian noise in our observation of the $`y_\mathrm{n}`$. Here the natural learning curve is the evolution of $`\chi ^2`$ for generalization as a function of the number of examples. Within the approximations discussed below, it is straightforward to show that as $`N`$ becomes large,
$`\chi ^2(N)={\displaystyle \frac{1}{\sigma ^2}}\left[yf(x;𝜶)\right]^2(2\mathrm{ln}2)\mathrm{\Lambda }(N)+1,`$ (15)
where $`\sigma ^2`$ is the variance of the noise. Thus a more conventional measure of performance at learning a function is equal to the universal learning curve defined purely by information theoretic criteria. In other words, if a learning curve is measured in the right units, then its integral represents the amount of the useful information accumulated. Then the subextensivity of $`S_1`$ guarantees that the learning curve decreases to zero as $`N\mathrm{}`$.
Different quantities related to the subextensive entropy have been discussed in several contexts. For example, the code length $`\mathrm{}(N)`$ has been defined as a learning curve in the specific case of neural networks (Opper and Haussler 1995) and has been termed the “thermodynamic dive” (Crutchfield and Shalizi 1998) and “$`N^{\mathrm{th}}`$ order block entropy” (Grassberger 1986). The universal learning curve $`\mathrm{\Lambda }(N)`$ has been studied as the “expected instantaneous information gain” by Haussler et al. (1994). Mutual information between all of the past and all of the future (both semi–infinite) is known also as the “excess entropy,” “effective measure complexity,” “stored information,” and so on \[see Shalizi and Crutchfield (1999) and references therein, as well as the discussion below\]. If the data allow a description by a model with finite (and in some cases also infinite) number of parameters, then mutual information between the data and the parameters is of interest. This is easily shown to be equal to the predictive information about all of the future, and it is also the “cumulative information gain” (Haussler et al. 1994) or the “cumulative relative entropy risk” (Haussler and Opper 1997). Investigation of this problem can be traced back to Renyi (1964) and Ibragimov and Hasminskii (1972), and some particular topics are still being discussed (Haussler and Opper 1995, Opper and Haussler 1995, Herschkowitz and Nadal 1999). In certain limits, decoding a signal from a population of $`N`$ neurons can be thought of as ‘learning’ a parameter from $`N`$ observations, with a parallel notion of information transmission (Brunel and Nadal 1998, Kang and Sompolinsky 2001). In addition, as noted already, the subextensive component of the description length (Rissanen 1978, 1989, 1996, Clarke and Barron 1990) averaged over a class of allowed models also is similar to the predictive information. What is important is that the predictive information or subextensive entropy is related to all these quantities, and that it can be defined for any process without a reference to a class of models. It is this universality that we find appealing, and this universality is strongest if we focus on the limit of long observation times. Qualitatively, in this regime ($`T\mathrm{}`$) we expect the predictive information to behave in one of three different ways, as illustrated by the Ising models above: it may either stay finite, or grow to infinity together with $`T`$; in the latter case the rate of growth may be slow (logarithmic) or fast (sublinear power) \[see Barron and Cover (1991) for a similar classification in the framework of the Minimal Description Length (MDL) analysis\].
The first possibility, $`lim_T\mathrm{}I_{\mathrm{pred}}(T)=`$ constant, means that no matter how long we observe we gain only a finite amount of information about the future. This situation prevails, for example, when the dynamics are too regular: for a purely periodic system, complete prediction is possible once we know the phase, and if we sample the data at discrete times this is a finite amount of information; longer period orbits intuitively are more complex and also have larger $`I_{\mathrm{pred}}`$, but this doesn’t change the limiting behavior $`lim_T\mathrm{}I_{\mathrm{pred}}(T)=`$ constant.
Alternatively, the predictive information can be small when the dynamics are irregular but the best predictions are controlled only by the immediate past, so that the correlation times of the observable data are finite \[see, for example, Crutchfield and Feldman (1997) and the fixed $`J`$ case in Fig. (3)\]. Imagine, for example, that we observe $`x(t)`$ at a series of discrete times $`\{t_\mathrm{n}\}`$, and that at each time point we find the value $`x_\mathrm{n}`$. Then we always can write the joint distribution of the $`N`$ data points as a product,
$`P(x_1,x_2,\mathrm{},x_N)`$ $`=`$ $`P(x_1)P(x_2|x_1)P(x_3|x_2,x_1)\mathrm{}.`$ (16)
For Markov processes, what we observe at $`t_\mathrm{n}`$ depends only on events at the previous time step $`t_{\mathrm{n}1}`$, so that
$`P(x_\mathrm{n}|\{x_{1\mathrm{i}\mathrm{n}1}\})`$ $`=`$ $`P(x_\mathrm{n}|x_{\mathrm{n}1}),`$ (17)
and hence the predictive information reduces to
$`I_{\mathrm{pred}}=\mathrm{log}_2\left[{\displaystyle \frac{P(x_\mathrm{n}|x_{\mathrm{n}1})}{P(x_\mathrm{n})}}\right].`$ (18)
The maximum possible predictive information in this case is the entropy of the distribution of states at one time step, which in turn is bounded by the logarithm of the number of accessible states. To approach this bound the system must maintain memory for a long time, since the predictive information is reduced by the entropy of the transition probabilities. Thus systems with more states and longer memories have larger values of $`I_{\mathrm{pred}}`$.
More interesting are those cases in which $`I_{\mathrm{pred}}(T)`$ diverges at large $`T`$. In physical systems we know that there are critical points where correlation times become infinite, so that optimal predictions will be influenced by events in the arbitrarily distant past. Under these conditions the predictive information can grow without bound as $`T`$ becomes large; for many systems the divergence is logarithmic, $`I_{\mathrm{pred}}(T\mathrm{})\mathrm{ln}T`$, as for the variable $`J_{\mathrm{ij}}`$, short range Ising model of Figs. (2, 3). Long range correlation also are important in a time series where we can learn some underlying rules. It will turn out that when the set of possible rules can be described by a finite number of parameters, the predictive information again diverges logarithmically, and the coefficient of this divergence counts the number of parameters. Finally, a faster growth is also possible, so that $`I_{\mathrm{pred}}(T\mathrm{})T^\alpha `$, as for the variable $`J_{\mathrm{ij}}`$ long range Ising model, and we shall see that this behavior emerges from, for example, nonparametric learning problems.
## 4 Learning and predictability
Learning is of interest precisely in those situations where correlations or associations persist over long periods of time. In the usual theoretical models, there is some rule underlying the observable data, and this rule is valid forever; examples seen at one time inform us about the rule, and this information can be used to make predictions or generalizations. The predictive information quantifies the average generalization power of examples, and we shall see that there is a direct connection between the predictive information and the complexity of the possible underlying rules.
### 4.1 A test case
Let us begin with a simple example already mentioned above. We observe two streams of data $`x`$ and $`y`$, or equivalently a stream of pairs $`(x_1,y_1)`$, $`(x_2,y_2)`$, $`\mathrm{}`$ , $`(x_\mathrm{N},y_\mathrm{N})`$. Assume that we know in advance that the $`x`$’s are drawn independently and at random from a distribution $`P(x)`$, while the $`y`$’s are noisy versions of some function acting on $`x`$,
$`y_\mathrm{n}=f(x_\mathrm{n};𝜶)+\eta _\mathrm{n},`$ (19)
where $`f(x;𝜶)`$ is a class of functions parameterized by $`𝜶`$, and $`\eta _\mathrm{n}`$ is noise, which for simplicity we will assume is Gaussian with known standard deviation $`\sigma `$. We can even start with a very simple case, where the function class is just a linear combination of basis functions, so that
$`f(x;𝜶)={\displaystyle \underset{\mu =1}{\overset{K}{}}}\alpha _\mu \varphi _\mu (x).`$ (20)
The usual problem is to estimate, from $`N`$ pairs $`\{x_\mathrm{i},y_\mathrm{i}\}`$, the values of the parameters $`𝜶`$; in favorable cases such as this we might even be able to find an effective regression formula. We are interested in evaluating the predictive information, which means that we need to know the entropy $`S(N)`$. We go through the calculation in some detail because it provides a model for the more general case.
To evaluate the entropy $`S(N)`$ we first construct the probability distribution $`P(x_1,y_1,x_2,y_2,\mathrm{},x_\mathrm{N},y_\mathrm{N})`$. The same set of rules applies to the whole data stream, which here means that the same parameters $`𝜶`$ apply for all pairs $`\{x_\mathrm{i},y_\mathrm{i}\}`$, but these parameters are chosen at random from a distribution $`𝒫(𝜶)`$ at the start of the stream. Thus we write
$`P(x_1,y_1,x_2,y_2,\mathrm{},x_\mathrm{N},y_\mathrm{N})`$
$`={\displaystyle d^K\alpha P(x_1,y_1,x_2,y_2,\mathrm{},x_\mathrm{N},y_\mathrm{N}|𝜶)𝒫(𝜶)},`$ (21)
and now we need to construct the conditional distributions for fixed $`𝜶`$. By hypothesis each $`x`$ is chosen independently, and once we fix $`𝜶`$ each $`y_\mathrm{i}`$ is correlated only with the corresponding $`x_\mathrm{i}`$, so that we have
$`P(x_1,y_1,x_2,y_2,\mathrm{},x_\mathrm{N},y_\mathrm{N}|𝜶)={\displaystyle \underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}}\left[P(x_\mathrm{i})P(y_\mathrm{i}|x_\mathrm{i};𝜶)\right].`$ (22)
Further, with the simple assumptions above about the class of functions and Gaussian noise, the conditional distribution of $`y_\mathrm{i}`$ has the form
$`P(y_\mathrm{i}|x_\mathrm{i};𝜶)={\displaystyle \frac{1}{\sqrt{2\pi \sigma ^2}}}\mathrm{exp}\left[{\displaystyle \frac{1}{2\sigma ^2}}\left(y_\mathrm{i}{\displaystyle \underset{\mu =1}{\overset{K}{}}}\alpha _\mu \varphi _\mu (x_\mathrm{i})\right)^2\right].`$ (23)
Putting all these factors together,
$`P(x_1,y_1,x_2,y_2,\mathrm{},x_\mathrm{N},y_\mathrm{N})`$
$`=\left[{\displaystyle \underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}}P(x_\mathrm{i})\right]\left({\displaystyle \frac{1}{\sqrt{2\pi \sigma ^2}}}\right)^N{\displaystyle d^K\alpha 𝒫(𝜶)\mathrm{exp}\left[\frac{1}{2\sigma ^2}\underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}y_\mathrm{i}^2\right]}`$
$`\times \mathrm{exp}\left[{\displaystyle \frac{N}{2}}{\displaystyle \underset{\mu ,\nu =1}{\overset{K}{}}}A_{\mu \nu }(\{x_\mathrm{i}\})\alpha _\mu \alpha _\nu +N{\displaystyle \underset{\mu =1}{\overset{K}{}}}B_\mu (\{x_\mathrm{i},y_\mathrm{i}\})\alpha _\mu \right],`$ (24)
where
$`A_{\mu \nu }(\{x_\mathrm{i}\})`$ $`=`$ $`{\displaystyle \frac{1}{\sigma ^2N}}{\displaystyle \underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}}\varphi _\mu (x_\mathrm{i})\varphi _\nu (x_\mathrm{i}),\mathrm{and}`$ (25)
$`B_\mu (\{x_\mathrm{i},y_\mathrm{i}\})`$ $`=`$ $`{\displaystyle \frac{1}{\sigma ^2N}}{\displaystyle \underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}}y_\mathrm{i}\varphi _\mu (x_\mathrm{i}).`$ (26)
Our placement of the factors of $`N`$ means that both $`A_{\mu \nu }`$ and $`B_\mu `$ are of order unity as $`N\mathrm{}`$. These quantities are empirical averages over the samples $`\{x_\mathrm{i},y_\mathrm{i}\}`$, and if the $`\varphi _\mu `$ are well behaved we expect that these empirical means converge to expectation values for most realizations of the series $`\{x_\mathrm{i}\}`$:
$`\underset{N\mathrm{}}{lim}A_{\mu \nu }(\{x_\mathrm{i}\})`$ $`=`$ $`A_{\mu \nu }^{\mathrm{}}={\displaystyle \frac{1}{\sigma ^2}}{\displaystyle 𝑑xP(x)\varphi _\mu (x)\varphi _\nu (x)},`$ (27)
$`\underset{N\mathrm{}}{lim}B_\mu (\{x_\mathrm{i},y_\mathrm{i}\})`$ $`=`$ $`B_\mu ^{\mathrm{}}={\displaystyle \underset{\nu =1}{\overset{K}{}}}A_{\mu \nu }^{\mathrm{}}\overline{\alpha }_\nu ,`$ (28)
where $`\overline{𝜶}`$ are the parameters that actually gave rise to the data stream $`\{x_\mathrm{i},y_\mathrm{i}\}`$. In fact we can make the same argument about the terms in $`y_\mathrm{i}^2`$,
$`\underset{N\mathrm{}}{lim}{\displaystyle \underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}}y_\mathrm{i}^2=N\sigma ^2\left[{\displaystyle \underset{\mu ,\nu =1}{\overset{K}{}}}\overline{\alpha }_\mu A_{\mu \nu }^{\mathrm{}}\overline{\alpha }_\nu +1\right].`$ (29)
Conditions for this convergence of empirical means to expectation values are at the heart of learning theory. Our approach here is first to assume that this convergence works, then to examine the consequences for the predictive information, and finally to address the conditions for and implications of this convergence breaking down.
Putting the different factors together, we obtain
$`P(x_1,y_1,x_2,y_2,\mathrm{},x_\mathrm{N},y_\mathrm{N})`$
$`\stackrel{~}{}\left[{\displaystyle \underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}}P(x_\mathrm{i})\right]\left({\displaystyle \frac{1}{\sqrt{2\pi \sigma ^2}}}\right)^N{\displaystyle d^K\alpha 𝒫(𝜶)\mathrm{exp}\left[NE_N(𝜶;\{x_\mathrm{i},y_\mathrm{i}\})\right]},`$
where the effective “energy” per sample is given by
$`E_N(𝜶;\{x_\mathrm{i},y_\mathrm{i}\})={\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu ,\nu =1}{\overset{K}{}}}(\alpha _\mu \overline{\alpha }_\mu )A_{\mu \nu }^{\mathrm{}}(\alpha _\nu \overline{\alpha }_\nu ).`$ (31)
Here we use the symbol $`\stackrel{~}{}`$ to indicate that we not only take the limit of large $`N`$, but also neglect the fluctuations. Note that in this approximation the dependence on the sample points themselves is hidden in the definition of $`\overline{𝜶}`$ as being the parameters that generated the samples.
The integral that we need to do in Eq. (4.1) involves an exponential with a large factor $`N`$ in the exponent; the energy $`E_N`$ is of order unity as $`N\mathrm{}`$. This suggests that we evaluate the integral by a saddle point or steepest descent approximation \[similar analyses were performed by Clarke and Barron (1990), by MacKay (1992), and by Balasubramanian (1997)\]:
$`{\displaystyle d^K\alpha 𝒫(𝜶)\mathrm{exp}\left[NE_N(𝜶;\{x_\mathrm{i},y_\mathrm{i}\})\right]}𝒫(𝜶_{\mathrm{cl}})`$
$`\times \mathrm{exp}\left[NE_N(𝜶_{\mathrm{cl}};\{x_\mathrm{i},y_\mathrm{i}\}){\displaystyle \frac{K}{2}}\mathrm{ln}{\displaystyle \frac{N}{2\pi }}{\displaystyle \frac{1}{2}}\mathrm{ln}det_N+\mathrm{}\right],`$ (32)
where $`𝜶_{\mathrm{cl}}`$ is the “classical” value of $`𝜶`$ determined by the extremal conditions
$`{\displaystyle \frac{E_N(𝜶;\{x_\mathrm{i},y_\mathrm{i}\})}{\alpha _\mu }}|_{𝜶=𝜶_{\mathrm{cl}}}=0,`$ (33)
the matrix $`_N`$ consists of the second derivatives of $`E_N`$,
$`_N={\displaystyle \frac{^2E_N(𝜶;\{x_\mathrm{i},y_\mathrm{i}\})}{\alpha _\mu \alpha _\nu }}|_{𝜶=𝜶_{\mathrm{cl}}},`$ (34)
and $`\mathrm{}`$ denotes terms that vanish as $`N\mathrm{}`$. If we formulate the problem of estimating the parameters $`𝜶`$ from the samples $`\{x_\mathrm{i},y_\mathrm{i}\}`$, then as $`N\mathrm{}`$ the matrix $`N_N`$ is the Fisher information matrix (Cover and Thomas 1991); the eigenvectors of this matrix give the principal axes for the error ellipsoid in parameter space, and the (inverse) eigenvalues give the variances of parameter estimates along each of these directions. The classical $`𝜶_{\mathrm{cl}}`$ differs from $`\overline{𝜶}`$ only in terms of order $`1/N`$; we neglect this difference and further simplify the calculation of leading terms as $`N`$ becomes large. After a little more algebra, then, we find the probability distribution we have been looking for:
$`P(x_1,y_1,x_2,y_2,\mathrm{},x_\mathrm{N},y_\mathrm{N})`$
$`\stackrel{~}{}\left[{\displaystyle \underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}}P(x_\mathrm{i})\right]{\displaystyle \frac{1}{Z_A}}𝒫(\overline{𝜶})\mathrm{exp}\left[{\displaystyle \frac{N}{2}}\mathrm{ln}(2\pi \mathrm{e}\sigma ^2){\displaystyle \frac{K}{2}}\mathrm{ln}N+\mathrm{}\right],`$ (35)
where the normalization constant
$`Z_A=\sqrt{(2\pi )^KdetA^{\mathrm{}}}.`$ (36)
Again we note that the sample points $`\{x_\mathrm{i},y_\mathrm{i}\}`$ are hidden in the value of $`\overline{𝜶}`$ that gave rise to these points.<sup>5</sup><sup>5</sup>5We emphasize once more that there are two approximations leading to Eq. (35). First, we have replaced empirical means by expectation values, neglecting fluctuations associated with the particular set of sample points $`\{x_\mathrm{i},y_\mathrm{i}\}`$. Second, we have evaluated the average over parameters in a saddle point approximation. At least under some conditions, both of these approximations become increasingly accurate as $`N\mathrm{}`$, so that this approach should yield the asymptotic behavior of the distribution and hence the subextensive entropy at large $`N`$. Although we give a more detailed analysis below, it is worth noting here how things can go wrong. The two approximations are independent, and we could imagine that fluctuations are important but saddle point integration still works, for example. Controlling the fluctuations turns out to be exactly the question of whether our finite parameterization captures the true dimensionality of the class of models, as discussed in the classic work of Vapnik, Chervonenkis, and others \[see Vapnik (1998) for a review\]. The saddle point approximation can break down because the saddle point becomes unstable or because multiple saddle points become important. It will turn out that instability is exponentially improbable as $`N\mathrm{}`$, while multiple saddle points are a real problem in certain classes of models, again when counting parameters doesn’t really measure the complexity of the model class.
To evaluate the entropy $`S(N)`$ we need to compute the expectation value of the (negative) logarithm of the probability distribution in Eq. (35); there are three terms. One is constant, so averaging is trivial. The second term depends only on the $`x_\mathrm{i}`$, and because these are chosen independently from the distribution $`P(x)`$ the average again is easy to evaluate. The third term involves $`\overline{𝜶}`$, and we need to average this over the joint distribution $`P(x_1,y_1,x_2,y_2,\mathrm{},x_\mathrm{N},y_\mathrm{N})`$. As above, we can evaluate this average in steps: first we choose a value of the parameters $`\overline{𝜶}`$, then we average over the samples given these parameters, and finally we average over parameters. But because $`\overline{𝜶}`$ is defined as the parameters that generate the samples, this stepwise procedure simplifies enormously. The end result is that
$$S(N)=N\left[S_x+\frac{1}{2}\mathrm{log}_2(2\pi \mathrm{e}\sigma ^2)\right]+\frac{K}{2}\mathrm{log}_2N+S_𝜶+\mathrm{log}_2Z_A_𝜶+\mathrm{},$$
(37)
where $`\mathrm{}_𝜶`$ means averaging over parameters, $`S_x`$ is the entropy of the distribution of $`x`$,
$`S_x={\displaystyle 𝑑xP(x)\mathrm{log}_2P(x)},`$ (38)
and similarly for the entropy of the distribution of parameters,
$`S_𝜶={\displaystyle d^K\alpha 𝒫(𝜶)\mathrm{log}_2𝒫(𝜶)}.`$ (39)
The different terms in the entropy Eq. (37) have a straightforward interpretation. First we see that the extensive term in the entropy,
$`𝒮_0=S_x+{\displaystyle \frac{1}{2}}\mathrm{log}_2(2\pi \mathrm{e}\sigma ^2),`$ (40)
reflects contributions from the random choice of $`x`$ and from the Gaussian noise in $`y`$; these extensive terms are independent of the variations in parameters $`𝜶`$, and these would be the only terms if the parameters were not varying (that is, if there were nothing to learn). There also is a term which reflects the entropy of variations in the parameters themselves, $`S_𝜶`$. This entropy is not invariant with respect to coordinate transformations in the parameter space, but the term $`\mathrm{log}_2Z_A_𝜶`$ compensates for this noninvariance. Finally, and most interestingly for our purposes, the subextensive piece of the entropy is dominated by a logarithmic divergence,
$$S_1(N)\frac{K}{2}\mathrm{log}_2N(\mathrm{bits}).$$
(41)
The coefficient of this divergence counts the number of parameters independent of the coordinate system that we choose in the parameter space. Furthermore, this result does not depend on the set of basis functions $`\{\varphi _\mu (x)\}`$. This is a hint that the result in Eq. (41) is more universal than our simple example.
### 4.2 Learning a parameterized distribution
The problem discussed above is an example of supervised learning: we are given examples of how the points $`x_\mathrm{n}`$ map into $`y_\mathrm{n}`$, and from these examples we are to induce the association or functional relation between $`x`$ and $`y`$. An alternative view is that pair of points $`(x,y)`$ should be viewed as a vector $`\stackrel{}{x}`$, and what we are learning is the distribution of this vector. The problem of learning a distribution usually is called unsupervised learning, but in this case supervised learning formally is a special case of unsupervised learning; if we admit that all the functional relations or associations that we are trying to learn have an element of noise or stochasticity, then this connection between supervised and unsupervised problems is quite general.
Suppose a series of random vector variables $`\{\stackrel{}{x}_\mathrm{i}\}`$ are drawn independently from the same probability distribution $`Q(\stackrel{}{x}|𝜶)`$, and this distribution depends on a (potentially infinite dimensional) vector of parameters $`𝜶`$. As above, the parameters are unknown, and before the series starts they are chosen randomly from a distribution $`𝒫(𝜶)`$. With no constraints on the densities $`𝒫(𝜶)`$ or $`Q(\stackrel{}{x}|𝜶)`$ it is impossible to derive any regression formulas for parameter estimation, but one can still say something about the entropy of the data series and thus the predictive information. For a finite dimensional vector of parameters $`𝜶`$ the literature on bounding similar quantities is especially rich (Haussler et al. 1994, Wong and Shen 1995, Haussler and Opper 1995, Haussler and Opper 1997 and references therein), and related asymptotic calculations have been done (Clarke and Barron 1990, MacKay 1992, Balasubramanian 1997).
We begin with the definition of entropy
$`S(N)`$ $``$ $`S[\left\{\stackrel{}{x}_\mathrm{i}\right\}]=`$ (42)
$`{\displaystyle 𝑑\stackrel{}{x}_1\mathrm{}𝑑\stackrel{}{x}_\mathrm{N}P(\stackrel{}{x}_1,\stackrel{}{x}_2,\mathrm{},\stackrel{}{x}_\mathrm{N})\mathrm{log}_2P(\stackrel{}{x}_1,\stackrel{}{x}_2,\mathrm{},\stackrel{}{x}_\mathrm{N})}.`$
By analogy with Eq. (21) we then write
$`P(\stackrel{}{x}_1,\stackrel{}{x}_2,\mathrm{},\stackrel{}{x}_\mathrm{N})={\displaystyle d^K\alpha 𝒫(𝜶)\underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}Q(\stackrel{}{x}_\mathrm{i}|𝜶)}.`$ (43)
Next, combining the last two equations and rearranging the order of integration, we can rewrite $`S(N)`$ as
$`S(N)`$ $`=`$ $`{\displaystyle d^K\overline{𝜶}𝒫(\overline{𝜶})\left\{𝑑\stackrel{}{x}_1\mathrm{}𝑑\stackrel{}{x}_\mathrm{N}\underset{\mathrm{j}=1}{\overset{\mathrm{N}}{}}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶})\mathrm{log}_2P(\{\stackrel{}{x}_\mathrm{i}\})\right\}}.`$
Equation (LABEL:Srewritten) allows an easy interpretation. There is the ‘true’ set of parameters $`\overline{𝜶}`$ that gave rise to the data sequence $`\stackrel{}{x}_1\mathrm{}\stackrel{}{x}_\mathrm{N}`$ with the probability $`_{\mathrm{j}=1}^\mathrm{N}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶})`$. We need to average $`\mathrm{log}_2P(\stackrel{}{x}_1\mathrm{}\stackrel{}{x}_\mathrm{N})`$ first over all possible realizations of the data keeping the true parameters fixed, and then over the parameters $`\overline{𝜶}`$ themselves. With this interpretation in mind, the joint probability density, the logarithm of which is being averaged, can be rewritten in the following useful way:
$`P(\stackrel{}{x}_1,\mathrm{},\stackrel{}{x}_\mathrm{N})`$ $`=`$ $`{\displaystyle \underset{\mathrm{j}=1}{\overset{\mathrm{N}}{}}}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶}){\displaystyle d^K\alpha 𝒫(𝜶)\underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}\left[\frac{Q(\stackrel{}{x}_\mathrm{i}|𝜶)}{Q(\stackrel{}{x}_\mathrm{i}|\overline{𝜶})}\right]}`$ (45)
$`=`$ $`{\displaystyle \underset{\mathrm{j}=1}{\overset{\mathrm{N}}{}}}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶}){\displaystyle d^K\alpha 𝒫(𝜶)\mathrm{exp}\left[N_N(𝜶;\{\stackrel{}{x}_\mathrm{i}\})\right]},`$
$`_N(𝜶;\{\stackrel{}{x}_\mathrm{i}\})`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{\mathrm{i}=1}{\overset{\mathrm{N}}{}}}\mathrm{ln}\left[{\displaystyle \frac{Q(\stackrel{}{x}_\mathrm{i}|𝜶)}{Q(\stackrel{}{x}_\mathrm{i}|\overline{𝜶})}}\right].`$ (46)
Since, by our interpretation, $`\overline{𝜶}`$ are the true parameters that gave rise to the particular data $`\{\stackrel{}{x}_\mathrm{i}\}`$, we may expect that the empirical average in Eq. (46) will converge to the corresponding expectation value, so that
$$_N(𝜶;\{\stackrel{}{x}_\mathrm{i}\})=d^DxQ(x|\overline{𝜶})\mathrm{ln}\left[\frac{Q(\stackrel{}{x}|𝜶)}{Q(\stackrel{}{x}|\overline{𝜶})}\right]\psi (𝜶,\overline{𝜶};\{x_\mathrm{i}\}),$$
(47)
where $`\psi 0`$ as $`N\mathrm{}`$; here we neglect $`\psi `$, and return to this term below.
The first term on the right hand side of Eq. (47) is the Kullback–Leibler divergence, $`D_{\mathrm{KL}}(\overline{𝜶}||𝜶)`$, between the true distribution characterized by parameters $`\overline{𝜶}`$ and the possible distribution characterized by $`𝜶`$. Thus at large $`N`$ we have
$`P(\stackrel{}{x}_1,\stackrel{}{x}_2,\mathrm{},\stackrel{}{x}_\mathrm{N})\stackrel{~}{}{\displaystyle \underset{\mathrm{j}=1}{\overset{\mathrm{N}}{}}}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶}){\displaystyle }d^K\alpha 𝒫(𝜶)\mathrm{exp}[ND_{\mathrm{KL}}(\overline{𝜶}||𝜶)],`$ (48)
where again the notation $`\stackrel{~}{}`$ reminds us that we are not only taking the limit of large $`N`$ but also making another approximation in neglecting fluctuations. By the same arguments as above we can proceed (formally) to compute the entropy of this distribution, and we find
$`S(N)`$ $``$ $`𝒮_0N+S_1^{(\mathrm{a})}(N),`$ (49)
$`𝒮_0`$ $`=`$ $`{\displaystyle d^K\alpha 𝒫(𝜶)\left[d^DxQ(\stackrel{}{x}|𝜶)\mathrm{log}_2Q(\stackrel{}{x}|𝜶)\right]},\mathrm{and}`$ (50)
$`S_1^{(\mathrm{a})}(N)`$ $`=`$ $`{\displaystyle d^K\overline{\alpha }𝒫(\overline{𝜶})\mathrm{log}_2\left[d^K\alpha P(𝜶)\mathrm{e}^{ND_{\mathrm{KL}}(\overline{𝜶}||𝜶)}\right]}.`$ (51)
Here $`S_1^{(\mathrm{a})}`$ is an approximation to $`S_1`$ that neglects fluctuations $`\psi `$. This is the same as the annealed approximation in the statistical mechanics of disordered systems, as has been used widely in the study of supervised learning problems (Seung et al. 1992). Thus we can identify the particular data sequence $`\stackrel{}{x}_1\mathrm{}\stackrel{}{x}_\mathrm{N}`$ with the disorder, $`_N(𝜶;\{\stackrel{}{x}_\mathrm{i}\})`$ with the energy of the quenched system, and $`D_{\mathrm{KL}}(\overline{𝜶}||𝜶)`$ with its annealed analogue.
The extensive term $`𝒮_0`$, Eq. (50), is the average entropy of a distribution in our family of possible distributions, generalizing the result of Eq. (40). The subextensive terms in the entropy are controlled by the $`N`$ dependence of the partition function
$$Z(\overline{𝜶};N)=d^K\alpha 𝒫(𝜶)\mathrm{exp}[ND_{\mathrm{KL}}(\overline{𝜶}||𝜶)],$$
(52)
and $`S_1(N)=\mathrm{log}_2Z(\overline{𝜶};N)_{\overline{𝜶}}`$ is analogous to the free energy. Since what is important in this integral is the Kullback–Leibler (KL) divergence between different distributions, it is natural to ask about the density of models that are KL divergence $`D`$ away from the target $`\overline{𝜶}`$,
$`\rho (D;\overline{𝜶})={\displaystyle }d^K\alpha 𝒫(𝜶)\delta [DD_{\mathrm{KL}}(\overline{𝜶}||𝜶)];`$ (53)
this density could be very different for different targets.<sup>6</sup><sup>6</sup>6If parameter space is compact, then a related description of the space of targets based on metric entropy, also called Kolmogorov’s $`ϵ`$-entropy, is used in the literature \[see, for example, Haussler and Opper (1997)\]. Metric entropy is the logarithm of the smallest number of disjoint parts of the size not greater than $`ϵ`$ into which the target space can be partitioned. The density of divergences is normalized because the original distribution over parameter space, $`P(𝜶)`$, is normalized,
$`{\displaystyle 𝑑D\rho (D;\overline{𝜶})}={\displaystyle d^K\alpha 𝒫(𝜶)}=1.`$ (54)
Finally, the partition function takes the simple form
$$Z(\overline{𝜶};N)=𝑑D\rho (D;\overline{𝜶})\mathrm{exp}[ND].$$
(55)
We recall that in statistical mechanics the partition function is given by
$`Z(\beta )={\displaystyle 𝑑E\rho (E)\mathrm{exp}[\beta E]},`$ (56)
where $`\rho (E)`$ is the density of states that have energy $`E`$, and $`\beta `$ is the inverse temperature. Thus the subextensive entropy in our learning problem is analogous to a system in which energy corresponds to the Kullback–Leibler divergence relative to the target model, and temperature is inverse to the number of examples. As we increase the length $`N`$ of the time series we have observed, we “cool” the system and hence probe models which approach the target; the dynamics of this approach is determined by the density of low energy states, that is the behavior of $`\rho (D;\overline{𝜶})`$ as $`D0`$.<sup>7</sup><sup>7</sup>7It may be worth emphasizing that this analogy to statistical mechanics emerges from the analysis of the relevant probability distributions, rather than being imposed on the problem through some assumptions about the nature of the learning algorithm.
The structure of the partition function is determined by a competition between the (Boltzmann) exponential term, which favors models with small $`D`$, and the density term, which favors values of $`D`$ that can be achieved by the largest possible number of models. Because there (typically) are many parameters, there are very few models with $`D0`$. This picture of competition between the Boltzmann factor and a density of states has been emphasized in previous work on supervised learning (Haussler et al. 1996).
The behavior of the density of states, $`\rho (D;\overline{𝜶})`$, at small $`D`$ is related to the more intuitive notion of dimensionality. In a parameterized family of distributions, the Kullback–Leibler divergence between two distributions with nearby parameters is approximately a quadratic form,
$`D_{\mathrm{KL}}(\overline{𝜶}||𝜶){\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu \nu }{}}(\overline{\alpha }_\mu \alpha _\mu )_{\mu \nu }(\overline{\alpha }_\nu \alpha _\nu )+\mathrm{},`$ (57)
where $`N`$ is the Fisher information matrix. Intuitively, if we have a reasonable parameterization of the distributions, then similar distributions will be nearby in parameter space, and more importantly points that are far apart in parameter space will never correspond to similar distributions; Clarke and Barron (1990) refer to this condition as the parameterization forming a “sound” family of distributions. If this condition is obeyed, then we can approximate the low $`D`$ limit of the density $`\rho (D;\overline{𝜶})`$:
$`\rho (D;\overline{𝜶})`$ $`=`$ $`{\displaystyle }d^K\alpha 𝒫(𝜶)\delta [DD_{\mathrm{KL}}(\overline{𝜶}||𝜶)]`$ (58)
$``$ $`{\displaystyle d^K\alpha 𝒫(𝜶)\delta \left[D\frac{1}{2}\underset{\mu \nu }{}(\overline{\alpha }_\mu \alpha _\mu )_{\mu \nu }(\overline{\alpha }_\nu \alpha _\nu )\right]}`$
$`=`$ $`{\displaystyle d^K\alpha 𝒫(\overline{𝜶}+𝒰𝝃)\delta \left[D\frac{1}{2}\underset{\mu }{}\mathrm{\Lambda }_\mu \xi _\mu ^2\right]},`$
where $`𝒰`$ is a matrix that diagonalizes $``$,
$`(𝒰^T𝒰)_{\mu \nu }=\mathrm{\Lambda }_\mu \delta _{\mu \nu }.`$ (59)
The delta function restricts the components of $`𝝃`$ in Eq. (58) to be of order $`\sqrt{D}`$ or less, and so if $`P(𝜶)`$ is smooth we can make a perturbation expansion. After some algebra the leading term becomes
$`\rho (D0;\overline{𝜶})𝒫(\overline{𝜶}){\displaystyle \frac{2\pi ^{K/2}}{\mathrm{\Gamma }(K/2)}}\left(det\right)^{1/2}D^{(K2)/2}.`$ (60)
Here, as before, $`K`$ is the dimensionality of the parameter vector. Computing the partition function from Eq. (55), we find
$$Z(\overline{𝜶};N\mathrm{})f(\overline{𝜶})\frac{\mathrm{\Gamma }(K/2)}{N^{K/2}},$$
(61)
where $`f(\overline{𝜶})`$ is some function of the target parameter values. Finally, this allows us to evaluate the subextensive entropy, from Eqs. (51, 52):
$`S_1^{(\mathrm{a})}(N)`$ $`=`$ $`{\displaystyle d^K\overline{\alpha }𝒫(\overline{𝜶})\mathrm{log}_2Z(\overline{𝜶};N)}`$ (62)
$``$ $`{\displaystyle \frac{K}{2}}\mathrm{log}_2N+\mathrm{}(\mathrm{bits}),`$ (63)
where $`\mathrm{}`$ are finite as $`N\mathrm{}`$. Thus, general $`K`$–parameter model classes have the same subextensive entropy as for the simplest example considered in the previous section. To the leading order, this result is independent even of the prior distribution $`𝒫(𝜶)`$ on the parameter space, so that the predictive information seems to count the number of parameters under some very general conditions \[cf. Fig. (3) for a very different numerical example of the logarithmic behavior\].
Although Eq. (63) is true under a wide range of conditions, this cannot be the whole story. Much of modern learning theory is concerned with the fact that counting parameters is not quite enough to characterize the complexity of a model class; the naive dimension of the parameter space $`K`$ should be viewed in conjunction with the pseudodimension (also known as the shattering dimension or Vapnik–Chervonenkis dimension $`d_{\mathrm{VC}}`$), which measures capacity of the model class, and with the phase space dimension $`d`$, which accounts for volumes in the space of models (Vapnik 1998, Opper 1994). Both of these dimensions can differ from the number of parameters. One possibility is that $`d_{\mathrm{VC}}`$ is infinite when the number of parameters is finite, a problem discussed below. Another possibility is that the determinant of $``$ is zero, and hence $`d_{\mathrm{VC}}`$ and $`d`$ both are smaller than the number of parameters because we have adopted a redundant description. This sort of degeneracy could occur over a finite fraction but not all of the parameter space, and this is one way to generate an effective fractional dimensionality. One can imagine multifractal models such that the effective dimensionality varies continuously over the parameter space, but it is not obvious where this would be relevant. Finally, models with $`d<d_{\mathrm{VC}}<\mathrm{}`$ also are possible \[see, for example, Opper (1994)\], and this list probably is not exhaustive.
The calculation above, Eq. (60), lets us actually define the phase space dimension through the exponent in the small $`D_{\mathrm{KL}}`$ behavior of the model density,
$`\rho (D0;\overline{𝜶})D^{(d2)/2},`$ (64)
and then $`d`$ appears in place of $`K`$ as the coefficient of the log divergence in $`S_1(N)`$ (Clarke and Barron 1990, Opper 1994). However, this simple conclusion can fail in two ways. First, it can happen that the density $`\rho (D;\overline{𝜶})`$ accumulates a macroscopic weight at some nonzero value of $`D`$, so that the small $`D`$ behavior is irrelevant for the large $`N`$ asymptotics. Second, the fluctuations neglected here may be uncontrollably large, so that the asymptotics are never reached. Since controllability of fluctuations is a function of $`d_{\mathrm{VC}}`$ (see Vapnik 1998, Haussler et al. 1994, and later in this paper), we may summarize this in the following way. Provided that the small $`D`$ behavior of the density function is the relevant one, the coefficient of the logarithmic divergence of $`I_{\mathrm{pred}}`$ measures the phase space or the scaling dimension $`d`$ and nothing else. This asymptote is valid, however, only for $`Nd_{\mathrm{VC}}`$. It is still an open question whether the two pathologies that can violate this asymptotic behavior are related.
### 4.3 Learning a parameterized process
Consider a process where samples are not independent, and our task is to learn their joint distribution $`Q(\stackrel{}{x}_1,\mathrm{},\stackrel{}{x}_\mathrm{N}|𝜶)`$. Again, $`𝜶`$ is an unknown parameter vector which is chosen randomly at the beginning of the series. If $`𝜶`$ is a $`K`$ dimensional vector, then one still tries to learn just $`K`$ numbers and there are still $`N`$ examples, even if there are correlations. Therefore, although such problems are much more general than those considered above, it is reasonable to expect that the predictive information still is measured by $`(K/2)\mathrm{log}_2N`$ provided that some conditions are met.
One might suppose that conditions for simple results on the predictive information are very strong, for example that the distribution $`Q`$ is a finite order Markov model. In fact all we really need are the following two conditions:
$`S\left[\{\stackrel{}{x}_\mathrm{i}\}|𝜶\right]`$ $``$ $`{\displaystyle d^N\stackrel{}{x}Q(\{\stackrel{}{x}_\mathrm{i}\}|𝜶)\mathrm{log}_2Q(\{\stackrel{}{x}_\mathrm{i}\}|𝜶)}`$ (65)
$``$ $`N𝒮_0+𝒮_0^{};𝒮_0^{}=O(1),`$
$`D_{\mathrm{KL}}[Q(\{\stackrel{}{x}_\mathrm{i}\}|\overline{𝜶})||Q(\{\stackrel{}{x}_\mathrm{i}\}|𝜶)]`$ $``$ $`N𝒟_{\mathrm{KL}}(\overline{𝜶}||𝜶)+o(N).`$ (66)
Here the quantities $`𝒮_0`$, $`𝒮_0^{}`$, and $`𝒟_{\mathrm{KL}}(\overline{𝜶}||𝜶)`$ are defined by taking limits $`N\mathrm{}`$ in both equations. The first of the constraints limits deviations from extensivity to be of order unity, so that if $`𝜶`$ is known there are no long range correlations in the data—all of the long range predictability is associated with learning the parameters.<sup>8</sup><sup>8</sup>8Suppose that we observe a Gaussian stochastic process and we try to learn the power spectrum. If the class of possible spectra includes ratios of polynomials in the frequency (rational spectra) then this condition is met. On the other hand, if the class of possible spectra includes $`1/f`$ noise, then the condition may not be met. For more on long range correlations, see below. The second constraint, Eq. (66), is a less restrictive one, and it ensures that the “energy” of our statistical system is an extensive quantity.
With these conditions it is straightforward to show that the results of the previous subsection carry over virtually unchanged. With the same cautious statements about fluctuations and the distinction between $`K`$, $`d`$, and $`d_{\mathrm{VC}}`$, one arrives at the result:
$`S(N)`$ $`=`$ $`𝒮_0N+S_1^{(\mathrm{a})}(N),`$ (67)
$`S_1^{(\mathrm{a})}(N)`$ $`=`$ $`{\displaystyle \frac{K}{2}}\mathrm{log}_2N+\mathrm{}(\mathrm{bits}),`$ (68)
where $`\mathrm{}`$ stands for terms of order one as $`N\mathrm{}`$. Note again that for the results Eq. (68) to be valid, the process considered is not required to be a finite order Markov process. Memory of all previous outcomes may be kept, provided that the accumulated memory does not contribute a divergent term to the subextensive entropy.
It is interesting to ask what happens if the condition in Eq. (65) is violated, so that there are long range correlations even in the conditional distribution $`Q(\stackrel{}{x}_1,\mathrm{},\stackrel{}{x}_\mathrm{N}|𝜶)`$. Suppose, for example, that $`𝒮_0^{}=(K^{}/2)\mathrm{log}_2N`$. Then the subextensive entropy becomes
$$S_1^{(\mathrm{a})}(N)=\frac{K+K^{}}{2}\mathrm{log}_2N+\mathrm{}(\mathrm{bits}).$$
(69)
We see the that the subextensive entropy makes no distinction between predictability that comes from unknown parameters and predictability that comes from intrinsic correlations in the data; in this sense, two models with the same $`K+K^{}`$ are equivalent. This, actually, must be so. As an example, consider a chain of Ising spins with long range interactions in one dimension. This system can order (magnetize) and exhibit long range correlations, and so the predictive information will diverge at the transition to ordering. In one view, there is no global parameter analogous to $`𝜶`$, just the long range interactions. On the other hand, there are regimes in which we can approximate the effect of these interactions by saying that all the spins experience a mean field which is constant across the whole length of the system, and then formally we can think of the predictive information as being carried by the mean field itself. In fact there are situations in which this is not just an approximation, but an exact statement. Thus we can trade a description in terms of long range interactions ($`K^{}0`$, but $`K=0`$) for one in which there are unknown parameters describing the system but given these parameters there are no long range correlations ($`K0,K^{}=0`$). The two descriptions are equivalent, and this is captured by the subextensive entropy.<sup>9</sup><sup>9</sup>9There are a number of interesting questions about how the coefficients in the diverging predictive information relate to the usual critical exponents, and we hope to return to this problem in a later paper.
### 4.4 Taming the fluctuations
The preceding calculations of the subextensive entropy $`S_1`$ are worthless unless we prove that the fluctuations $`\psi `$ are controllable. In this subsection we are going to discuss when and if this, indeed, happens. We limit the discussion to the case of finding a probability density (Section 4.2); the case of learning a process (Section 4.3) is very similar.
Clarke and Barron (1990) solved essentially the same problem. They did not make a separation into the annealed and the fluctuation term, and the quantity they were interested in was a bit different from ours, but, interpreting loosely, they proved that, modulo some reasonable technical assumptions on differentiability of functions in question, the fluctuation term always approaches zero. However, they did not investigate the speed of this approach, and we believe that, by doing so, they missed some important qualitative distinctions between different problems that can arise due to a difference between $`d`$ and $`d_{\mathrm{VC}}`$. In order to illuminate these distinctions, we here go through the trouble of analyzing fluctuations all over again.
Returning to Eqs. (45, 47) and the definition of entropy, we can write the entropy $`S(N)`$ exactly as
$`S(N)`$ $`=`$ $`{\displaystyle d^K\overline{\alpha }𝒫(\overline{𝜶})\underset{j=1}{\overset{N}{}}\left[d\stackrel{}{x}_\mathrm{j}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶})\right]}`$
$`\times \mathrm{log}_2\left[{\displaystyle \underset{i=1}{\overset{N}{}}}Q(\stackrel{}{x}_\mathrm{i}|\overline{𝜶}){\displaystyle d^K\alpha 𝒫(𝜶)\mathrm{e}^{ND_{\mathrm{KL}}(\overline{𝜶}||𝜶)+N\psi (𝜶,\overline{𝜶};\{\stackrel{}{x}_\mathrm{i}\})}}\right].`$
This expression can be decomposed into the terms identified above, plus a new contribution to the subextensive entropy that comes from the fluctuations alone, $`S_1^{(\mathrm{f})}(N)`$:
$`S(N)`$ $`=`$ $`𝒮_0N+S_1^{(\mathrm{a})}(N)+S_1^{(\mathrm{f})}(N),`$ (71)
$`S_1^{(\mathrm{f})}`$ $`=`$ $`{\displaystyle d^K\overline{\alpha }𝒫(\overline{𝜶})\underset{j=1}{\overset{N}{}}\left[d\stackrel{}{x}_\mathrm{j}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶})\right]}`$ (72)
$`\times \mathrm{log}_2\left[{\displaystyle \frac{d^K\alpha 𝒫(𝜶)}{Z(\overline{𝜶};N)}\mathrm{e}^{ND_{\mathrm{KL}}(\overline{𝜶}||𝜶)+N\psi (𝜶,\overline{𝜶};\{\stackrel{}{x}_\mathrm{i}\})}}\right],`$
where $`\psi `$ is defined as in Eq. (47), and $`Z`$ as in Eq. (52).
Some loose but useful bounds can be established. First, the predictive information is a positive (semidefinite) quantity, and so the fluctuation term may not be smaller (more negative) than the value of $`S_1^{(\mathrm{a})}`$ as calculated in Eqs. (63, 68). Second, since fluctuations make it more difficult to generalize from samples, the predictive information should always be reduced by fluctuations, so that $`S^{(\mathrm{f})}`$ is negative. This last statement corresponds to the fact that for the statistical mechanics of disordered systems, the annealed free energy always is less than the average quenched free energy, and may be proven rigorously by applying Jensen’s inequality to the (concave) logarithm function in Eq. (72); essentially the same argument was given by Haussler and Opper (1997). A related Jensen’s inequality argument allows us to show that the total $`S_1(N)`$ is bounded,
$`S_1(N)`$ $``$ $`N{\displaystyle }d^K\alpha {\displaystyle }d^K\overline{\alpha }𝒫(𝜶)𝒫(\overline{𝜶})D_{\mathrm{KL}}(\overline{𝜶}||𝜶)`$ (73)
$``$ $`ND_{\mathrm{KL}}(\overline{𝜶}||𝜶)_{\overline{𝜶},𝜶},`$
so that if we have a class of models (and a prior $`𝒫(𝜶)`$) such that the average Kullback–Leibler divergence among pairs of models is finite, then the subextensive entropy necessarily is properly defined. In its turn, finiteness of the average KL divergence or of similar quantities is a usual constraint in the learning problems of this type \[see, for example, Haussler and Opper (1997)\]. Note also that $`D_{\mathrm{KL}}(\overline{𝜶}||𝜶)_{\overline{𝜶},𝜶}`$ includes $`𝒮_0`$ as one of its terms, so that usually $`𝒮_0`$ and $`S_1`$ are well– or ill–defined together.
Tighter bounds require nontrivial assumptions about the class of distributions considered. The fluctuation term would be zero if $`\psi `$ were zero, and $`\psi `$ is the difference between an expectation value (KL divergence) and the corresponding empirical mean. There is a broad literature that deals with this type of difference (see, for example, Vapnik 1998).
We start with the case when the pseudodimension ($`d_{\mathrm{VC}}`$) of the set of probability densities $`\{Q(\stackrel{}{x}|𝜶)\}`$ is finite. Then for any reasonable function $`F(\stackrel{}{x};\beta )`$, deviations of the empirical mean from the expectation value can be bounded by probabilistic bounds of the form
$`P\left\{\underset{\beta }{sup}\left|{\displaystyle \frac{\frac{1}{N}_\mathrm{j}F(\stackrel{}{x}_\mathrm{j};\beta )𝑑\stackrel{}{x}Q(\stackrel{}{x}|\overline{𝜶})F(\stackrel{}{x};\beta )}{L[F]}}\right|>ϵ\right\}`$
$`<M(ϵ,N,d_{\mathrm{VC}})\mathrm{e}^{cNϵ^2},`$ (74)
where $`c`$ and $`L[F]`$ depend on the details of the particular bound used. Typically, $`c`$ is a constant of order one, and $`L[F]`$ either is some moment of $`F`$ or the range of its variation. In our case, $`F`$ is the log–ratio of two densities, so that $`L[F]`$ may be assumed bounded for almost all $`\beta `$ without loss of generality in view of Eq. (73). In addition, $`M(ϵ,N,d_{\mathrm{VC}})`$ is finite at zero, grows at most subexponentially in its first two arguments, and depends exponentially on $`d_{\mathrm{VC}}`$. Bounds of this form may have different names in different contexts: Glivenko–Cantelli, Vapnik–Chervonenkis, Hoeffding, Chernoff, …; for review see Vapnik (1998) and the references therein.
To start the proof of finiteness of $`S_1^{(\mathrm{f})}`$ in this case, we first show that only the region $`𝜶\overline{𝜶}`$ is important when calculating the inner integral in Eq. (72). This statement is equivalent to saying that at large values of $`𝜶\overline{𝜶}`$ the KL divergence almost always dominates the fluctuation term, that is, the contribution of sequences of $`\{\stackrel{}{x}_\mathrm{i}\}`$ with atypically large fluctuations is negligible (atypicality is defined as $`\psi \delta `$, where $`\delta `$ is some small constant independent of $`N`$). Since the fluctuations decrease as $`1/\sqrt{N}`$ \[see Eq. (74)\], and $`D_{\mathrm{KL}}`$ is of order one, this is plausible. To show this, we bound the logarithm in Eq. (72) by $`N`$ times the supremum value of $`\psi `$. Then we realize that the averaging over $`\overline{𝜶}`$ and $`\{\stackrel{}{x}_\mathrm{i}\}`$ is equivalent to integration over all possible values of the fluctuations. The worst case density of the fluctuations may be estimated by differentiating Eq. (74) with respect to $`ϵ`$ (this brings down an extra factor of $`Nϵ`$). Thus the worst case contribution of these atypical sequences is
$`S_1^{(\mathrm{f}),\mathrm{atypical}}{\displaystyle _\delta ^{\mathrm{}}}𝑑ϵN^2ϵ^2M(ϵ)\mathrm{e}^{cNϵ^2}\mathrm{e}^{cN\delta ^2}1\mathrm{for}\mathrm{large}N.`$ (75)
This bound lets us focus our attention on the region $`𝜶\overline{𝜶}`$. We expand the exponent of the integrand of Eq. (72) around this point and perform a simple Gaussian integration. In principle, large fluctuations might lead to an instability (positive or zero curvature) at the saddle point, but this is atypical and therefore is accounted for already. Curvatures at the saddle points of both numerator and denominator are of the same order, and throwing away unimportant additive and multiplicative constants of order unity, we obtain the following result for the contribution of typical sequences:
$`S_1^{(\mathrm{f}),\mathrm{typical}}`$ $``$ $`{\displaystyle d^K\overline{\alpha }𝒫(\overline{𝜶})d^N\stackrel{}{x}\underset{\mathrm{j}}{}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶})N(𝐁𝒜^1𝐁)};`$ (76)
$`B_\mu `$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{\mathrm{i}}{}}{\displaystyle \frac{\mathrm{log}Q(\stackrel{}{x}_\mathrm{i}|\overline{𝜶})}{\overline{\alpha }_\mu }},𝐁_\stackrel{}{x}=0;`$
$`(𝒜)_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{\mathrm{i}}{}}{\displaystyle \frac{^2\mathrm{log}Q(\stackrel{}{x}_\mathrm{i}|\overline{𝜶})}{\overline{\alpha }_\mu \overline{\alpha }_\nu }},𝒜_\stackrel{}{x}=.`$
Here $`\mathrm{}_\stackrel{}{x}`$ means an averaging with respect to all $`\stackrel{}{x}_\mathrm{i}`$’s keeping $`\overline{𝜶}`$ constant. One immediately recognizes that $`𝐁`$ and $`𝒜`$ are, respectively, first and second derivatives of the empirical KL divergence that was in the exponent of the inner integral in Eq. (72).
We are dealing now with typical cases. Therefore, large deviations of $`𝒜`$ from $``$ are not allowed, and we may bound Eq. (76) by replacing $`𝒜^1`$ with $`^1(1+\delta )`$, where $`\delta `$ again is independent of $`N`$. Now we have to average a bunch of products like
$`{\displaystyle \frac{\mathrm{log}Q(\stackrel{}{x}_\mathrm{i}|\overline{𝜶})}{\overline{\alpha }_\mu }}(^1)_{\mu \nu }{\displaystyle \frac{\mathrm{log}Q(\stackrel{}{x}_\mathrm{j}|\overline{𝜶})}{\overline{\alpha }_\nu }}`$ (77)
over all $`\stackrel{}{x}_\mathrm{i}`$’s. Only the terms with $`\mathrm{i}=\mathrm{j}`$ survive the averaging. There are $`K^2N`$ such terms, each contributing of order $`N^1`$. This means that the total contribution of the typical fluctuations is bounded by a number of order one and does not grow with $`N`$.
This concludes the proof of controllability of fluctuations for $`d_{\mathrm{VC}}<\mathrm{}`$. One may notice that we never used the specific form of $`M(ϵ,N,d_{\mathrm{VC}})`$, which is the only thing dependent on the precise value of the dimension. Actually, a more thorough look at the proof shows that we do not even need the strict uniform convergence enforced by the Glivenko–Cantelli bound. With some modifications the proof should still hold if there exist some a priori improbable values of $`𝜶`$ and $`\overline{𝜶}`$ that lead to violation of the bound. That is, if the prior $`𝒫(𝜶)`$ has sufficiently narrow support, then we may still expect fluctuations to be unimportant even for VC–infinite problems.
A proof of this can be found in the realm of the Structural Risk Minimization (SRM) theory (Vapnik 1998). SRM theory assumes that an infinite structure $`𝒞`$ of nested subsets $`C_1C_2C_3\mathrm{}`$ can be imposed onto the set $`C`$ of all admissible solutions of a learning problem, such that $`C=C_\mathrm{n}`$. The idea is that, having a finite number of observations $`N`$, one is confined to the choices within some particular structure element $`C_\mathrm{n},\mathrm{n}=\mathrm{n}(N),`$ when looking for an approximation to the true solution; this prevents overfitting and poor generalization. Then, as the number of samples increases and one is able to distinguish within more and more complicated subsets, $`\mathrm{n}`$ grows. If $`d_{\mathrm{VC}}`$ for learning in any $`C_\mathrm{n},\mathrm{n}<\mathrm{},`$ is finite, then one can show convergence of the estimate to the true value as well as the absolute smallness of fluctuations (Vapnik 1998). It is remarkable that this result holds even if the capacity of the whole set $`C`$ is not described by a finite $`d_{\mathrm{VC}}`$.
In the context of SRM, the role of the prior $`P(𝜶)`$ is to induce a structure on the set of all admissible densities, and the fight between the number of samples $`N`$ and the narrowness of the prior is precisely what determines how the capacity of the current element of the structure $`C_\mathrm{n},\mathrm{n}=\mathrm{n}(N),`$ grows with $`N`$. A rigorous proof of smallness of the fluctuations can be constructed based on well known results, as detailed elsewhere (Nemenman 2000). Here we focus on the question of how narrow the prior should be so that every structure element is of finite VC–dimension, and one can guarantee eventual convergence of fluctuations to zero.
Consider two examples. A variable $`x`$ is distributed according to the following probability density functions:
$`Q(x|\alpha )`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}\left(x\alpha \right)^2\right],x(\mathrm{};+\mathrm{});`$ (78)
$`Q(x|\alpha )`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}\left(\mathrm{sin}\alpha x\right)}{_0^{2\pi }𝑑x\mathrm{exp}\left(\mathrm{sin}\alpha x\right)}},x[0;2\pi ).`$ (79)
Learning the parameter in the first case is a $`d_{\mathrm{VC}}=1`$ problem, while in the second case $`d_{\mathrm{VC}}=\mathrm{}`$. In the first example, as we have shown above, one may construct a uniform bound on fluctuations irrespective of the prior $`P(𝜶)`$. The second one does not allow this. Suppose that the prior is uniform in a box $`0<\alpha <\alpha _{\mathrm{max}}`$, and zero elsewhere, with $`\alpha _{\mathrm{max}}`$ rather large. Then for not too many sample points $`N`$, the data would be better fitted not by some value in the vicinity of the actual parameter, but by some much larger value, for which almost all data points are at the crests of $`\mathrm{sin}\alpha x`$. Adding a new data point would not help, until that best, but wrong, parameter estimate is less than $`\alpha _{\mathrm{max}}`$.<sup>10</sup><sup>10</sup>10Interestingly, since for the model Eq. (79) KL divergence is bounded from below and from above, for $`\alpha _{\mathrm{max}}\mathrm{}`$ the weight in $`\rho (D;\overline{𝜶})`$ at small $`D_{\mathrm{KL}}`$ vanishes, and a finite weight accumulates at some nonzero value of $`D`$. Thus, even putting the fluctuations aside, the asymptotic behavior based on the phase space dimension is invalidated, as mentioned above. So the fluctuations are large, and the predictive information is small in this case. Eventually, however, data points would overwhelm the box size, and the best estimate of $`\alpha `$ would swiftly approach the actual value. At this point the argument of Clarke and Barron would become applicable, and the leading behavior of the subextensive entropy would converge to its asymptotic value of $`(1/2)\mathrm{log}N`$. On the other hand, there is no uniform bound on the value of $`N`$ for which this convergence will occur—it is guaranteed only for $`Nd_{\mathrm{VC}}`$, which is never true if $`d_{\mathrm{VC}}=\mathrm{}`$. For some sufficiently wide priors this asymptotically correct behavior would be never reached in practice. Further, if we imagine a thermodynamic limit where the box size and the number of samples both become large, then by analogy with problems in supervised learning (Seung et al. 1992, Haussler et al. 1996) we expect that there can be sudden changes in performance as a function of the number of examples. The arguments of Clarke and Barron cannot encompass these phase transitions or “aha!” phenomena. A further bridge between VC dimension and the information theoretic approach to learning may be found in Haussler et al. (1994), where the authors bounded predictive information–like quantities with loose but useful bounds explicitly proportional to $`d_{\mathrm{VC}}`$.
While much of learning theory has focused on problems with finite VC dimension, it might be that the conventional scenario in which the number of examples eventually overwhelms the number of parameters or dimensions is too weak to deal with many real world problems. Certainly in the present context there is not only a quantitative, but also a qualitative difference between reaching the asymptotic regime in just a few measurements, or in many millions of them. Finitely parameterizable models with finite or infinite $`d_{\mathrm{VC}}`$ fall in essentially different universality classes with respect to the predictive information.
### 4.5 Beyond finite parameterization: <br>general considerations
The previous sections have considered learning from time series where the underlying class of possible models is described with a finite number of parameters. If the number of parameters is not finite then in principle it is impossible to learn anything unless there is some appropriate regularization of the problem. If we let the number of parameters stay finite but become large, then there is more to be learned and correspondingly the predictive information grows in proportion to this number, as in Eq. (63). On the other hand, if the number of parameters becomes infinite without regularization, then the predictive information should go to zero since nothing can be learned. We should be able to see this happen in a regularized problem as the regularization weakens: eventually the regularization would be insufficient and the predictive information would vanish. The only way this can happen is if the subextensive term in the entropy grows more and more rapidly with $`N`$ as we weaken the regularization, until finally it becomes extensive at the point where learning becomes impossible. More precisely, if this scenario for the breakdown of learning is to work, there must be situations in which the predictive information grows with $`N`$ more rapidly than the logarithmic behavior found in the case of finite parameterization.
Subextensive terms in the entropy are controlled by the density of models as a function of their Kullback–Leibler divergence to the target model. If the models have finite VC and phase space dimensions then this density vanishes for small divergences as $`\rho D^{(d2)/2}`$. Phenomenologically, if we let the number of parameters increase, the density vanishes more and more rapidly. We can imagine that beyond the class of finitely parameterizable problems there is a class of regularized infinite dimensional problems in which the density $`\rho (D0)`$ vanishes more rapidly than any power of $`D`$. As an example, we could have
$`\rho (D0)A\mathrm{exp}\left[{\displaystyle \frac{B}{D^\mu }}\right],\mu >0;`$ (80)
that is, an essential singularity at $`D=0`$. For simplicity we assume that the constants $`A`$ and $`B`$ can depend on the target model, but that the nature of the essential singularity ($`\mu `$) is the same everywhere. Before providing an explicit example, let us explore the consequences of this behavior.
From Equation (55) above, we can write the partition function as
$`Z(\overline{𝜶};N)`$ $`=`$ $`{\displaystyle 𝑑D\rho (D;\overline{𝜶})\mathrm{exp}[ND]}`$ (81)
$``$ $`A(\overline{𝜶}){\displaystyle 𝑑D\mathrm{exp}\left[\frac{B(\overline{𝜶})}{D^\mu }ND\right]}`$
$``$ $`\stackrel{~}{A}(\overline{𝜶})\mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu +2}{\mu +1}}\mathrm{ln}NC(\overline{𝜶})N^{\mu /(\mu +1)}\right],`$
where in the last step we use a saddle point or steepest descent approximation which is accurate at large $`N`$, and the coefficients are
$`\stackrel{~}{A}(\overline{𝜶})`$ $`=`$ $`A(\overline{𝜶})\left({\displaystyle \frac{2\pi \mu ^{1/(\mu +1)}}{\mu +1}}\right)^{1/2}[B(\overline{𝜶})]^{1/(2\mu +2)}`$ (82)
$`C(\overline{𝜶})`$ $`=`$ $`[B(\overline{𝜶})]^{1/(\mu +1)}\left({\displaystyle \frac{1}{\mu ^{\mu /(\mu +1)}}}+\mu ^{1/(\mu +1)}\right).`$ (83)
Finally we can use Eqs. (62, 81) to compute the subextensive term in the entropy, keeping only the dominant term at large $`N`$,
$$S_1^{(\mathrm{a})}(N)\frac{1}{\mathrm{ln}2}C(\overline{𝜶})_{\overline{𝜶}}N^{\mu /(\mu +1)}(\mathrm{bits}),$$
(84)
where $`\mathrm{}_{\overline{𝜶}}`$ denotes an average over all the target models.
This behavior of the first subextensive term is qualitatively different from everything we have observed so far. A power law divergence is much stronger than a logarithmic one. Therefore, a lot more predictive information is accumulated in an “infinite parameter” (or nonparametric) system; the system is much richer and more complex, both intuitively and quantitatively.
Subextensive entropy also grows as a power law in a finitely parameterizable system with a growing number of parameters. For example, suppose that we approximate the distribution of a random variable by a histogram with $`K`$ bins, and we let $`K`$ grow with the quantity of available samples as $`KN^\nu `$. A straightforward generalization of Eq. (63) seems to imply then that $`S_1(N)N^\nu \mathrm{log}N`$ (Hall and Hannan 1988, Barron and Cover 1991). While not necessarily wrong, analogies of this type are dangerous. If the number of parameters grows constantly, then the scenario where the data overwhelms all the unknowns is far from certain. Observations may provide much less information about features that were introduced into the model at some large $`N`$ than about those that have existed since the very first measurements. Indeed, in a $`K`$–parameter system, the $`N^{\mathrm{th}}`$ sample point contributes $`K/2N`$ bits to the subextensive entropy \[cf. Eq. (63)\]. If $`K`$ changes as mentioned, the $`N^{\mathrm{th}}`$ example then carries $`N^{\nu 1}`$ bits. Summing this up over all samples, we find $`S_1^{(\mathrm{a})}N^\nu `$, and if we let $`\nu =\mu /(\mu +1)`$ we obtain Eq. (84); note that the growth of the number of parameters is slower than $`N`$ ($`\nu =\mu /(\mu +1)<1`$), which makes sense. Rissanen et al. (1992) made a similar observation. According to them, for models with increasing number of parameters, predictive codes, which are optimal at each particular $`N`$ \[cf. Eq. (11)\], provide a strictly shorter coding length than nonpredictive codes optimized for all data simultaneously. This has to be contrasted with the finite parameter model classes, for which these codes are asymptotically equal.
Power law growth of the predictive information illustrates the point made earlier about the transition from learning more to finally learning nothing as the class of investigated models becomes more complex. As $`\mu `$ increases, the problem becomes richer and more complex, and this is expressed in the stronger divergence of the first subextensive term of the entropy; for fixed large $`N`$, the predictive information increases with $`\mu `$. However, if $`\mu \mathrm{}`$ the problem is too complex for learning—in our model example the number of bins grows in proportion to the number of samples, which means that we are trying to find too much detail in the underlying distribution. As a result, the subextensive term becomes extensive and stops contributing to predictive information. Thus, at least to the leading order, predictability is lost, as promised.
### 4.6 Beyond finite parameterization: example
While literature on problems in the logarithmic class is reasonably rich, the research on establishing the power law behavior seems to be in its early stages. Some authors have found specific learning problems for which quantities similar to, but sometimes very nontrivially related to $`S_1`$, are bounded by power–law functions (Haussler and Opper 1997, 1998, Cesa–Bianchi and Lugosi 2000). Others have chosen to study finite dimensional models, for which the optimal number of parameters \[usually determined by the MDL criterion of Rissanen (1989)\] grows as a power law (Hall and Hannan 1988, Rissanen et al. 1992). In addition to the questions raised earlier about the danger of copying finite dimensional intuition to the infinite dimensional setup, these are not examples of truly nonparametric Bayesian learning. Instead these authors make use of a priori constraints to restrict learning to codes of particular structure (histogram codes), while a non–Bayesian inference is conducted within the class. Without Bayesian averaging and with restrictions on the coding strategy it may happen that a realization of the code length is substantially different from the predictive information. Similar conceptual problems plague even true nonparametric examples, as considered, for example, by Barron and Cover (1991). In summary, we don’t know of a complete calculation in which a family of power–law scalings of the predictive information is derived from a Bayesian formulation.
The discussion in the previous section suggests that we should look for power–law behavior of the predictive information in learning problems where rather than learning ever more precise values for a fixed set of parameters, we learn a progressively more detailed description—effectively increasing the number of parameters—as we collect more data. One example of such a problem is learning the distribution $`Q(x)`$ for a continuous variable $`x`$, but rather than writing a parametric form of $`Q(x)`$ we assume only that this function itself is chosen from some distribution that enforces a degree of smoothness. There are some natural connections of this problem to the methods of quantum field theory (Bialek, Callan, and Strong 1996) which we can exploit to give a complete calculation of the predictive information, at least for a class of smoothness constraints.
We write $`Q(x)=(1/l_0)\mathrm{exp}[\varphi (x)]`$ so that positivity of the distribution is automatic, and then smoothness may be expressed by saying that the ‘energy’ (or action) associated with a function $`\varphi (x)`$ is related to an integral over its derivatives, like the strain energy in a stretched string. The simplest possibility following this line of ideas is that the distribution of functions is given by
$$𝒫[\varphi (x)]=\frac{1}{𝒵}\mathrm{exp}\left[\frac{l}{2}𝑑x\left(\frac{\varphi }{x}\right)^2\right]\delta \left[\frac{1}{l_0}𝑑x\mathrm{e}^{\varphi (x)}1\right],$$
(85)
where $`𝒵`$ is the normalization constant for $`𝒫[\varphi ]`$, the delta function insures that each distribution $`Q(x)`$ is normalized, and $`l`$ sets a scale for smoothness. If distributions are chosen from this distribution, then the optimal Bayesian estimate of $`Q(x)`$ from a set of samples $`x_1,x_2,\mathrm{},x_N`$ converges to the correct answer, and the distribution at finite $`N`$ is nonsingular, so that the regularization provided by this prior is strong enough to prevent the development of singular peaks at the location of observed data points (Bialek, Callan, and Strong 1996). Further developments of the theory, including alternative choices of $`P[\varphi (x)]`$, have been given by Periwal (1997, 1998), Holy (1997) and Aida (1998); for a detailed numerical investigation of this problem see Nemenman and Bialek (2001). Our goal here is to be illustrative rather than exhaustive.<sup>11</sup><sup>11</sup>11We caution the reader that our discussion in this section is less self–contained than in other sections. Since the crucial steps exactly parallel those in the earlier work, here we just give references.
From the discussion above we know that the predictive information is related to the density of Kullback–Leibler divergences, and that the power–law behavior we are looking for comes from an essential singularity in this density function. Thus we calculate $`\rho (D,\overline{\varphi })`$ in the model defined by Eq. (85).
With $`Q(x)=(1/l_0)\mathrm{exp}[\varphi (x)]`$, we can write the KL divergence as
$$D_{\mathrm{KL}}[\overline{\varphi }(x)||\varphi (x)]=\frac{1}{l_0}dx\mathrm{exp}[\overline{\varphi }(x)][\varphi (x)\overline{\varphi }(x)].$$
(86)
We want to compute the density,
$`\rho (D;\overline{\varphi })`$ $`=`$ $`{\displaystyle }[d\varphi (x)]𝒫[\varphi (x)]\delta (DD_{\mathrm{KL}}[\overline{\varphi }(x)||\varphi (x)])`$ (87)
$`=`$ $`M{\displaystyle }[d\varphi (x)]𝒫[\varphi (x)]\delta (MDMD_{\mathrm{KL}}[\overline{\varphi }(x)||\varphi (x)]),`$ (88)
where we introduce a factor $`M`$ which we will allow to become large so that we can focus our attention on the interesting limit $`D0`$. To compute this integral over all functions $`\varphi (x)`$, we introduce a Fourier representation for the delta function, and then rearrange the terms:
$`\rho (D;\overline{\varphi })`$ $`=`$ $`M{\displaystyle \frac{dz}{2\pi }\mathrm{exp}(izMD)[d\varphi (x)]𝒫[\varphi (x)]\mathrm{exp}(izMD_{\mathrm{KL}})}`$ (90)
$`=`$ $`M{\displaystyle \frac{dz}{2\pi }\mathrm{exp}\left(izMD+\frac{izM}{l_0}𝑑x\overline{\varphi }(x)\mathrm{exp}[\overline{\varphi }(x)]\right)}`$
$`\times {\displaystyle }[d\varphi (x)]𝒫[\varphi (x)]\mathrm{exp}({\displaystyle \frac{izM}{l_0}}{\displaystyle }dx\varphi (x)\mathrm{exp}[\overline{\varphi }(x)]).`$
The inner integral over the functions $`\varphi (x)`$ is exactly the integral which was evaluated in the original discussion of this problem (Bialek, Callan and Strong 1996); in the limit that $`zM`$ is large we can use a saddle point approximation, and standard field theoretic methods allow us to compute the fluctuations around the saddle point. The result is that
$$[d\varphi (x)]𝒫[\varphi (x)]\mathrm{exp}\left(\frac{izM}{l_0}𝑑x\varphi (x)\mathrm{exp}[\overline{\varphi }(x)]\right)$$
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{izM}{l_0}}{\displaystyle 𝑑x\overline{\varphi }(x)\mathrm{exp}[\overline{\varphi }(x)]}S_{\mathrm{eff}}[\overline{\varphi }(x);zM]\right),`$
$`S_{\mathrm{eff}}[\overline{\varphi };zM]`$ $`=`$ $`{\displaystyle \frac{l}{2}}{\displaystyle 𝑑x\left(\frac{\overline{\varphi }}{x}\right)^2}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{izM}{ll_0}}\right)^{1/2}{\displaystyle 𝑑x\mathrm{exp}[\overline{\varphi }(x)/2]}.`$
Now we can do the integral over $`z`$, again by a saddle point method. The two saddle point approximations are both valid in the limit that $`D0`$ and $`MD^{3/2}\mathrm{}`$; we are interested precisely in the first limit, and we are free to set $`M`$ as we wish, so this gives us a good approximation for $`\rho (D0;\overline{\varphi })`$. This results in
$`\rho (D0;\overline{\varphi })`$ $`=`$ $`A[\overline{\varphi }(x)]D^{3/2}\mathrm{exp}\left({\displaystyle \frac{B[\overline{\varphi }(x)]}{D}}\right),`$ (93)
$`A[\overline{\varphi }(x)]`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{16\pi ll_0}}}\mathrm{exp}\left[{\displaystyle \frac{l}{2}}{\displaystyle 𝑑x\left(_x\overline{\varphi }\right)^2}\right]{\displaystyle 𝑑x\mathrm{exp}[\overline{\varphi }(x)/2]}`$ (94)
$`B[\overline{\varphi }(x)]`$ $`=`$ $`{\displaystyle \frac{1}{16ll_0}}\left({\displaystyle 𝑑x\mathrm{exp}[\overline{\varphi }(x)/2]}\right)^2.`$ (95)
Except for the factor of $`D^{3/2}`$, this is exactly the sort of essential singularity that we considered in the previous section, with $`\mu =1`$. The $`D^{3/2}`$ prefactor does not change the leading large $`N`$ behavior of the predictive information, and we find that
$$S_1^{(\mathrm{a})}(N)\frac{1}{2\mathrm{ln}2\sqrt{ll_0}}𝑑x\mathrm{exp}[\overline{\varphi }(x)/2]_{\overline{\varphi }}N^{1/2},$$
(96)
where $`\mathrm{}_{\overline{\varphi }}`$ denotes an average over the target distributions $`\overline{\varphi }(x)`$ weighted once again by $`𝒫[\overline{\varphi }(x)]`$. Notice that if $`x`$ is unbounded then the average in Eq. (96) is infrared divergent; if instead we let the variable $`x`$ range from $`0`$ to $`L`$ then this average should be dominated by the uniform distribution. Replacing the average by its value at this point, we obtain the approximate result
$$S_1^{(\mathrm{a})}(N)\frac{1}{2\mathrm{ln}2}\sqrt{N}\left(\frac{L}{l}\right)^{1/2}\mathrm{bits}.$$
(97)
To understand the result in Eq. (97), we recall that this field theoretic approach is more or less equivalent to an adaptive binning procedure in which we divide the range of $`x`$ into bins of local size $`\sqrt{l/NQ(x)}`$ (Bialek, Callan, and Strong 1996). From this point of view, Eq. (97) makes perfect sense: the predictive information is directly proportional to the number of bins that can be put in the range of $`x`$. This also is in direct accord with a comment from the previous subsection that power law behavior of predictive information arises from the number of parameters in the problem depending on the number of samples.
This counting of parameters allows a schematic argument about the smallness of fluctuations in this particular nonparametric problem. If we take the hint that at every step $`\sqrt{N}`$ bins are being investigated, then we can imagine that the field theoretic prior in Eq. (85) has imposed a structure $`𝒞`$ on the set of all possible densities, so that the set $`C_\mathrm{n}`$ is formed of all continuous piecewise linear functions that have not more than $`\mathrm{n}`$ kinks. Learning such functions for finite $`\mathrm{n}`$ is a $`d_{\mathrm{VC}}=\mathrm{n}`$ problem. Now, as $`N`$ grows, the elements with higher capacities $`\mathrm{n}\sqrt{N}`$ are admitted. The fluctuations in such a problem are known to be controllable (Vapnik 1998), as discussed in more detail elsewhere (Nemenman 2000).
One thing which remains troubling is that the predictive information depends on the arbitrary parameter $`l`$ describing the scale of smoothness in the distribution. In the original work it was proposed that one should integrate over possible values of $`l`$ (Bialek, Callan and Strong 1996). Numerical simulations demonstrate that this parameter can be learned from the data itself (Nemenman and Bialek 2000), but perhaps even more interesting is a formulation of the problem by Periwal (1997, 1998) which recovers complete coordinate invariance by effectively allowing $`l`$ to vary with $`x`$. In this case the whole theory has no length scale, and there also is no need to confine the variable $`x`$ to a box (here of size $`L`$). We expect that this coordinate invariant approach will lead to a universal coefficient multiplying $`\sqrt{N}`$ in the analog of Eq. (97), but this remains to be shown.
In summary, the field theoretic approach to learning a smooth distribution in one dimension provides us with a concrete, calculable example of a learning problem with power–law growth of the predictive information. The scenario is exactly as suggested in the previous section, where the density of KL divergences develops an essential singularity. Heuristic considerations (Bialek, Callan, and Strong 1996; Aida 1999) suggest that different smoothness penalties and generalizations to higher dimensional problems will have sensible effects on the predictive information—all have power–law growth, smoother functions have smaller powers (less to learn), and higher dimensional problems have larger powers (more to learn)—but real calculations for these cases remain challenging.
## 5 $`I_{\mathrm{pred}}`$ as a measure of complexity
The problem of quantifying complexity is very old \[see Grassberger (1991) for a short review\]. Solomonoff (1964), Kolmogorov (1965), and Chaitin (1975) investigated a mathematically rigorous notion of complexity that measures (roughly) the minimum length of a computer program that simulates the observed time series \[see also Li and Vitányi (1993)\]. Unfortunately there is no algorithm that can calculate the Kolmogorov complexity of all data sets. In addition, algorithmic or Kolmogorov complexity is closely related to the Shannon entropy, which means that it measures something closer to our intuitive concept of randomness than to the intuitive concept of complexity \[as discussed, for example, by Bennett (1990)\]. These problems have fueled continued research along two different paths, representing two major motivations for defining complexity. First, one would like to make precise an impression that some systems—such as life on earth or a turbulent fluid flow—evolve toward a state of higher complexity, and one would like to be able to classify those states. Second, in choosing among different models that describe an experiment, one wants to quantify a preference for simpler explanations or, equivalently, provide a penalty for complex models that can be weighed against the more conventional goodness of fit criteria. We bring our readers up to date with some developments in both of these directions, and then discuss the role of predictive information as a measure of complexity. This also gives us an opportunity to discuss more carefully the relation of our results to previous work.
### 5.1 Complexity of statistical models
The construction of complexity penalties for model selection is a statistics problem. As far as we know, however, the first discussions of complexity in this context belong to philosophical literature. Even leaving aside the early contributions of William of Occam on the need for simplicity, Hume on the problem of induction, and Popper on falsifiability, Kemeney (1953) suggested explicitly that it would be possible to create a model selection criterion that balances goodness of fit versus complexity. Wallace and Burton (1968) hinted that this balance may result in the model with “the briefest recording of all attribute information.” Even though he probably had a somewhat different motivation, Akaike (1974a, 1974b) made the first quantitative step along these lines. His ad hoc complexity term was independent of the number of data points and was proportional to the number of free independent parameters in the model.
These ideas were rediscovered and developed systematically by Rissanen in a series of papers starting from 1978. He has emphasized strongly (Rissanen 1984, 1986, 1987, 1989) that fitting a model to data represents an encoding of those data, or predicting future data, and that in searching for an efficient code we need to measure not only the number of bits required to describe the deviations of the data from the model’s predictions (goodness of fit), but also the number of bits required to specify the parameters of the model (complexity). This specification has to be done to a precision supported by the data.<sup>12</sup><sup>12</sup>12Within this framework Akaike’s suggestion can be seen as coding the model to (suboptimal) fixed precision. Rissanen (1984) and Clarke and Barron (1990) in full generality were able to prove that the optimal encoding of a model requires a code with length asymptotically proportional to the number of independent parameters and logarithmically dependent on the number of data points we have observed. The minimal amount of space one needs to encode a data string (minimum description length or MDL) within a certain assumed model class was termed by Rissanen stochastic complexity, and in recent work he refers to the piece of the stochastic complexity required for coding the parameters as the model complexity (Rissanen 1996). This approach was further strengthened by a recent result (Vitányi and Li 2000) that an estimation of parameters using the MDL principle is equivalent to Bayesian parameter estimations with a “universal” prior (Li and Vitányi 1993).
There should be a close connection between Rissanen’s ideas of encoding the data stream and the subextensive entropy. We are accustomed to the idea that the average length of a code word for symbols drawn from a distribution $`P`$ is given by the entropy of that distribution; thus it is tempting to say that an encoding of a stream $`x_1,x_2,\mathrm{},x_N`$ will require an amount of space equal to the entropy of the joint distribution $`P(x_1,x_2,\mathrm{},x_N)`$. The situation here is a bit more subtle, because the usual proofs of equivalence between code length and entropy rely on notions of typicality and asymptotics as we try to encode sequences of many symbols; here we already have $`N`$ symbols and so it doesn’t really make sense to talk about a stream of streams. One can argue, however, that atypical sequences are not truly random within a considered distribution since their coding by the methods optimized for the distribution is not optimal. So atypical sequences are better considered as typical ones coming from a different distribution \[a point also made by Grassberger (1986)\]. This allows us to identify properties of an observed (long) string with the properties of the distribution it comes from, as was done by Vitányi and Li (2000). If we accept this identification of entropy with code length, then Rissanen’s stochastic complexity should be the entropy of the distribution $`P(x_1,x_2,\mathrm{},x_N)`$.
As emphasized by Balasubramanian (1997), the entropy of the joint distribution of $`N`$ points can be decomposed into pieces that represent noise or errors in the model’s local predictions—an extensive entropy—and the space required to encode the model itself, which must be the subextensive entropy. Since in the usual formulation all long–term predictions are associated with the continued validity of the model parameters, the dominant component of the subextensive entropy must be this parameter coding, or model complexity in Rissanen’s terminology. Thus the subextensive entropy should be the model complexity, and in simple cases where we can describe the data by a $`K`$–parameter model both quantities are equal to $`(K/2)\mathrm{log}_2N`$ bits to the leading order.
The fact that the subextensive entropy or predictive information agrees with Rissanen’s model complexity suggests that $`I_{\mathrm{pred}}`$ provides a reasonable measure of complexity in learning problems. This agreement might lead the reader to wonder if all we have done is to rewrite the results of Rissanen et al. in a different notation. To calm these fears we recall again that our approach distinguishes infinite VC problems from finite ones and treats nonparametric cases as well. Indeed, the predictive information is defined without reference to the idea that we are learning a model, and thus we can make a link to physical aspects of the problem, as discussed below.
The MDL principle was introduced as a procedure for statistical inference from a data stream to a model. In contrast, we take the predictive information to be a characterization of the data stream itself. To the extent that we can think of the data stream as arising from a model with unknown parameters, as in the examples of Section 4, all notions of inference are purely Bayesian and there is no additional “penalty for complexity.” In this sense our discussion is much closer in spirit to Balasubramanian (1997) than to Rissanen (1978). On the other hand, we can always think about the ensemble of data streams that can be generated by a class of models, provided that we have a proper Bayesian prior on the members of the class. Then the predictive information measures the complexity of the class, and this characterization can be used to understand why inference within some (simpler) model classes will be more efficient; for practical examples along these lines see Nemenman (2000) and Nemenman and Bialek (2001).
### 5.2 Complexity of dynamical systems
While there are a few attempts to quantify complexity in terms of deterministic predictions (Wolfram 1984), the majority of efforts to measure the complexity of physical systems starts with a probabilistic approach. In addition, there is a strong prejudice that the complexity of physical systems should be measured by quantities that are not only statistical, but are also at least related to more conventional thermodynamic quantities (temperature, entropy, $`\mathrm{}`$), since this is the only way one will be able to calculate complexity within the framework of statistical mechanics. Most proposals define complexity as an entropy–like quantity, but an entropy of some unusual ensemble. For example, Lloyd and Pagels (1988) identified complexity as thermodynamic depth, the entropy of the state sequences that lead to the current state. The idea clearly is in the same spirit as the measurement of predictive information, but this depth measure does not completely discard the extensive component of the entropy (Crutchfield and Shalizi 1999) and thus fails to resolve the essential difficulty in constructing complexity measures for physical systems: distinguishing genuine complexity from randomness (entropy), the complexity should be zero both for purely regular and for purely random systems.
New definitions of complexity that try to satisfy these criteria (Lopez–Ruiz et al. 1995, Gell–Mann and Lloyd 1996, Shiner et al. 1999, Sole and Luque 1999, Adami and Cerf 2000) and criticisms of these proposals (Crutchfield et al. 1999, Feldman and Crutchfield 1998, Sole and Luque 1999) continue to emerge even now. Aside from the obvious problems of not actually eliminating the extensive component for all or a part of the parameter space or not expressing complexity as an average over a physical ensemble, the critiques often are based on a clever argument first mentioned explicitly by Feldman and Crutchfield (1998). In an attempt to create a universal measure, the constructions can be made over–universal: many proposed complexity measures depend only on the entropy density $`𝒮_0`$ and thus are functions only of disorder—not a desired feature. In addition, many of these and other definitions are flawed because they fail to distinguish among the richness of classes beyond some very simple ones.
In a series of papers, Crutchfield and coworkers identified statistical complexity with the entropy of causal states, which in turn are defined as all those microstates (or histories) that have the same conditional distribution of futures (Crutchfield and Young 1989, Shalizi and Crutchfield 1999). The causal states provide an optimal description of a system’s dynamics in the sense that these states make as good a prediction as the histories themselves. Statistical complexity is very similar to predictive information, but Shalizi and Crutchfield (1999) define a quantity which is even closer to the spirit of our discussion: their excess entropy is exactly the mutual information between the semi–infinite past and future. Unfortunately, by focusing on cases in which the past and future are infinite but the excess entropy is finite, their discussion is limited to systems for which (in our language) $`I_{\mathrm{pred}}(T\mathrm{})=\mathrm{constant}`$.
In our view, Grassberger (1986, 1991) has made the clearest and the most appealing definitions. He emphasized that the slow approach of the entropy to its extensive limit is a sign of complexity, and has proposed several functions to analyze this slow approach. His effective measure complexity is the subextensive entropy term of an infinite data sample. Unlike Crutchfield et al., he allows this measure to grow to infinity. As an example, for low dimensional dynamical systems, the effective measure complexity is finite whether the system exhibits periodic or chaotic behavior, but at the bifurcation point that marks the onset of chaos, it diverges logarithmically. More interestingly, Grassberger also notes that simulations of specific cellular automaton models that are capable of universal computation indicate that these systems can exhibit an even stronger, power–law, divergence.
Grassberger (1986, 1991) also introduces the true measure complexity, or the forecasting complexity, which is the minimal information one needs to extract from the past in order to provide optimal prediction. Interestingly, another complexity measure, the logical depth (Bennett 1985), which measures the time needed to decode the optimal compression of the observed data, is bounded from below by this quantity because decoding requires reading all of this information. In addition, true measure complexity is exactly the statistical complexity of Crutchfield et al., and the two approaches are actually much closer than they seem. The relation between the true and the effective measure complexities, or between the statistical complexity and the excess entropy, closely parallels the idea of extracting or compressing relevant information (Tishby et al. 1999, Bialek and Tishby, in preparation), as discussed below.
### 5.3 A unique measure of complexity?
We recall that entropy provides a measure of information that is unique in satisfying certain plausible constraints (Shannon 1948). It would be attractive if we could prove a similar uniqueness theorem for the predictive information, or any part of it, as a measure of the complexity or richness of a time dependent signal $`x(0<t<T)`$ drawn from a distribution $`P[x(t)]`$. Before proceeding along these lines we have to be more specific about what we mean by “complexity.” In most cases, including the learning problems discussed above, it is clear that we want to measure complexity of the dynamics underlying the signal, or equivalently the complexity of a model that might be used to describe the signal.<sup>13</sup><sup>13</sup>13The problem of finding this model or of reconstructing the underlying dynamics may also be complex in the computational sense, so that there may not exist an efficient algorithm. More interestingly, the computational effort required may grow with the duration $`T`$ of our observations. We leave these algorithmic issues aside for the present discussion. There remains a question, however, whether we want to attach measures of complexity to a particular signal $`x(t)`$, or whether we are interested in measures (like the entropy itself) that are defined by an average over the ensemble $`P[x(t)]`$.
One problem in assigning complexity to single realizations is that there can be atypical data streams. Either we must treat atypicality explicitly, arguing that atypical data streams from one source should be viewed as typical streams from another source, as discussed by Vitányi and Li (2000), or we have to look at average quantities. Grassberger (1986) in particular has argued that our visual intuition about the complexity of spatial patterns is an ensemble concept, even if the ensemble is only implicit; see also Tong in the discussion session of Rissanen (1987). In Rissanen’s formulation of MDL, one tries to compute the description length of a single string with respect to some class of possible models for that string, but if these models are probabilistic we can always think about these models as generating an ensemble of possible strings. The fact that we admit probabilistic models is crucial: even at a colloquial level, if we allow for probabilistic models then there is a simple description for a sequence of truly random bits, but if we insist on a deterministic model then it may be very complicated to generate precisely the observed string of bits.<sup>14</sup><sup>14</sup>14This is the statement that the Kolmogorov complexity of a random sequence is large: the programs or algorithms considered in the Kolmogorov formulation are deterministic, and the program must generate precisely the observed string. Furthermore, in the context of probabilistic models it hardly makes sense to ask for a dynamics that generates a particular data stream; we must ask for dynamics that generate the data with reasonable probability, which is more or less equivalent to asking that the given string be a typical member of the ensemble generated by the model. All of these paths lead us to thinking not about single strings but about ensembles in the tradition of statistical mechanics, and so we shall search for measures of complexity that are averages over the distribution $`P[x(t)]`$.
Once we focus on average quantities, we can start by adopting Shannon’s postulates as constraints on a measure of complexity: if there are $`N`$ equally likely signals, then the measure should be monotonic in $`N`$; if the signal is decomposable into statistically independent parts then the measure should be additive with respect to this decomposition; and if the signal can be described as a leaf on a tree of statistically independent decisions then the measure should be a weighted sum of the measures at each branching point. We believe that these constraints are as plausible for complexity measures as for information measures, and it is well known from Shannon’s original work that this set of constraints leaves the entropy as the only possibility. Since we are discussing a time dependent signal, this entropy depends on the duration of our sample, $`S(T)`$. We know of course that this cannot be the end of the discussion, because we need to distinguish between randomness (entropy) and complexity. The path to this distinction is to introduce other constraints on our measure.
First we notice that if the signal $`x`$ is continuous, then the entropy is not invariant under transformations of $`x`$. It seems reasonable to ask that complexity be a function of the process we are observing and not of the coordinate system in which we choose to record our observations. The examples above show us, however, that it is not the whole function $`S(T)`$ which depends on the coordinate system for $`x`$;<sup>15</sup><sup>15</sup>15Here we consider instantaneous transformations of $`x`$, not filtering or other transformations that mix points at different times. it is only the extensive component of the entropy that has this noninvariance. This can be seen more generally by noting that subextensive terms in the entropy contribute to the mutual information among different segments of the data stream (including the predictive information defined here), while the extensive entropy cannot; mutual information is coordinate invariant, so all of the noninvariance must reside in the extensive term. Thus, any measure complexity that is coordinate invariant must discard the extensive component of the entropy.
The fact that extensive entropy cannot contribute to complexity is discussed widely in the physics literature (Bennett 1990), as our short review above shows. To statisticians and computer scientists, who are used to Kolmogorov’s ideas, this is less obvious. However, Rissanen (1986, 1987) also talks about “noise” and “useful information” in a data sequence, which is similar to splitting entropy into its extensive and the subextensive parts. His “model complexity,” aside from not being an average as required above, is essentially equal to the subextensive entropy. Similarly, Whittle \[in the discussion of Rissanen (1987)\] talks about separating the predictive part of the data from the rest.
If we continue along these lines, we can think about the asymptotic expansion of the entropy at large $`T`$. The extensive term is the first term in this series, and we have seen that it must be discarded. What about the other terms? In the context of learning a parameterized model, most of the terms in this series depend in detail on our prior distribution in parameter space, which might seem odd for a measure of complexity. More generally, if we consider transformations of the data stream $`x(t)`$ that mix points within a temporal window of size $`\tau `$, then for $`T>>\tau `$ the entropy $`S(T)`$ may have subextensive terms which are constant, and these are not invariant under this class of transformations. On the other hand, if there are divergent subextensive terms, these are invariant under such temporally local transformations.<sup>16</sup><sup>16</sup>16Throughout this discussion we assume that the signal $`x`$ at one point in time is finite dimensional. There are subtleties if we allow $`x`$ to represent the configuration of a spatially infinite system. So if we insist that measures of complexity be invariant not only under instantaneous coordinate transformations, but also under temporally local transformations, then we can discard both the extensive and the finite subextensive terms in the entropy, leaving only the divergent subextensive terms as a possible measure of complexity.
An interesting example of these arguments is provided by the statistical mechanics of polymers. It is conventional to make models of polymers as random walks on a lattice, with various interactions or self avoidance constraints among different elements of the polymer chain. If we count the number $`𝒩`$ of walks with $`N`$ steps, we find that $`𝒩(N)AN^\gamma z^N`$ (de Gennes 1979). Now the entropy is the logarithm of the number of states, and so there is an extensive entropy $`𝒮_0=\mathrm{log}_2z`$, a constant subextensive entropy $`\mathrm{log}_2A`$, and a divergent subextensive term $`S_1(N)\gamma \mathrm{log}_2N`$. Of these three terms, only the divergent subextensive term (related to the critical exponent $`\gamma `$) is universal, that is independent of the detailed structure of the lattice. Thus, as in our general argument, it is only the divergent subextensive terms in the entropy that are invariant to changes in our description of the local, small scale dynamics.
We can recast the invariance arguments in a slightly different form using the relative entropy. We recall that entropy is defined cleanly only for discrete processes, and that in the continuum there are ambiguities. We would like to write the continuum generalization of the entropy of a process $`x(t)`$ distributed according to $`P[x(t)]`$ as
$`S_{\mathrm{cont}}={\displaystyle Dx(t)P[x(t)]\mathrm{log}_2P[x(t)]},`$ (98)
but this is not well defined because we are taking the logarithm of a dimensionful quantity. Shannon gave the solution to this problem: we use as a measure of information the relative entropy or KL divergence between the distribution $`P[x(t)]`$ and some reference distribution $`Q[x(t)]`$,
$`S_{\mathrm{rel}}={\displaystyle Dx(t)P[x(t)]\mathrm{log}_2\left(\frac{P[x(t)]}{Q[x(t)]}\right)},`$ (99)
which is invariant under changes of our coordinate system on the space of signals. The cost of this invariance is that we have introduced an arbitrary distribution $`Q[x(t)]`$, and so really we have a family of measures. We can find a unique complexity measure within this family by imposing invariance principles as above, but in this language we must make our measure invariant to different choices of the reference distribution $`Q[x(t)]`$.
The reference distribution $`Q[x(t)]`$ embodies our expectations for the signal $`x(t)`$; in particular, $`S_{\mathrm{rel}}`$ measures the extra space needed to encode signals drawn from the distribution $`P[x(t)]`$ if we use coding strategies that are optimized for $`Q[x(t)]`$. If $`x(t)`$ is a written text, two readers who expect different numbers of spelling errors will have different $`Q`$s, but to the extent that spelling errors can be corrected by reference to the immediate neighboring letters we insist that any measure of complexity be invariant to these differences in $`Q`$. On the other hand, readers who differ in their expectations about the global subject of the text may well disagree about the richness of a newspaper article. This suggests that complexity is a component of the relative entropy that is invariant under some class of local translations and misspellings.
Suppose that we leave aside global expectations, and construct our reference distribution $`Q[x(t)]`$ by allowing only for short ranged interactions—certain letters tend to follow one another, letters form words, and so on, but we bound the range over which these rules are applied. Models of this class cannot embody the full structure of most interesting time series (including language), but in the present context we are not asking for this. On the contrary, we are looking for a measure that is invariant to differences in this short ranged structure. In the terminology of field theory or statistical mechanics, we are constructing our reference distribution $`Q[x(t)]`$ from local operators. Because we are considering a one dimensional signal (the one dimension being time), distributions constructed from local operators cannot have any phase transitions as a function of parameters; again it is important that the signal $`x`$ at one point in time is finite dimensional. The absence of critical points means that the entropy of these distributions (or their contribution to the relative entropy) consists of an extensive term (proportional to the time window $`T`$) plus a constant subextensive term, plus terms that vanish as $`T`$ becomes large. Thus, if we choose different reference distributions within the class constructible from local operators, we can change the extensive component of the relative entropy, and we can change constant subextensive terms, but the divergent subextensive terms are invariant.
To summarize, the usual constraints on information measures in the continuum produce a family of allowable complexity measures, the relative entropy to an arbitrary reference distribution. If we insist that all observers who choose reference distributions constructed from local operators arrive at the same measure of complexity, or if we follow the first line of arguments presented above, then this measure must be the divergent subextensive component of the entropy or, equivalently, the predictive information. We have seen that this component is connected to learning in a straightforward way, quantifying the amount that can be learned about dynamics that generate the signal, and to measures of complexity that have arisen in statistics and in dynamical systems theory.
## 6 Discussion
We have presented predictive information as a characterization of various data streams. In the context of learning, predictive information is related directly to generalization. More generally, the structure or order in a time series or a sequence is related almost by definition to the fact that there is predictability along the sequence. The predictive information measures the amount of such structure, but doesn’t exhibit the structure in a concrete form. Having collected a data stream of duration $`T`$, what are the features of these data that carry the predictive information $`I_{\mathrm{pred}}(T)`$? From Equation (10) we know that most of what we have seen over the time $`T`$ must be irrelevant to the problem of prediction, so that the predictive information is a small fraction of the total information; can we separate these predictive bits from the vast amount of nonpredictive data?
The problem of separating predictive from nonpredictive information is a special case of the problem discussed recently (Tishby et al. 1999, Bialek and Tishby, in preparation): given some data $`x`$, how do we compress our description of $`x`$ while preserving as much information as possible about some other variable $`y`$? Here we identify $`x=x_{\mathrm{past}}`$ as the past data and $`y=x_{\mathrm{future}}`$ as the future. When we compress $`x_{\mathrm{past}}`$ into some reduced description $`\widehat{x}_{\mathrm{past}}`$ we keep a certain amount of information about the past, $`I(\widehat{x}_{\mathrm{past}};x_{\mathrm{past}})`$, and we also preserve a certain amount of information about the future, $`I(\widehat{x}_{\mathrm{past}};x_{\mathrm{future}})`$. There is no single correct compression $`x_{\mathrm{past}}\widehat{x}_{\mathrm{past}}`$; instead there is a one parameter family of strategies which trace out an optimal curve in the plane defined by these two mutual informations, $`I(\widehat{x}_{\mathrm{past}};x_{\mathrm{future}})`$ vs. $`I(\widehat{x}_{\mathrm{past}};x_{\mathrm{past}})`$.
The predictive information preserved by compression must be less than the total, so that $`I(\widehat{x}_{\mathrm{past}};x_{\mathrm{future}})I_{\mathrm{pred}}(T)`$. Generically no compression can preserve all of the predictive information so that the inequality will be strict, but there are interesting special cases where equality can be achieved. If prediction proceeds by learning a model with a finite number of parameters, we might have a regression formula that specifies the best estimate of the parameters given the past data; using the regression formula compresses the data but preserves all of the predictive power. In cases like this (more generally, if there exist sufficient statistics for the prediction problem) we can ask for the minimal set of $`\widehat{x}_{\mathrm{past}}`$ such that $`I(\widehat{x}_{\mathrm{past}};x_{\mathrm{future}})=I_{\mathrm{pred}}(T)`$. The entropy of this minimal $`\widehat{x}_{\mathrm{past}}`$ is the true measure complexity defined by Grassberger (1986) or the statistical complexity defined by Crutchfield and Young (1989) \[in the framework of the causal states theory a very similar comment was made recently by Shalizi and Crutchfield (2000)\].
In the context of statistical mechanics, long range correlations are characterized by computing the correlation functions of order parameters, which are coarse–grained functions of the system’s microscopic variables. When we know something about the nature of the order parameter (e. g., whether it is a vector or a scalar), then general principles allow a fairly complete classification and description of long range ordering and the nature of the critical points at which this order can appear or change. On the other hand, defining the order parameter itself remains something of an art. For a ferromagnet, the order parameter is obtained by local averaging of the microscopic spins, while for an antiferromagnet one must average the staggered magnetization to capture the fact that the ordering involves an alternation from site to site, and so on. Since the order parameter carries all the information that contributes to long range correlations in space and time, it might be possible to define order parameters more generally as those variables that provide the most efficient compression of the predictive information, and this should be especially interesting for complex or disordered systems where the nature of the order is not obvious intuitively; a first try in this direction was made by Bruder (1998). At critical points the predictive information will diverge with the size of the system, and the coefficients of these divergences should be related to the standard scaling dimensions of the order parameters, but the details of this connection need to be worked out.
If we compress or extract the predictive information from a time series we are in effect discovering “features” that capture the nature of the ordering in time. Learning itself can be seen as an example of this, where we discover the parameters of an underlying model by trying to compress the information that one sample of $`N`$ points provides about the next, and in this way we address directly the problem of generalization (Bialek and Tishby, in preparation). The fact that (as mentioned above) nonpredictive information is useless to the organism suggests that one crucial goal of neural information processing is to separate predictive information from the background. Perhaps rather than providing an efficient representation of the current state of the world—as suggested by Attneave (1954), Barlow (1959, 1961), and others (Atick 1992)—the nervous system provides an efficient representation of the predictive information.<sup>17</sup><sup>17</sup>17If, as seems likely, the stream of data reaching our senses has diverging predictive information then the space required to write down our description grows and grows as we observe the world for longer periods of time. In particular, if we can observe for a very long time then the amount that we know about the future will exceed, by an arbitrarily large factor, the amount that we know about the present. Thus representing the predictive information may require many more neurons than would be required to represent the current data. If we imagine that the goal of primary sensory cortex is to represent the current state of the sensory world, then it is difficult to understand why these cortices have so many more neurons than they have sensory inputs. In the extreme case, the region of primary visual cortex devoted to inputs from the fovea has nearly 30,000 neurons for each photoreceptor cell in the retina (Hawken and Parker 1991); although much remains to be learned about these cells, it is difficult to imagine that the activity of so many neurons constitutes an efficient representation of the current sensory inputs. But if we live in a world where the predictive information in the movies reaching our retina diverges, it is perfectly possible that an efficient representation of the predictive information available to us at one instant requires thousands of times more space than an efficient representation of the image currently falling on our retina. It should be possible to test this directly by studying the encoding of reasonably natural signals and asking if the information which neural responses provide about the future of the input is close to the limit set by the statistics of the input itself, given that the neuron only captures a certain number of bits about the past. Thus we might ask if, under natural stimulus conditions, a motion sensitive visual neuron captures features of the motion trajectory that allow for optimal prediction or extrapolation of that trajectory; by using information theoretic measures we both test the “efficient representation” hypothesis directly and avoid arbitrary assumptions about the metric for errors in prediction. For more complex signals such as communication sounds, even identifying the features that capture the predictive information is an interesting problem.
It is natural to ask if these ideas about predictive information could be used to analyze experiments on learning in animals or humans. We have emphasized the problem of learning probability distributions or probabilistic models rather than learning deterministic functions, associations or rules. It is known that the nervous system adapts to the statistics of its inputs, and this adaptation is evident in the responses of single neurons (Smirnakis et al. 1996, Brenner et al. 2000); these experiments provide a simple example of the system learning a parameterized distribution. When making saccadic eye movements, human subjects alter their distribution of reaction times in relation to the relative probabilities of different targets, as if they had learned an estimate of the relevant likelihood ratios (Carpenter and Williams 1995). Humans also can learn to discriminate almost optimally between random sequences (fair coin tosses) and sequences that are correlated or anticorrelated according to a Markov process; this learning can be accomplished from examples alone, with no other feedback (Lopes and Oden 1987). Acquisition of language may require learning the joint distribution of successive phonemes, syllables, or words, and there is direct evidence for learning of conditional probabilities from artificial sound sequences, both by infants and by adults (Saffran et al. 1996; 1999). These examples, which are not exhaustive, indicate that the nervous system can learn an appropriate probabilistic model,<sup>18</sup><sup>18</sup>18As emphasized above, many other learning problems, including learning a function from noisy examples, can be seen as the learning of a probabilistic model. Thus we expect that this description applies to a much wider range of biological learning tasks. and this offers the opportunity to analyze the dynamics of this learning using information theoretic methods: What is the entropy of $`N`$ successive reaction times following a switch to a new set of relative probabilities in the saccade experiment? How much information does a single reaction time provide about the relevant probabilities? Following the arguments above, such analysis could lead to a measurement of the universal learning curve $`\mathrm{\Lambda }(N)`$.
The learning curve $`\mathrm{\Lambda }(N)`$ exhibited by a human observer is limited by the predictive information in the time series of stimulus trials itself. Comparing $`\mathrm{\Lambda }(N)`$ to this limit defines an efficiency of learning in the spirit of the discussion by Barlow (1983). While it is known that the nervous system can make efficient use of available information in signal processing tasks \[cf. Chapter 4 of Rieke et al. (1997)\], and that it can represent this information efficiently in the spike trains of individual neurons \[cf. Chapter 3 of Rieke et al. (1997), as well as Berry, Warland and Meister (1997), Strong et al. (1998), and Reinagel and Reid (2000)\], it is not known whether the brain is an efficient learning machine in the analogous sense. Given our classification of learning tasks by their complexity, it would be natural to ask if the efficiency of learning were a critical function of task complexity: perhaps we can even identify a limit beyond which efficient learning fails, indicating a limit to the complexity of the internal model used by the brain during a class of learning tasks. We believe that our theoretical discussion here at least frames a clear question about the complexity of internal models, even if for the present we can only speculate about the outcome of such experiments.
An important result of our analysis is the characterization of time series or learning problems beyond the class of finitely parameterizable models, that is the class with power–law divergent predictive information. Qualitatively this class is more complex than any parametric model, no matter how many parameters there may be, because of the more rapid asymptotic growth of $`I_{\mathrm{pred}}(N)`$. On the other hand, with a finite number of observations $`N`$, the actual amount of predictive information in such a nonparametric problem may be smaller than in a model with a large but finite number of parameters. Specifically, if we have two models, one with $`I_{\mathrm{pred}}(N)AN^\nu `$ and one with $`K`$ parameters so that $`I_{\mathrm{pred}}(N)(K/2)\mathrm{log}_2N`$, the infinite parameter model has less predictive information for all $`N`$ smaller than some critical value
$$N_c\left[\frac{K}{2A\nu }\mathrm{log}_2\left(\frac{K}{2A}\right)\right]^{1/\nu }.$$
(100)
In the regime $`N<<N_c`$, it is possible to achieve more efficient prediction by trying to learn the (asymptotically) more complex model, as illustrated concretely in simulations of the density estimation problem (Nemenman and Bialek 2000). Even if there are a finite number of parameters—such as the finite number of synapses in a small volume of the brain—this number may be so large that we always have $`NN_c`$, so that it may be more effective to think of the many parameter model as approximating a continuous or nonparametric one.
It is tempting to suggest that the regime $`N<<N_c`$ is the relevant one for much of biology. If we consider, for example, 10 mm<sup>2</sup> of inferotemporal cortex devoted to object recognition (Logothetis and Sheinberg 1996), the number of synapses is $`K5\times 10^9`$. On the other hand, object recognition depends on foveation, and we move our eyes roughly three times per second throughout perhaps 10 years of waking life during which we master the art of object recognition. This limits us to at most $`N10^9`$ examples. Remembering that we must have $`\nu <1`$, even with large values of $`A`$ Eq. (100) suggests that we operate with $`N<N_c`$. One can make similar arguments about very different brains, such as the mushroom bodies in insects (Capaldi, Robinson and Fahrbach 1999). If this identification of biological learning with the regime $`N<<N_c`$ is correct, then the success of learning in animals must depend on strategies that implement sensible priors over the space of possible models.
There is one clear empirical hint that humans can make effective use of models that are beyond finite parameterization (in the sense that predictive information diverges as a power–law), and this comes from studies of language. Long ago, Shannon (1951) used the knowledge of native speakers to place bounds on the entropy of written English, and his strategy made explicit use of predictability. Shannon showed $`N`$–letter sequences to native speakers (readers), asked them to guess the next letter, and recorded how many guesses were required before they got the right answer. Thus each letter in the text is turned into a number, and the entropy of the distribution of these numbers is an upper bound on the conditional entropy $`\mathrm{}(N)`$ \[cf. Eq. (11)\]. Shannon himself thought that the convergence as $`N`$ becomes large was rather quick, and quoted an estimate of the extensive entropy per letter $`𝒮_0`$. Many years later, Hilberg (1990) reanalyzed Shannon’s data and found that the approach to extensivity in fact was very slow: certainly there is evidence for a large component $`S_1(N)N^{1/2}`$, and this may even dominate the extensive component for accessible $`N`$. Ebeling and Pöschel (1994; see also Pöschel, Ebeling, and Rosé 1995) studied the statistics of letter sequences in long texts (like Moby Dick) and found the same strong subextensive component. It would be attractive to repeat Shannon’s experiments with a design that emphasizes the detection of subextensive terms at large $`N`$.<sup>19</sup><sup>19</sup>19Associated with the slow approach to extensivity is a large mutual information between words or characters separated by long distances, and several groups have found that this mutual information declines as a power law. Cover and King (1978) criticize such observations by noting that this behavior is impossible in Markov chains of arbitrary order. While it is possible that existing mutual information data have not reached asymptotia, the criticism of Cover and King misses the possibility that language is not a Markov process. Of course it cannot be Markovian if it has a power–law divergence in the predictive information.
In summary, we believe that our analysis of predictive information solves the problem of measuring the complexity of time series. This analysis unifies ideas from learning theory, coding theory, dynamical systems, and statistical mechanics. In particular we have focused attention on a class of processes that are qualitatively more complex than those treated in conventional learning theory, and there are several reasons to think that this class includes many examples of relevance to biology and cognition.
### Acknowledgements
We thank V. Balasubramanian, A. Bell, S. Bruder, C. Callan, A. Fairhall, G. Garcia de Polavieja Embid, R. Koberle, A. Libchaber, A. Melikidze, A. Mikhailov, O. Motrunich, M. Opper, R. Rumiati, R. de Ruyter van Steveninck, N. Slonim, T. Spencer, S. Still, S. Strong, and A. Treves for many helpful discussions. We also thank M. Nemenman and R. Rubin for a help with the numerical simulations, and an anonymous referee for pointing out yet more opportunities to connect our work with the earlier literature. Our collaboration was aided in part by a grant from the US–Israel Binational Science Foundation to the Hebrew University of Jerusalem, and work at Princeton was supported in part by funds from NEC.
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# 1 Introduction
## 1 Introduction
Last year, with Bering we proposed a new method to digitize the sum of planar diagrams selected by ’t Hooft’s $`N_c\mathrm{}`$ limit of $`SU(N_c)`$ gauge theories . The proposal, based on the light-cone or infinite momentum frame description of the dynamics, involved discretization of both the $`p^+`$ carried by each line of the diagram and the propagation time $`\tau =ix^+`$, as in . But the main advantage of the new version was a coherent prescription for resolving most of the ambiguities due to $`p^+=0`$ divergences that typically plague the light-cone description.
We hope that our formalism will eventually allow an improved understanding of QCD in 4 dimensional space-time. But in this article, we merely wish to test the proposal in the context of the well-understood case of large $`N_c`$ gauge theories in two space-time dimensions, namely the ’t Hooft model . Our purpose is not to unearth new aspects of the model, but rather to see how its well known properties can be obtained from our new discretization.
The physical content of the ’t Hooft model boils down to an integral equation, essentially a Bethe-Salpeter equation , that determines the mass spectrum of $`q\overline{q}`$ mesons. The reason the limit $`N_c\mathrm{}`$ reduces to ladder diagrams (albeit with self-energy corrected quark propagators), is that the 2 dimensional gluon is not dynamical (there are no transverse polarizations). Thus, as with any axial gauge, the light-cone gauge $`A_{}=0`$ eliminates all gluon self-interactions, so $`A_+`$ can be integrated out inducing an instantaneous Coulomb potential. But the ’t Hooft limit $`N_c\mathrm{}`$ further eliminates all quark loops and all non-planar diagrams, leaving only the planar self energy corrections to the quark propagator, and the ladder bare gluon exchanges (Coulomb interaction) between quark anti-quark lines in the singlet $`q\overline{q}`$ channel. In light-cone parameters the Bethe-Salpeter equation summing these ladder $`q\overline{q}`$ diagrams simplifies to the single variable ’t Hooft integral equation .
$`^2\phi (x)=\left({\displaystyle \frac{1}{x}}+{\displaystyle \frac{1}{1x}}\right)\mu ^2\phi (x){\displaystyle \frac{g_s^2N_c}{2\pi }}P{\displaystyle _0^1}𝑑y{\displaystyle \frac{\phi (y)\phi (x)}{(yx)^2}},`$ (1.1)
where the integral is understood to be evaluated by the principal value prescription. The variable $`x`$ is the fraction carried by the quark of the total $`P^+`$ of the system (the anti-quark carries $`P^+`$ fraction $`1x`$). Also $``$ is the mass of the meson bound state and $`\phi `$ satisfies the boundary conditions, $`\phi (0)=\phi (1)=0`$.
Since the new formalism discretizes $`\tau ix^+=ka`$ in addition to $`p^+=lm`$, the corresponding simplifications lead to an equation that is not a straightforward discretization of this integral equation. In particular, the continuum limit can be taken in different ways depending on the ratio $`T_0=m/a`$ (which would be infinite for continuous $`\tau `$), and we want to explore to what extent these different continuum limits lead to the same physics. We shall find that some care must be taken with the setup of the discrete $`\tau `$ dynamics in order for this to be true. Indeed, a numerical study shows that the most simple-minded treatment leads to a ground state that becomes unstable at moderate ’t Hooft coupling even with relatively small $`P^+/mM`$ unless the ratio $`a/m=1/T_0`$ is tuned to be sufficiently small (perhaps infinitesimal for large $`M`$). If this feature were robust, it would cast doubt on any potential utility of the discretization of $`\tau `$.
To overcome this difficulty, we find it necessary to veto some of the “densest” discretized Feynman diagrams: a quark must be forbidden to emit 2 gluons at immediately successive time steps, with a similar veto on two successive absorptions. With this simple veto (which is prescribed locally in time), we shall show that the continuum limit reduces to the ’t Hooft model provided only that the total $`P^+`$ of the $`q\overline{q}`$ system is large compared to the discretization unit $`m`$. In particular it is not necessary that the ratio $`T_0=m/a`$ be large. Keeping $`T_0`$ finite in the continuum limit leads to the ’t Hooft equation with a non-trivial renormalization of the coupling. Because of this effect, it turns out that the effective (renormalized) coupling is small for both large and small bare coupling, reminiscent of strong/weak coupling duality. The strong coupling limit favors the densest diagrams, so vetoing some of the densest ones has a dramatic effect on the strong coupling behavior of the theory. This possibility was anticipated and discussed in in connection with the nature of the fishnet diagrams in higher dimensional space-time.
The rest of the paper is organized as follows. In Section 2 we set up the discretized ’t Hooft model. We analyze it using a single time-step transfer matrix in Section 3 and using a Bethe-Salpeter approach in Section 4. In Section 5 we discuss and implement the veto which allows a satisfactory continuum limit at fixed $`T_0`$. Discussion and concluding remarks are the subject of the final Section.
## 2 Discretized ’t Hooft Model
The Lagrange density for $`SU(N_c)`$ gauge fields coupled to quarks in the fundamental representation is given by
$`={\displaystyle \frac{1}{4}}\mathrm{Tr}F^{\mu \nu }F_{\mu \nu }+\overline{q}\left[i\gamma (igA)\mu _0\right]q,`$ (2.1)
where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu ig[A_\mu ,A_\nu ]`$. We remind the reader that the normalization of gauge fields appropriate for matrix fields and dictated by the gluon kinetic term differs by a factor $`1/\sqrt{2}`$ from the more standard one:
$$\frac{1}{4}\underset{a}{}F_a^{\mu \nu }F_{a\mu \nu }=\frac{1}{2}\mathrm{Tr}F_s^{\mu \nu }F_{s\mu \nu },$$
with $`F_s_a\frac{\lambda _a}{2}F_a`$. Thus $`A_s=A/\sqrt{2}`$, and we conclude that $`g=g_s/\sqrt{2}`$. In 2 space-time dimensions we choose the representation of $`\gamma `$ matrices for which the light-like components are
$`\gamma ^+=\sqrt{2}\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)\gamma ^{}=\sqrt{2}\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).`$ (2.6)
With this choice the field equation for the upper component of the quark spinor does not involve the “time” derivative and is an equation of constraint relating the upper component, $`q_1`$, to the lower component, $`q_2`$. Working in light-cone gauge ($`A_{}=A^+=0`$), we can eliminate the upper component in favor of the lower component yielding the light-cone gauge Lagrange density
$`=+{\displaystyle \frac{1}{2}}\mathrm{Tr}(_{}A_+)^2+i\psi ^{}\left[_+igA_++{\displaystyle \frac{\mu _0^2}{2_{}}}\right]\psi ,`$ (2.7)
where $`\psi =2^{1/4}q_2`$.
Our discretization of Feynman diagrams is based on the $`x^+`$ representation of each bare propagator
$`D(p^+,x^+)={\displaystyle \frac{dp^{}}{2\pi }\stackrel{~}{D}(p^+,p^{})e^{ix^+p^{}}}.`$ (2.8)
Performing the $`p^{}`$ integral gives the following Feynman rules for the continuum theory
$`D_\psi (p^+,x^+)`$ $`=`$ $`e^{ix^+\mu _0^2/2p^+}e^{\tau \mu _0^2/2p^+}`$
$`D_A(p^+,x^+)`$ $`=`$ $`i{\displaystyle \frac{\delta (x^+)}{p_{}^{+}{}_{}{}^{2}}}{\displaystyle \frac{\delta (\tau )}{p_{}^{+}{}_{}{}^{2}}}`$ (2.9)
$`V_{\psi ^{}\psi A}`$ $`=`$ $`igg,`$
where the arrows indicate the rules to use with imaginary time.
One way to digitize the ’t Hooft equation (1.1) is to put the variables $`x,y`$ on a grid, which amounts to discrete light-cone quantization , where one discretizes the amount of $`P^+`$ each line of the ladder diagram carries in quanta of $`m`$
$$p^+=lml=1,2,3,\mathrm{}.$$
One can then focus on a state of the system of interest (in our case a $`q\overline{q}`$ system) with total $`P^+=Mm`$. The continuum theory is recovered by taking the combined limits $`m0`$ and $`M\mathrm{}`$ while keeping $`P^+=Mm`$ fixed. Following , in addition to discretizing the $`p^+`$ of each particle, we also discretize imaginary light-cone time, $`\tau =ix^+=ka`$ ($`k=1,2,3,\mathrm{}`$). This discretization (which also serves as an ultraviolet cutoff) allows the continuum limit to be reached by keeping $`T_0m/a`$ fixed and taking both $`m,a0`$ and $`M\mathrm{}`$ simultaneously. Actually, since the physics of the discretized model depends only on the ratio $`m/a`$, the continuum limit is nothing but the large $`M`$ limit, where $`M`$ measures the total $`P^+`$ of the system state. The conventional continuous time DLCQ approach (see and references therein) corresponds to the special case $`T_0\mathrm{}`$.
Discretization of the quark propagator poses no difficulty. However, for the instantaneous interaction induced by integrating out $`A_+`$, we allow for ambiguities as in . The only constraint is that the discretized propagator become that of Eq. 2.9 in the continuum limit. This allows us to spread out the instantaneous interaction away from $`\tau =0`$ (see for further discussion). Thus the gauge propagator can be expressed as
$`D_A(Mm,ika)=f_k{\displaystyle \frac{T_0}{M^2}}\mathrm{where}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}f_k=1.`$ (2.10)
We require that these arbitrary parameters $`f_k`$ rapidly vanish with increasing $`k`$. Using this discretization, the Feynman rules for the discrete theory are summarized in Fig. 1.
For the purposes of this paper we shall not exploit the full generality of the set of $`\{f_k\}`$’s. We restrict attention to the simplest version where the spread out interaction propagates only one unit in light-cone time, this corresponds to setting $`f_1=1,f_{k>1}=0`$. The Feynman rules of Fig. 1 can be further simplified if we absorb the negative sign from the anti-quark propagator into the corresponding vertex factor. We define new parameters
$`\alpha e^{\mu _0^2/2T_0}\mathrm{and}\kappa \sqrt{{\displaystyle \frac{g^2N_c}{2\pi T_0}}}.`$ (2.11)
We also recall that in ’t Hooft’s large $`N_c`$ limit every additional pair of cubic vertices in the ladder sum corresponds to a completed color index loop, which produces a factor $`N_c`$. Thus we shall also absorb a factor of $`\sqrt{N_c}`$ into each vertex. Simply put, all terms in the ladder sum are only dependent on the ’t Hooft coupling $`g^2N_c`$. The simplified Feynman rules are presented in Fig. 2.
## 3 Single Time-Step Transfer Matrix
Using the Feynman rules of Fig. 2 we can now proceed to set up a transfer matrix which evolves a singlet $`q\overline{q}`$ system one step forward in $`x^+`$-time. Once the matrix has been determined as a function of the coupling, $`\kappa `$, solving the eigenvalue problem will yield the bound state energies as functions of coupling. Since the scalar particle which mediates the Coulomb interaction only lives one time-step, any state can have at most two intermediate scalars. Thus for the simplest systems with $`P^+/mM=3,4,5,6`$ the number of states are $`3,7,14,25`$ (the number of states is $`(M1)(M^22M+6)/6`$ for general $`M`$). For illustrative purposes we shall explicitly present the transfer matrix for $`M=4`$.
For $`M=4`$ there are 7 states namely:
$`\begin{array}{ccccc}|3,1=b_3^{}d_1^{}|0,& & |2,2=b_2^{}d_2^{}|0,& & |1,3=b_1^{}d_3^{}|0,\\ |2,1,1=b_2^{}a_1^{}d_1^{}|0,& & |1,2,1=b_1^{}a_2^{}d_1^{}|0,& & |1,1,2=b_1^{}a_1^{}d_2^{}|0,\\ & & |1,1,1,1=b_1^{}a_{1}^{}{}_{}{}^{2}d_1^{}|0,& & \end{array}`$ (3.4)
where $`b^{}`$, $`d^{}`$, and $`a^{}`$ are creation operators for the quark, anti-quark and intermediate gauge particle states (the subscript on these operators denotes $`p^+/m`$).
By construction each of the quark and anti-quark states has at least one unit of $`p^+/m`$. The matrix that evolves the system forward in $`x^+`$ can be factored into a matrix $`A`$ that involves only propagators and a matrix $`V`$ that involves vertices. Writing the state of the system as a column vector, $`\mathrm{{\rm Y}}`$, with 7 components corresponding to the seven states in Eq. 3.4, the transfer matrix equation is
$`t\mathrm{{\rm Y}}=AV\mathrm{{\rm Y}},`$ (3.5)
where
$`A`$ $`=`$ $`\mathrm{diag}[\alpha ^{4/3},\alpha ,\alpha ^{4/3},\alpha ^{3/2}\alpha ^2/4,\alpha ^{3/2},\alpha ^2],`$ (3.6)
$`V`$ $`=`$ $`\left(\begin{array}{ccccccc}1& 0& 0& \kappa & \kappa & 0& 0\\ 0& 1& 0& \kappa & 0& \kappa & \kappa ^2\\ 0& 0& 1& 0& \kappa & \kappa & 0\\ \kappa & \kappa & 0& 0& 0& \kappa ^2& 0\\ \kappa & 0& \kappa & 0& 0& 0& 0\\ 0& \kappa & \kappa & \kappa ^2& 0& 0& 0\\ 0& \kappa ^2& 0& 0& 0& 0& 0\end{array}\right),`$ (3.14)
and the eigenvalue is $`t=e^{aE}`$. Solving this eigenvalue problem will yield energy eigenvalues as a function of the coupling $`\kappa `$. Note that the matrix $`AV`$ is not hermitian, and because of the negative diagonal entries in $`A`$, the equivalent matrix $`\sqrt{A}V\sqrt{A}`$ is not hermitian either. Thus there will, in general be complex eigenvalues $`t`$. The best one can hope for is that the lowest lying energy eigenvalues (highest lying positive real part for $`t`$) are real. A satisfactory outcome for the continuum limit $`M\mathrm{}`$ would be that the ground state energy and all the energy values with real parts of order $`1/M`$ above the ground state energy are real. Then the complex eigenvalues would be strict lattice artifacts.
The existence of complex $`t`$ eigenvalues is already evident at $`M=4`$ as shown in Fig. 3, where we have chosen $`\alpha =0.5`$ which for definiteness we use in subsequent graphs unless otherwise indicated. The ground state (highest value) of $`t`$ stays real and positive for all coupling. However the next 2 excited states stay real only for $`\kappa <\kappa _c`$ when they collide with eigenvalues that have emerged from $`t=0`$ (infinite energy) after which the eigenvalues become complex conjugate pairs. The hope is that for increasing $`M`$ the number of lowest lying energy levels that remain real all the way to strong coupling should increase. For $`M=4`$ analysis shows that the lowest energy eigenvalue (that of the ground state) stays well-separated from the other states (real and complex) for all couplings, see Fig. 3.
We also see the eigenvalue solutions (again see Fig. 3) which are well behaved at weak coupling can merge with $`t=0`$ solutions (solutions which have $`t=0`$ at zero coupling correspond to infinite energy lattice artifacts) and become complex. Complex $`t`$ solutions are not physical as they correspond to complex energies. This behavior is generic for our discretization, but as we shall see later, when the problem has been set up correctly, we can separate the lowest lying states which survive the continuum limit from the lattice artifacts.
However, when one performs a similar analysis for the $`M=5`$ and $`M=6`$ systems the lowest eigenstate at weak coupling does not remain the ground state for all coupling. In both cases a complex solution at weaker coupling becomes real at larger coupling with a lower energy than the weak coupling ground state. Comparing this behavior for $`M=5`$ and $`M=6`$ suggests that for increasing $`M`$ this probably occurs at weaker coupling. Thus for large $`M`$ the weak coupling ground state might only be valid for extremely weak (perhaps only infinitesimal) coupling.
Conventional continuous time DLCQ corresponds in our discretization to $`\kappa ^20`$ since then $`T_0\mathrm{}`$. In order for our light-cone time discretization to be useful, the solution should work for all coupling (corresponding to all values of $`T_0`$). Here, in this single time-step analysis, we see that our most naive discretization does not satisfy this requirement. We shall have to modify the discretization in order to fix this.
Since the continuum limit requires $`M\mathrm{}`$ the single time-step analysis is also inefficient because the rank of the matrix to diagonalize is of $`𝒪(M^3)`$. However, as we shall show in the following section, writing the ladder equation in the form of a Bethe-Salpeter equation (exchange-to-exchange rather than single time-step) will reduce the complexity of the eigenvalue problem to a matrix of rank of $`𝒪(M)`$.
## 4 Bethe-Salpeter equation
A more efficient way to solve the discretized ’t Hooft model is by setting up a Bethe-Salpeter equation . Instead of a matrix equation that evolves the $`q\overline{q}`$ system one step forward in time, we can write down a system of equations (also a matrix equation) which evolves the system exchange to exchange. The simplification is that the intermediate state involves two (dressed) particles ($`M1`$ possible states for general $`M`$) rather than two, three, and four bare particles as in the case of the single time-step transfer matrix. The trade-off is that the equations become more complicated because of the dressed propagators.
In order to set up the Bethe-Salpeter equation it is necessary to work out the dressed quark propagator. In the context of this discretization the dressed quark propagator is just the sum of all possible iterated bubbles. There is no room for nested bubbles because $`f_{k>1}=0`$. While the bubbles extend only one time-step in $`x^+`$, we must still allow for all possible $`P^+`$ routings through each bubble.
The energy representation of the bare quark propagator carrying $`p^+/m=l`$ (without bubbles), obtained by multiplying by $`u^k`$ and summing over all $`k>0`$, is given by
$`D_q(l)={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left(u\alpha ^{1/l}\right)^k={\displaystyle \frac{u\alpha ^{1/l}}{1u\alpha ^{1/l}}},`$ (4.1)
where $`u1/t=e^{aE}`$. The contribution of a single bubble is
$`\kappa ^2\mathrm{\Sigma }_l\kappa ^2{\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{1}{r^2}}\alpha ^{1/(lr)}.`$ (4.2)
The full propagator is given by iterations of Eq. 4.1 and Eq. 4.2 as displayed in Fig. 5
$`D_q^{\mathrm{full}}(l)=D_q(l){\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left(u\kappa ^2\mathrm{\Sigma }_lD_q(l)\right)^k={\displaystyle \frac{t\alpha ^{1/l}}{t^2t\alpha ^{1/l}+\alpha ^{1/l}\kappa ^2\mathrm{\Sigma }_l}}.`$ (4.3)
The denominator of the full propagator can be factored in two roots so that
$`D_q^{\mathrm{full}}(l)={\displaystyle \frac{t\alpha ^{1/l}}{(tt_+)(tt_{})}},`$ (4.4)
where
$`t_\pm ={\displaystyle \frac{\alpha ^{1/l}}{2}}\left[1\pm \sqrt{14\alpha ^{1/l}\kappa ^2\mathrm{\Sigma }_l}\right].`$ (4.5)
We can now partial fraction the full propagator
$`D_q^{\mathrm{full}}(l)`$ $`=`$ $`{\displaystyle \frac{\alpha ^{1/l}}{(t_+t_{})}}\left[{\displaystyle \frac{t_+}{(tt_+)}}{\displaystyle \frac{t_{}}{(tt_{})}}\right]`$ (4.6)
$`=`$ $`{\displaystyle \frac{\alpha ^{1/l}}{(t_+t_{})}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}u^k\left(t_+^kt_{}^k\right).`$
Expressing the full quark propagator as the sum in Eq. 4.6 allows us to read off the time representation of the full quark propagator for discrete $`\tau =ka`$.
What we really need in order to set up the Bethe-Salpeter equation is a ‘propagator’ which propagates the $`q\overline{q}`$ system, including bubbles, between exchanges between the quark and anti-quark, see Fig. 6.
The ‘propagator’ which evolves the system forward between exchanges is then
$`𝒟_{q\overline{q}}(l)={\displaystyle \frac{\alpha ^{M/l(Ml)}}{(t_+t_{})(s_+s_{})}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}u^k\left(t_+^kt_{}^k\right)\left(s_+^ks_{}^k\right),`$ (4.7)
where $`s_\pm `$ are the roots for the anti-quark (obtained simply by replacing $`l`$ in Eq. 4.5 by $`Ml`$). With some manipulation this can be simplified to
$`𝒟_{q\overline{q}}(l)={\displaystyle \frac{u\alpha ^{M/l(Ml)}\left(1u^2\kappa ^4\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{}\right)}{\left(1u^2\kappa ^4\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{}\right)^2u\alpha ^{M/l(Ml)}\left(1u\kappa ^2\mathrm{\Sigma }_{Ml}^{}\right)\left(1u\kappa ^2\mathrm{\Sigma }_l^{}\right)}},`$ (4.8)
where for brevity, we have defined
$`\mathrm{\Sigma }_l^{}\alpha ^{1/(Ml)}\mathrm{\Sigma }_l.`$ (4.9)
We can now now set up the Bethe-Salpeter equations
$`\mathrm{\Psi }_{\text{}}(l)`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{Ml2}{}}}{\displaystyle \frac{\kappa ^2}{r^2}}𝒟_{\text{}}(l+r)\mathrm{\Psi }_{\text{}}(l+r)+{\displaystyle \underset{r=1}{\overset{Ml1}{}}}{\displaystyle \frac{\kappa ^2}{r^2}}𝒟_{\text{}}(l+r)\mathrm{\Psi }_{\text{}}(l+r)`$
$`\mathrm{\Psi }_{\text{}}(l)`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{\kappa ^2}{r^2}}𝒟_{\text{}}(lr)\mathrm{\Psi }_{\text{}}(lr)+{\displaystyle \underset{r=1}{\overset{l2}{}}}{\displaystyle \frac{\kappa ^2}{r^2}}𝒟_{\text{}}(lr)\mathrm{\Psi }_{\text{}}(lr).`$ (4.10)
$`\mathrm{\Psi }_{\text{}}`$, $`\mathrm{\Psi }_{\text{}}`$ label two-particle states where the last ladder rung propagated forward in time from left to right or right to left, respectively. The first equation is graphically portrayed in Fig. 7.
Since each of the quark/anti-quark propagators must carry a minimum of one unit of $`p^+/m`$ there are only $`2(M2)`$ possible states
$`\mathrm{\Psi }_{\text{}}(l):1lM2,\mathrm{\Psi }_{\text{}}(l):2lM1.`$ (4.11)
Eq. 4.10 is constructed by evolving the system from a state just after one exchange in the ladder sum to just after the next. The various $`𝒟`$’s in Eq. 4.10 correspond to the Feynman diagram contributions which are either parallelogram or trapezoidal sections which take a $`\mathrm{\Psi }_{\text{}}`$ or $`\mathrm{\Psi }_{\text{}}`$ to a $`\mathrm{\Psi }_{\text{}}`$ or $`\mathrm{\Psi }_{\text{}}`$. The parallelogram propagator sections are simply related to Eq. 4.8,
$`𝒟_{\text{}}(l)`$ $`=`$ $`𝒟_{q\overline{q}}(l)`$
$`𝒟_{\text{}}(l)`$ $`=`$ $`𝒟_{q\overline{q}}(Ml)=𝒟_{q\overline{q}}(l).`$ (4.12)
However, the trapezoidal segments must be independently determined
$`𝒟_{\text{}}(l)`$ $`=`$ $`\overline{𝒟}_{q\overline{q}}(l)`$
$`𝒟_{\text{}}(l)`$ $`=`$ $`\overline{𝒟}_{q\overline{q}}(Ml),`$ (4.13)
where
$`\overline{𝒟}_{q\overline{q}}(l)`$ $`=`$ $`{\displaystyle \frac{\alpha ^{M/l(Ml)}}{(t_+t_{})(s_+s_{})}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}u^{k+1}\left(t_+^kt_{}^k\right)\left(s_+^{k+2}s_{}^{k+2}\right)`$ (4.14)
$`=`$ $`{\displaystyle \frac{u\alpha ^{2/(Ml)}\left[u\alpha ^{M/l(Ml)}\left(1u\kappa ^2\mathrm{\Sigma }_{Ml}^{}\right)\kappa ^2\mathrm{\Sigma }_{Ml}^{}\left(1u^2\kappa ^4\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{}\right)\right]}{\left(1u^2\kappa ^4\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{}\right)^2u\alpha ^{M/l(Ml)}\left(1u\kappa ^2\mathrm{\Sigma }_{Ml}^{}\right)\left(1u\kappa ^2\mathrm{\Sigma }_l^{}\right)}}.`$
In order to solve the matrix equation in Eq. 4.10 we would like to write it in the form of an eigenvalue problem yielding $`t=1/u`$ as a function of $`\kappa ^2`$. This is slightly complicated since the propagator segments involve $`\mathrm{\Sigma }^{}`$’s which appear together with factors of $`\kappa ^2`$. By setting $`\chi =u\kappa ^2`$ we can manipulate the equation to isolate $`t`$ as the eigenvalue, with solutions $`t_n(\chi )`$. This is achieved by rescaling $`\mathrm{\Psi }_{\text{}}`$, $`\mathrm{\Psi }_{\text{}}`$ by the denominator factor common to all $`𝒟`$’s, yielding
$`\alpha _l\left(1\chi \mathrm{\Sigma }_l^{}\right)\left(1\chi \mathrm{\Sigma }_{Ml}^{}\right)\mathrm{\Psi }_{\text{}}^{}(l)+{\displaystyle \underset{r=1}{\overset{Ml1}{}}}{\displaystyle \frac{\chi }{r^2}}\alpha _{l+r}\alpha ^{2/(Mlr)}\left(1\chi \mathrm{\Sigma }_{Mlr}^{}\right)\mathrm{\Psi }_{\text{}}^{}(l+r)`$
$`=t[(1\chi ^2\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{})^2\mathrm{\Psi }_{\text{}}^{}(l){\displaystyle \underset{r=1}{\overset{Ml2}{}}}{\displaystyle \frac{\chi }{r^2}}\alpha _{l+r}(1\chi ^2\mathrm{\Sigma }_{l+r}^{}\mathrm{\Sigma }_{Mlr}^{})\mathrm{\Psi }_{\text{}}^{}(l+r)`$
$`+{\displaystyle \underset{r=1}{\overset{Ml1}{}}}{\displaystyle \frac{\chi ^2}{r^2}}\alpha ^{2/(Mlr)}(1\chi ^2\mathrm{\Sigma }_{l+r}^{}\mathrm{\Sigma }_{Mlr}^{})\mathrm{\Psi }_{\text{}}^{}(l+r)]`$
$`\alpha _l\left(1\chi \mathrm{\Sigma }_l^{}\right)\left(1\chi \mathrm{\Sigma }_{Ml}^{}\right)\mathrm{\Psi }_{\text{}}^{}(l)+{\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{\chi }{r^2}}\alpha _{lr}\alpha ^{2/(lr)}\left(1\chi \mathrm{\Sigma }_{lr}^{}\right)\mathrm{\Psi }_{\text{}}^{}(lr)`$
$`=t[(1\chi ^2\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{})^2\mathrm{\Psi }_{\text{}}^{}(l){\displaystyle \underset{r=1}{\overset{l2}{}}}{\displaystyle \frac{\chi }{r^2}}\alpha _{lr}(1\chi ^2\mathrm{\Sigma }_{lr}^{}\mathrm{\Sigma }_{Ml+r}^{})\mathrm{\Psi }_{\text{}}^{}(lr)`$
$`+{\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{\chi ^2}{r^2}}\alpha ^{2/(lr)}(1\chi ^2\mathrm{\Sigma }_{lr}^{}\mathrm{\Sigma }_{Ml+r}^{})\mathrm{\Psi }_{\text{}}^{}(lr)],`$ (4.15)
where $`\alpha _l\alpha ^{M/l(Ml)}`$.
This discretized equation has roughly twice the complexity of a straightforward discretization of the ’t Hooft equation. The reason is that a rung propagating forward from left to right can couple to subsequent evolutions forbidden to a rung from right to left (and vice versa). See Fig. 8, for the graphs responsible for this asymmetry.
This is the reason we had to introduce a two-component Bethe-Salpeter wave function. An immediate consequence is that at $`\kappa =0`$ each energy value is at least doubly degenerate, including the ground state. This feature is evident in Fig. 9 where the solutions of the BS equation are displayed for $`M=4,5`$.
All of the solutions seen in Fig. 3 are present, but in addition there are extra spurious solutions. For example, with $`M=4`$, there is a second curve emerging from the $`\kappa =0`$ ground state eigenvalue. For $`\kappa >0`$ this extra eigenvalue curve lies below (in $`t`$) and well separated from the true ground level curve for all coupling. Similarly, for other values of $`M`$ the Bethe-Salpeter method consistently reproduces all the solutions of the transfer matrix method, but it also adds spurious solutions due to the two-component nature of the wave function.
One way to avoid these unwanted solutions is to slightly modify the discretized Feynman rules so that the rung will attach to the same lines whichever way the exchanged gluon propagates. As seen in Fig. 8, the asymmetry stems from the possibility of consecutive gluon emissions (absorptions) on immediately successive time steps. If this possibility is disallowed, the basic exchange rung can be taken to be the sum of the two different exchanges as in Fig. 10. In addition to removing unwanted solutions this veto rule also leads to simpler equations, with a more transparent continuum limit. As we shall see in the next section, it also produces a more physical strong coupling behavior than our original discretization.
## 5 Bethe-Salpeter with Veto
The Bethe-Salpeter equation for the discretized ’t Hooft model, with the veto imposed as described at the end of the previous section, is
$`\mathrm{\Psi }(l)`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{Ml1}{}}}{\displaystyle \frac{\kappa ^2}{r^2}}u\alpha ^{1/l+1/(Mlr)}𝒟_{q\overline{q}}(l+r)\mathrm{\Psi }(l+r)`$ (5.1)
$`+{\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{\kappa ^2}{r^2}}u\alpha ^{1/(lr)+1/(Ml)}𝒟_{q\overline{q}}(lr)\mathrm{\Psi }(lr),`$
where $`𝒟_{q\overline{q}}`$ is defined in Eq. 4.8. After re-indexing both sums the equation can be written as
$`\mathrm{\Psi }(l)={\displaystyle \underset{r=l+1}{\overset{M1}{}}}{\displaystyle \frac{u\kappa ^2}{(lr)^2}}\alpha ^{1/l+1/(Mr)}𝒟_{q\overline{q}}(r)\mathrm{\Psi }(r)+{\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{u\kappa ^2}{(lr)^2}}\alpha ^{1/r+1/(Ml)}𝒟_{q\overline{q}}(r)\mathrm{\Psi }(r).`$ (5.2)
By imposing the veto we have reduced the rank of the eigenvalue problem from $`2(M2)`$ to $`M1`$.
The new discretized equation is much easier to analyze in the formal continuum limit $`M\mathrm{}`$ than the original. First define $`\mathrm{\Phi }(r)𝒟_{q\overline{q}}(r)\mathrm{\Psi }(r)`$, and rearrange Eq. 5.2 to read
$`K(l)\mathrm{\Phi }(l)`$ $``$ $`\left({\displaystyle \frac{1}{u𝒟_{q\overline{q}}(l)}}\kappa ^2(\mathrm{\Sigma }_l^{}+\mathrm{\Sigma }_{Ml}^{})\right)\mathrm{\Phi }(l)`$
$`=`$ $`{\displaystyle \underset{r=l+1}{\overset{M1}{}}}{\displaystyle \frac{\kappa ^2}{(lr)^2}}\alpha ^{1/l+1/(Mr)}(\mathrm{\Phi }(r)\mathrm{\Phi }(l))+{\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{\kappa ^2}{(lr)^2}}\alpha ^{1/r+1/(Ml)}(\mathrm{\Phi }(r)\mathrm{\Phi }(l)).`$
To formally examine the continuum limit we suppose that each discrete $`p^+`$ variable is large putting each $`lxM`$,<sup>§</sup><sup>§</sup>§Of course even for $`M`$ large the equation does contain terms where $`l`$ and $`Ml`$ are small (i.e. close to 1). In order for these contributions to not affect the solution to the continuum Bethe-Salpeter equation, the wavefunction must vanish at the endpoints. We shall see how this occurs when we evaluate the numerics later. and take $`M\mathrm{}`$ at fixed $`x`$. Then the right hand side of Eq. 5 is set up to go to $`1/M`$ times the r.h.s. of the continuum ’t Hooft equation:
$`{\displaystyle \underset{r=l+1}{\overset{M1}{}}}{\displaystyle \frac{\kappa ^2}{(lr)^2}}\alpha ^{1/l+1/(Mr)}(\mathrm{\Phi }(r)\mathrm{\Phi }(l))+{\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{\kappa ^2}{(lr)^2}}\alpha ^{1/r+1/(Ml)}(\mathrm{\Phi }(r)\mathrm{\Phi }(l))`$
$`{\displaystyle \frac{\kappa ^2}{M}}P{\displaystyle _0^1}𝑑y{\displaystyle \frac{\mathrm{\Phi }(y)\mathrm{\Phi }(x)}{(yx)^2}}.`$ (5.4)
Clearly, $`u`$ must be chosen so that the l.h.s. is also of order $`1/M`$. Next, it is easy to verify that $`\mathrm{\Sigma }_l^{}=\alpha ^{1/l+1/(Ml)}(\pi ^2/61/l+𝒪(\mathrm{ln}l/l^2))`$, so that the inverse propagator can be simplified, neglecting terms of order $`\mathrm{ln}M/M^2`$,
$`{\displaystyle \frac{1}{u_l𝒟_{q\overline{q}}(l)}}{\displaystyle \frac{1}{u_l^2}}\kappa ^4\left({\displaystyle \frac{\pi ^2}{6}}{\displaystyle \frac{M}{2l(Ml)}}\right)^2{\displaystyle \frac{1}{u_l}}{\displaystyle \frac{1u_l\kappa ^2(\pi ^2/6M/2l(Ml))}{1+u_l\kappa ^2(\pi ^2/6M/2l(Ml))}},`$ (5.5)
where we have defined $`u_l=u\alpha ^{M/l(Ml)}`$. The factor $`K`$ multiplying $`\mathrm{\Phi }`$ on the l.h.s. of Eq. 5 can now be simplified to
$`K(l)`$ $``$ $`\alpha ^{M/l(Ml)}[{\displaystyle \frac{1}{u_l^2}}{\displaystyle \frac{\kappa ^4\pi ^4}{36}}{\displaystyle \frac{1}{u_l}}{\displaystyle \frac{1u_l\kappa ^2(\pi ^2/6M/2l(Ml))}{1+u_l\kappa ^2(\pi ^2/6M/2l(Ml))}}{\displaystyle \frac{\kappa ^2\pi ^2}{3}}`$ (5.6)
$`+{\displaystyle \frac{\kappa ^4\pi ^2}{6}}{\displaystyle \frac{M}{l(Ml)}}+\kappa ^2{\displaystyle \frac{M}{l(Ml)}}]`$
$``$ $`\alpha ^{M/l(Ml)}[{\displaystyle \frac{1}{u_l^2}}{\displaystyle \frac{\kappa ^4\pi ^4}{36}}{\displaystyle \frac{1}{u_l}}{\displaystyle \frac{1u_l\kappa ^2\pi ^2/6}{1+u_l\kappa ^2\pi ^2/6}}{\displaystyle \frac{\kappa ^2\pi ^2}{3}}`$
$`\text{ }+{\displaystyle \frac{M}{l(Ml)}}[{\displaystyle \frac{\kappa ^4\pi ^2}{6}}+\kappa ^2{\displaystyle \frac{\kappa ^2}{(1+u_l\kappa ^2\pi ^2/6)^2}}]]`$
Now write $`u=u_0e^{a\mathrm{\Delta }}`$, where $`a\mathrm{\Delta }`$ will be determined to be of order $`1/M`$, so that $`u_l=u_0(1+a\mathrm{\Delta }+(M/l(Ml))\mathrm{ln}\alpha )`$ to order $`1/M`$. Then $`u_0`$ must satisfy
$`f(u_0){\displaystyle \frac{1}{u_0^2}}{\displaystyle \frac{\kappa ^4\pi ^4}{36}}{\displaystyle \frac{1}{u_0}}{\displaystyle \frac{1u_0\kappa ^2\pi ^2/6}{1+u_0\kappa ^2\pi ^2/6}}{\displaystyle \frac{\kappa ^2\pi ^2}{3}}=0.`$ (5.7)
Then the continuum limit reads
$`\left[a\mathrm{\Delta }+{\displaystyle \frac{1}{Mx(1x)}}\left\{\mathrm{ln}\alpha +{\displaystyle \frac{1}{u_0f^{}(u_0)}}\left[{\displaystyle \frac{\kappa ^4\pi ^2}{6}}+\kappa ^2{\displaystyle \frac{\kappa ^2}{(1+u_0\kappa ^2\pi ^2/6)^2}}\right]\right\}\right]\mathrm{\Phi }(x)`$
$`={\displaystyle \frac{\kappa ^2}{Mu_0f^{}(u_0)}}P{\displaystyle _0^1}𝑑y{\displaystyle \frac{\mathrm{\Phi }(y)\mathrm{\Phi }(x)}{(yx)^2}}.`$
The energy of the system is $`E=(\mathrm{ln}u_0)/a+\mathrm{\Delta }`$, but the divergent first term is simply a physically irrelevant $`M`$ independent constant, so it is consistent to identify $`P^{}=\mathrm{\Delta }`$. Then $`^2=2P^+P^{}=2Mm\mathrm{\Delta }=2MT_0a\mathrm{\Delta }`$. We also identify
$$\mu ^2=2T_0\left\{\mathrm{ln}\alpha +\frac{1}{u_0f^{}(u_0)}\left[\frac{\kappa ^4\pi ^2}{6}+\kappa ^2\frac{\kappa ^2}{(1+u_0\kappa ^2\pi ^2/6)^2}\right]\right\},$$
and we obtain the continuum ’t Hooft equation
$`\left[^2\mu ^2{\displaystyle \frac{1}{x(1x)}}\right]\mathrm{\Phi }(x)`$ $`=`$ $`{\displaystyle \frac{2T_0\kappa ^2}{u_0f^{}(u_0)}}P{\displaystyle _0^1}𝑑y{\displaystyle \frac{\mathrm{\Phi }(y)\mathrm{\Phi }(x)}{(yx)^2}}`$ (5.9)
$`=`$ $`{\displaystyle \frac{g_s^2N_c}{2\pi u_0f^{}(u_0)}}P{\displaystyle _0^1}𝑑y{\displaystyle \frac{\mathrm{\Phi }(y)\mathrm{\Phi }(x)}{(yx)^2}}.`$
Comparing with Eq. 1.1, we see that the only effect on the continuum limit of keeping $`T_0`$ finite is a finite renormalization of the gauge coupling $`g^2g^2/u_0f^{}(u_0)`$, and a coupling constant dependent shift in $`\mu ^2`$. Thus, the only requirement for identical continuum physics is that $`u_0f^{}(u_0)`$ be negative. Since $`\alpha `$ is a free parameter, we can access all positive values of $`\mu ^2`$ by tuning it.
Eq. 5.7 implicitly relates $`u_0`$ to $`\kappa `$ via a cubic equation. Instead of solving this equation, it is more illuminating to use it to relate $`u_0`$ to the combination $`\eta u_0\kappa ^2\pi ^2/6`$
$`u_0={\displaystyle \frac{(1\eta ^2)(1+\eta )}{1+\eta +2\eta ^2}},\kappa ^2={\displaystyle \frac{6\eta }{u_0\pi ^2}}={\displaystyle \frac{6\eta (1+\eta +2\eta ^2)}{(1\eta ^2)(1+\eta )\pi ^2}}.`$ (5.10)
We can also obtain the charge renormalization factor $`u_0f^{}(u_0)`$ in terms of $`\eta `$:
$`u_0f^{}(u_0)={\displaystyle \frac{(1+\eta +2\eta ^2)(1+\eta +7\eta ^2\eta ^3)}{(1\eta ^2)^2(1+\eta )^2}},`$ (5.11)
the effective coupling in the ’t Hooft equation
$`{\displaystyle \frac{g_{\mathrm{eff}}^2N_c}{\pi }}={\displaystyle \frac{2\kappa ^2T_0}{u_0f^{}(u_0)}}={\displaystyle \frac{12\eta (1\eta ^2)(1+\eta )T_0}{\pi ^2(1+\eta +7\eta ^2\eta ^3)}},`$ (5.12)
and the renormalized mass parameter
$`\mu ^2=\mu _0^2+{\displaystyle \frac{12\eta ^2(3+\eta ^2)T_0}{\pi ^2(1+\eta +7\eta ^2\eta ^3)}},`$ (5.13)
where we have used $`\alpha =e^{\mu _0^2/2T_0}`$.
As a check, note that the continuous time limit corresponds to $`T_0\mathrm{}`$ or $`\kappa ^20`$, whence $`u_01`$ and $`\eta 0`$. Then the effective coupling Eq. 5.12 goes to $`12T_0\eta /\pi ^2=2T_0\kappa ^2=g^2N_c/\pi =g_s^2N_c/2\pi `$ as it should. Next, with discrete time, we see that, in order to have real energy and $`\kappa `$ ($`u_0>0`$ and $`\kappa ^2>0`$), we must place the restriction $`0<\eta <1`$. Small $`\kappa `$ corresponds to small $`\eta `$, and large $`\kappa `$ corresponds to $`\eta `$ near unity. Interestingly, we note that the effective coupling in the ’t Hooft equation is small in both the small and large $`\kappa `$ regimes.
It is easy to understand the small effective coupling at large $`\kappa `$ in terms of our discrete time Feynman diagrams. With discrete time, $`\kappa ^2\mathrm{}`$ causes the diagrams with a maximal number of powers of $`\kappa ^2`$ per time step to dominate. For example the $`q\overline{q}`$ propagator $`𝒟_{q\overline{q}}`$ behaves in this limit as
$`𝒟_{q\overline{q}}(l){\displaystyle \frac{u\alpha ^{M/l(Ml)}}{1u^2\kappa ^4\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{}}}=u\alpha ^{M/l(Ml)}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(u\kappa ^2)^{2k}(\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{})^k,`$ (5.14)
so that the propagator for $`2k+1`$ time steps is $`\alpha ^{M/l(Ml)}(\kappa ^4\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{})^k(\kappa ^2\pi ^2/6)^{2k}`$ in the continuum limit. We see that away from the endpoints there is a factor of $`\kappa ^2\pi ^2/6`$ per time step in the continuum limit, which corresponds to each quark propagating exactly one time unit between interactions. Since this is the eigenvalue of the transfer matrix, we immediately infer the strong coupling value of $`u_0=6/\kappa ^2\pi ^2`$. Because of our veto, every exchange between quark lines occupies precisely two time steps and contributes only a single factor of $`\kappa ^2`$. Thus each exchange costs a relative factor of $`1/\kappa ^2`$ in the strong coupling limit, and this relative factor is proportional to the effective coupling in the ’t Hooft equation. More precisely, separating out the factor corresponding to the strong coupling propagation of the quark and anti-quark for two time steps, we have $`\kappa ^2=(\kappa ^2\pi ^2/6)^2(36/\kappa ^2\pi ^4)`$, so the effective coupling for a single exchange is $`36/\kappa ^2\pi ^4`$ for large $`\kappa `$, in accord with the $`\eta 1`$ limit of Eqs. 5.105.12.
Now we turn to a numerical analysis of our discretized dynamics in order to understand how the continuum limit is approached in practice. As with the no-veto case in section 4 we can write this equation as an eigenvalue problem by rescaling $`\mathrm{\Psi }`$ and isolating the eigenvalue $`t`$ as a function of $`\chi u\kappa ^2`$. The resulting eigenvalue problem to solve is
$`t\mathrm{\Phi }(l)`$ $`=`$ $`{\displaystyle \frac{\alpha _l}{\left(1\chi ^2\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{}\right)}}[{\displaystyle \frac{\left(1\chi \mathrm{\Sigma }_l^{}\right)\left(1\chi \mathrm{\Sigma }_{Ml}^{}\right)}{\left(1\chi ^2\mathrm{\Sigma }_l^{}\mathrm{\Sigma }_{Ml}^{}\right)}}\mathrm{\Phi }(l)+\chi {\displaystyle \underset{r=1}{\overset{l1}{}}}{\displaystyle \frac{\alpha ^{1/r+1/(Ml)}}{(lr)^2}}\mathrm{\Phi }(r)`$ (5.15)
$`+\chi {\displaystyle \underset{r=l+1}{\overset{M1}{}}}{\displaystyle \frac{\alpha ^{1/l+1/(Mr)}}{(lr)^2}}\mathrm{\Phi }(r)].`$
We use numerical procedures in MAPLE and MATLAB to find the eigenvalues $`t_n(\chi )`$ of the matrix on the right hand side of this equation as a function of $`\chi `$. The value of $`\kappa ^2`$ is different for each $`t_n`$ since $`\kappa ^2=\chi t_n`$. However by varying $`0\chi \mathrm{}`$ we can generate the real solutions, $`t_n`$, for all $`\kappa ^2`$. In order to solve for complex $`t_n`$’s we would need to vary $`\chi `$ in the complex plane rather than just over positive real numbers.
The problem of contamination of the lowest lying states by complex solutions has been solved by our veto prescription: The lowest lying state for $`M=6`$ for Eq. 5.15 remains intact for all coupling $`\kappa ^2`$, see Fig. 11, which should be compared against Fig. 4 where the lowest lying state was only the ground state for $`\kappa ^23`$. When we analyze Eq. 5.15 for increasing $`M`$ (see Fig. 11 for $`M=16,32,64`$) we see that the number of low lying states that remain uncrossed for all couplings increases with increasing $`M`$. We also see that the spacing between these states decreases as $`M`$ increases. Recall that the solutions in Fig. 11 have been generated for $`\alpha =0.5`$.
In order to compare our numerical results for large values of $`M`$ (hopefully close to the continuum limit) with the numerical results of ’t Hooft we solve the Bethe-Salpeter equation in Eq. 5.15 for $`\kappa ^2=0.5`$ and
$`\alpha =1.16433\alpha =1.04167\alpha =0.70930.`$ (5.16)
These three choices of $`\alpha `$ correspond to values of ’t Hooft parameter, $`\stackrel{~}{\mu }^2\pi \mu ^2/g_{\mathrm{eff}}^2N_c`$, taken to be $`0`$, $`1`$ and $`2.11^2`$ respectively. These values of $`\stackrel{~}{\mu }`$ were used in . Fixing $`\kappa ^2`$ is equivalent to fixing $`\eta `$, $`u_0`$ and $`g_{\mathrm{eff}}^2N_c/\pi `$, thus choosing a value for $`\stackrel{~}{\mu }`$ determines $`\alpha `$ in Eq. 5.13.
As we can see in Fig. 12 plots of the three lowest lying energy levels against $`1/M`$ show curves that become linear with increasing $`M`$. These results can be fitted to the functional form
$`aE=\mathrm{ln}(u_0)+{\displaystyle \frac{c_1}{M}}\mathrm{exp}\left({\displaystyle \frac{c_2}{M}}+{\displaystyle \frac{c_3}{M^2}}\right),`$ (5.17)
where $`c_2`$ and $`c_3`$ parameterize the departure from $`1/M`$ behaviour away from large $`M`$. We used the data of Fig. 12 in the range $`128M2048`$ to fit this equation. With the fitted value of $`c_1`$ we can calculate the mass square of the corresponding bound state. As discussed previously, the $`M`$ independent term in Eq. 5.17 is dropped in identifying $`P^{}`$. Since
$`^2=2P^+P^{}=2MT_0(aE\mathrm{ln}u_0),`$ (5.18)
we have, for $`\kappa ^2=0.5`$,
$`^2={\displaystyle \frac{2c_1}{0.22265}}+\mathrm{},`$ (5.19)
in units of $`g_{\mathrm{eff}}^2N_c/\pi `$. The results of the fits are tabulated in Table 1 against the results of ’t Hooft .
We see that for $`\stackrel{~}{\mu }=1`$ and $`2.11`$, the results of our discretization match quite well those of . However, for $`\stackrel{~}{\mu }=0`$ we increased the range of $`M`$ to 4096, which still yielded a poor match. What we did note was that even for these sizable values of $`M`$, convergence for $`\stackrel{~}{\mu }=0`$ is slow. When fitting the data for $`\stackrel{~}{\mu }=0`$ for the ground state to Eq. 5.17 we are trying to force it to fit a coefficient to a $`1/M`$ term which is not supposed to be there. It is more appropriate to use the form
$`aE=\mathrm{ln}(u_0)+{\displaystyle \frac{c_1}{M^\beta }}\mathrm{exp}\left({\displaystyle \frac{c_2}{M}}+{\displaystyle \frac{c_3}{M^2}}\right),`$ (5.20)
where the power $`\beta `$ of the leading behavior is fitted dynamically. We performed this refined fit to the three lowest lying states for $`\stackrel{~}{\mu }=0`$ which yielded the results assembled in Table 2.
These results provide numerical evidence that for $`\stackrel{~}{\mu }=0`$, the 1st and 2nd excited states do have a nonzero meson mass (i.e. the leading behavior is $`1/M`$). However, the leading behavior for the ground state decreases more rapidly than $`1/M`$ and is consistent with zero meson mass.
We next address the issue of slow convergence for $`\stackrel{~}{\mu }=0`$ by examining the form of the ground energy eigenvector for increasing values of $`M`$. It is well known that the solutions of Eq. 1.1 for $`\mu =0`$ do not vanish at the endpoints $`x=0,1`$; indeed the exact ground state is simply a constant. As we can see in Fig. 13, at finite large $`M`$ the ground state solution of our discretized equation is ever smaller at the endpoints, and the progression of shapes is toward a more square profile. But even for $`M=4096`$ the eigenvector has not yet converged to its limiting form. This should be compared with the solution for $`\stackrel{~}{\mu }=1`$ which rapidly approaches it’s limiting form (see r.h.s. of Fig. 13). We see that, for our discretized equation, the solution for the ground state decreases more rapidly near the endpoints ($`x=0`$ and $`x=1`$) as $`M`$ increases, consistently with the shape eventually approaching a square profile at $`M\mathrm{}`$. However, it is not hard to show that consistency of the continuum limit requires that the range in $`x`$ over which the fall-off occurs must decrease less rapidly than $`1/\sqrt{M}`$. This still allows an approach to a square profile but convergence is necessarily slower than one might have expected. In fact all solutions of the continuum ’t Hooft equation with $`\stackrel{~}{\mu }=0`$ have non-zero values at the end points. Thus we should expect slow convergence for all solutions of the $`\stackrel{~}{\mu }=0`$ equation because the discrete solution tends to vanish at the endpoints but the limiting form does not. This effect does not occur for $`\stackrel{~}{\mu }>0`$ because then the continuum solution vanishes at the endpoints, so a decent approximation to it can be achieved with relatively smaller $`M`$.
## 6 Discussion and Conclusion
In this paper we have explored the efficacy of the discretization of large $`N_c`$ QCD proposed in by applying it to the well-understood ’t Hooft model. For a smooth continuum limit over the whole range of bare coupling $`\kappa `$, we had to introduce a refinement of the discrete time gluon emission vertex. This amounted to insisting that after an emission, at least 2 time steps had to intervene before the next emission, with a similar restriction on consecutive absorptions. In contrast, an absorption is allowed to immediately follow an emission and vice versa. With this refinement in place we found that the continuum ’t Hooft equation describes the mass spectrum for all real $`\kappa `$. However, the parameters that occur in the equation are renormalized from their bare values, as summarized in Eqs. 5.105.125.13.
An amusing outcome of this renormalization phenomenon is that the effective coupling goes to zero in both the small and large $`\kappa `$ limits. Perhaps this feature is a version of weak/strong coupling duality, much celebrated in recent developments in string/M theory. However, we must concede that 2 dimensional QCD may be too trivial to expect anything other than the usual continuum theory to emerge from any continuum limit. Another caveat against attributing much significance to this “duality” phenomenon, is that the physics of the continuum limit really only depends on the ratio $`\mu ^2/N_cg^2`$. This is because one can always choose the effective coupling as the fundamental unit of energy. Then the theories at different coupling but with the same value of this ratio (0 for example) are physically identical: any differences in description can be removed by a change of units.
At any rate, we conclude that the discretization of can be meaningfully applied to QCD in 2 space-time dimensions, with some intriguing hints about the nature of weak/strong coupling duality. An obvious and important limitation of the 2 dimensional case, however, is that the gluon has no dynamical degrees of freedom. Thus there is no opportunity for the $`P^+`$ of the system to be shared amongst an infinite number of gluons. This must occur for the fishnet diagrams to be relevant, and is allowed in higher dimensional space-time. The next step is to study the three dimensional case, the simplest gauge theory where fishnet diagrams can be relevant.
Acknowledgements: We thank Klaus Bering for his helpful contributions in the early stages of this project.
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# Left ordered groups with no nonabelian free subgroups
## 1. Introduction
A group $`G`$ is left ordered if it has a total ordering $``$ such that $`xygxgy`$ whenever $`g,x,yG`$. For much information on left ordered groups, see the books . Of course we say that a group $`G`$ is *right* ordered if it has a total ordering $``$ such that $`xyxgyg`$ whenever $`g,x,yG`$. However using the involution $`gg^1`$ of $`G`$, it is easy to see that a group is right orderable if and only if it is left orderable.
Recall that a group is *locally indicable* if and only if every finitely generated subgroup $`1`$ has an infinite cyclic quotient. Every locally indicable group is left orderable \[18, theorem 7.3.1\], but the converse is not true, as has been shown by Bergman and Tararin . On the other hand Chiswell and Kropholler \[4, theorem A\] showed that a solvable-by-finite left ordered group is locally indicable; also Tararin \[21, theorem 3\] has proved that if $`AG1`$ are groups with $`G/A`$ finitely generated and solvable, $`A`$ abelian and $`G`$ left orderable, then $`G`$ has a quotient isomorphic to an infinite subgroup of $``$. Further results in this direction were obtained in . In it was proved that every left ordered elementary amenable group is locally indicable, and the question was raised of whether every left ordered amenable group is locally indicable. Let $`\mathrm{NF}`$ denote the class of groups which contain no nonabelian free subgroup. We shall consider the following stronger statement.
###### Conjecture 1.1.
A left ordered $`\mathrm{NF}`$-group is locally indicable.
We shall now describe the class of groups for which we shall prove Conjecture 1.1. Let $`𝒫`$ denote the group of piecewise linear orientation preserving self homeomorphisms of the unit interval $`[0,1]`$ with multiplication defined as composition of functions. Thus if $`f,g𝒫`$, then $`f`$ is differentiable at all but a finite number of points, and $`(fg)(x)=f(g(x))`$ for all $`x[0,1]`$. Also let $`\mathrm{NS}`$ denote the class of groups which have no nonabelian free subsemigroup. Now define $`𝒞`$ to be the smallest class of groups which contains $`\mathrm{NS}`$ and $`𝒫`$, and is closed under taking subgroups, homomorphic images, group extensions and directed unions. Clearly $`𝒞`$ contains all elementary amenable groups and in particular all solvable by finite groups, and it is not difficult to show that $`𝒞\mathrm{NF}`$ (see Corollary 4.8). Moreover $`𝒞`$ contains groups which are not elementary amenable, such as the ubiquitous Thompson’s group (see for more information on this topic). Presumably not every $`\mathrm{NF}`$-group lies in the class $`𝒞`$, though I know of no explicit example in the literature. We can now state
###### Theorem 1.2.
A $`𝒞`$-group is left orderable if and only if it is locally indicable.
Of course the result that $`G`$ is locally indicable (whether or not $`G𝒞`$) implies that $`G`$ is left orderable has already been noted above. For the reverse implication, we prove a stronger result Theorem 4.12 which states that if $`F1`$ is a finitely generated left orderable group and $`FG𝒞`$, then there exists a left-relatively convex subgroup $`HF`$ (see Section 2) such that $`HGG`$, and $`G/HG`$ has a self centralizing torsion free normal abelian subgroup $`A/HG`$ such that $`G/A`$ is torsion free abelian. Thus in the special case $`F=G`$ (so $`G`$ is finitely generated and $`1`$), we see that $`G`$ has a quotient isomorphic to $``$.
In Section 5 we shall use Theorem 1.2 to prove the following result about $`\mathrm{Homeo}_+(S^1)`$, the group of orientation preserving homeomorphisms of the circle.
###### Corollary 1.3.
Let $`G`$ be a finitely generated subgroup of $`\mathrm{Homeo}_+(S^1)`$ such that $`G𝒞`$. Then
1. If $`G`$ is finite, then $`G`$ is cyclic.
2. If $`G`$ is infinite, then there exists $`KHG`$ such that $`G/H`$ is cyclic and $`H/K`$.
This should be compared with \[8, theorem 1.1\], where by considering the smaller group of orientation preserving *$`C^{\mathrm{}}`$ diffeomorphisms* of $`S^1`$, similar but stronger results were obtained.
In Section 6 we shall use some of the techniques in this paper to show that certain free products with amalgamation are left orderable. For example, we shall show in Theorem 6.3 that the free product of a left orderable group and a torsion free nilpotent group with an amalgamated cyclic subgroup is left orderable. In the final section we shall briefly consider some examples of left ordered groups which are not locally indicable.
Part of this work was carried out while I was at the Sonderforschungsbereich in Münster. I would like to thank Wolfgang Lück for organizing my visit to Münster, and the Sonderforschungsbereich for financial support.
## 2. Notation, Terminology and Assumed Results
As usual $``$, $``$, $``$ and $``$ will denote the rational numbers, real numbers, integers and natural numbers $`\{1,2,\mathrm{}\}`$ respectively. We shall use the notation $`G^{}`$ for the commutator subgroup of the group $`G`$, and if $`g,xG`$ and $`XG`$, then $`x^g=gxg^1`$, $`X^g=gXg^1`$, $`\text{C}_G(X)=\{gGx^g=x`$ for all $`xX\}`$, $`\text{C}_G(x)=\text{C}_G(\{x\})`$, and $`X`$ denotes the subgroup generated by $`X`$. Also if $`HG`$, then $`\mathrm{core}_G(H)=_{gG}H^g`$, the largest normal subgroup of $`G`$ contained in $`H`$.
All mappings will be written on the left, in particular all group actions will have the group acting on the left of the set. If $`G`$ is acting on a set $`Y`$ and $`ZY`$, then $`\mathrm{Stab}_G(Z)`$ will always denote the pointwise stablizer of $`Z`$ in $`G`$: thus $`\mathrm{Stab}_G(Z)=\{gGgz=z`$ for all $`zZ\}`$, and we write $`\mathrm{Stab}_G(y)`$ for $`\mathrm{Stab}_G(\{y\})`$ when $`yY`$. Also if $`HG`$, then $`\mathrm{Fix}_Y(H)`$ is the fixed points of $`H`$, that is $`\{yYhy=y`$ for all $`hH\}`$, and when $`gG`$ we write $`\mathrm{Fix}_Y(g)`$ for $`\mathrm{Fix}_Y(\{g\})`$. Then obviously $`\mathrm{Fix}_Y(X)=\mathrm{Fix}_Y(X)`$ whenever $`XG`$.
A totally ordered set $`X`$ is a set with a binary relation $``$ such that for $`x,y,zX`$, either $`xy`$ or $`yx`$, $`xy`$ and $`yx`$ implies $`x=y`$, and $`xyz`$ implies $`xz`$. Given totally ordered sets $`X`$ and $`Y`$, the map $`\theta :XY`$ is said to be order preserving if $`x<y`$ implies $`\theta x<\theta y`$ whenever $`x,yX`$. We shall let $`\mathrm{Aut}(X)`$ denote the group of all order preserving permutations $`XX`$. Note that if $`\theta `$ is an order preserving bijection $`XY`$, then $`\theta ^1`$ is also order preserving and thus $`\mathrm{Aut}(X)`$ is indeed a group. Also if $`X`$ and $`X`$ is given the order induced by the natural order on $``$, then the elements of $`\mathrm{Aut}(X)`$ are homeomorphisms of $`X`$.
If $`(G,)`$ is a left ordered group and $`KG`$, then we say that $`K`$ is a convex subgroup of $`G`$ if $`gG`$, $`j,kK`$ and $`jgk`$ implies $`gK`$. In this case the left cosets of $`K`$ in $`G`$, which we denote by $`G/K`$, is naturally a totally ordered set under the definition $`gK<hK`$ if and only if $`g<h`$ and $`gKhK`$, for $`g,hG`$. Furthermore $`G`$ then acts as order preserving permutations on $`G/K`$ according to the rule $`g(hK)=ghK`$. We say that a subgroup of $`G`$ is a *left-relatively convex* subgroup \[14, p. 127\] if it is convex with respect to some left order on $`G`$. Conversely suppose $`G`$ acts faithfully as order preserving permutations on some totally ordered set $`X`$. Then, as described in \[18, theorem 7.1.2\], we can make $`G`$ into a left ordered group as follows. Well order $`X`$, and then for $`f,gG`$ with $`fg`$, we say that $`f<g`$ if and only if $`f(x)<g(x)`$ where $`x`$ is the least element of $`X`$ such that $`f(x)g(x)`$. Note that if $`yX`$ and $`Y`$ is the set of all elements less than $`y`$, then $`\mathrm{Stab}_G(Y)`$ is a convex subgroup of $`G`$ under this order and consequently $`\mathrm{Stab}_G(Y)`$ is a left-relatively convex subgroup. Therefore $`\mathrm{Stab}_G(Y_0)`$ is a left-relatively convex subgroup of $`G`$ for any subset $`Y_0`$ of $`X`$. We need the following basic results about left-relatively convex subgroups.
###### Lemma 2.1.
Let $`G`$ be a left ordered group, let $`H`$ be a normal convex subgroup of $`G`$, and let $``$ be a set of left-relatively convex subgroups of $`G`$. Then
1. $`_BB`$ is a left-relatively convex subgroup of $`G`$.
2. If $``$ is totally ordered by inclusion, then $`_BB`$ is a left-relatively convex subgroup of $`G`$.
3. If $`B/H`$ is a left-relatively convex subgroup of $`G/H`$, then $`B`$ is a left-relatively convex subgroup of $`G`$.
###### Proof.
For (i) see \[14, proposition 5.1.10\] or \[16, lemma 2.2(i)\]. For (ii) see \[14, proposition 5.1.7\] or \[16, lemma 2.2(ii)\]. Finally for (iii), see \[16, lemma 2.1\]. ∎
We define $`\mathrm{PLO}`$ to be the class of groups which act faithfully as piecewise linear orientation preserving self homeomorphisms of $`[0,1]`$, and $`\mathrm{PLT}`$ the class of groups which act faithfully as piecewise linear orientation preserving self homeomorphisms of $`[0,1]`$ which do not have a common fixed point in $`(0,1)`$. Thus $`\mathrm{PLT}\mathrm{PLO}`$, $`G\mathrm{PLO}`$ if and only if $`G`$ is isomorphic to a subgroup of $`𝒫`$, and $`G\mathrm{PLT}`$ if and only if $`G`$ acts faithfully as piecewise linear orientation preserving self homeomorphisms of $`[0,1]`$, and given $`ϵ>0`$ and $`x(0,1)`$, there exist $`f,gG`$ such that $`f(x)<ϵ`$ and $`g(x)>1ϵ`$.
Finally in this section, we need the following refinement of the well known fact that a countable left ordered group can be considered as a subgroup of $`\mathrm{Aut}()`$ (see for example, \[23, lemma 2.2\]).
###### Lemma 2.2.
Let $`G`$ be a countable left ordered group, and let $`H`$ be a convex subgroup of $`G`$ such that $`HG`$. Then there is an order preserving action of $`G`$ on $``$ with kernel $`\mathrm{core}_G(H)`$ such that $`\mathrm{Stab}_G(0)=H`$ and $`\mathrm{Stab}_G(v)G`$ for all $`v`$.
###### Proof.
This follows from \[16, lemmas 2.4 and 2.3\]. ∎
## 3. Extension Closed Classes of Groups
Very similar results to the next lemma have been proved before, see for example \[16, lemma 3.1\]. We shall prove a more general result, which will hopefully avoid the need for further similar results. If $`𝒟`$ is a class of groups which is closed under taking subgroups, then we shall define $`\overline{𝒟}`$ to be the smallest class of groups containing $`𝒟`$ which is closed under group extension and is closed under directed unions. For arbitrary classes of groups $`𝒳`$ and $`𝒴`$, we shall let $`H\mathrm{L}𝒳`$ mean that every finite subset of the group $`H`$ is contained in an $`𝒳`$-subgroup, $`H\mathrm{Q}𝒳`$ to mean that $`H`$ is isomorphic to a quotient group of an $`𝒳`$-group, and $`H𝒳𝒴`$ mean that $`H`$ has a normal $`𝒳`$-subgroup $`X`$ such that $`H/X𝒴`$. If $`𝒳`$ is subgroup closed, then $`H\mathrm{L}𝒳`$ if and only if every finitely generated subgroup of $`H`$ is an $`𝒳`$-group. For each ordinal $`\alpha `$, the class of groups $`𝒟_\alpha `$ is defined inductively by $`𝒟_0=\{1\}`$, $`𝒟_{\alpha +1}=(\mathrm{L}𝒟_\alpha )𝒟`$ and $`𝒟_\beta =_{\alpha <\beta }𝒟_\alpha `$ if $`\beta `$ is a limit ordinal. Setting $`𝒳=_{\alpha 0}𝒟_\alpha `$, we can state
###### Lemma 3.1.
1. Each $`𝒟_\alpha `$ and $`𝒳`$ is subgroup closed.
2. If $`𝒟`$ is quotient group closed, then each $`𝒟_\alpha `$ and $`𝒳`$ is also quotient group closed.
3. $`𝒳=\overline{𝒟}`$.
###### Proof.
(i) This is easily proved by induction on $`\alpha `$, using the fact that $`𝒟`$ is subgroup closed.
(ii) This is also easily proved by induction on $`\alpha `$, using the fact that $`𝒟`$ is quotient group closed.
(iii) Clearly $`𝒳\overline{𝒟}`$, $`𝒳𝒟`$, and $`𝒳`$ is closed under directed unions. Therefore we need to prove that $`𝒳`$ is extension closed.
We show by induction on $`\beta `$ that $`𝒟_\alpha 𝒟_\beta 𝒟_{\alpha +\beta }`$; the case $`\beta =0`$ being obvious. If $`\beta =\gamma +1`$ for some ordinal $`\gamma `$, then
$`𝒟_\alpha 𝒟_\beta `$ $`=𝒟_\alpha ((\mathrm{L}𝒟_\gamma )𝒟)(𝒟_\alpha (\mathrm{L}𝒟_\gamma ))𝒟(\mathrm{L}(𝒟_\alpha 𝒟_\gamma ))𝒟`$
$`(\mathrm{L}𝒟_{\alpha +\gamma })𝒟\text{(by induction)}`$
$`=𝒟_{\alpha +\beta }.`$
On the other hand if $`\beta `$ is a limit ordinal, then $`𝒟_\beta =_{\gamma <\beta }𝒟_\gamma `$ and
$`𝒟_\alpha 𝒟_\beta `$ $`=𝒟_\alpha \left({\displaystyle \underset{\gamma <\beta }{}}𝒟_\gamma \right)={\displaystyle \underset{\gamma <\beta }{}}𝒟_\alpha 𝒟_\gamma `$
$`{\displaystyle \underset{\gamma <\beta }{}}𝒟_{\alpha +\gamma }\text{(by induction)}`$
$`𝒟_{\alpha +\beta }`$
as required. ∎
For the rest of this paper, we let $`𝒟=\mathrm{NS}\mathrm{Q}(\mathrm{PLO})`$. Then clearly $`\overline{𝒟}𝒞`$ and $`𝒟`$ is closed under taking subgroups, quotient groups, consequently $`𝒞=\overline{𝒟}=𝒳`$.
## 4. Proof of the Main Theorem
The statement and proof of the next lemma is just a reformulation of \[20, assertion 2.1\].
###### Lemma 4.1.
Let $`W`$ be a nonempty closed subset of $``$ and let $`\alpha ,\beta ,z\mathrm{Aut}(W)`$. Suppose $`\mathrm{Fix}_W(z)=\mathrm{}`$ and $`z`$ commutes with $`\alpha `$ and $`\beta `$. If $`\mathrm{Fix}_W(\alpha )\mathrm{}\mathrm{Fix}_W(\beta )`$ and $`\mathrm{Fix}_W(\alpha )\mathrm{Fix}_W(\beta )=\mathrm{}`$, then $`\alpha ,\beta `$ contains a nonabelian free subgroup.
###### Proof.
Without loss of generality we may assume that $`0,1W`$, $`\beta (0)=0`$, and $`z(0)=1`$. Then $`\beta (1)=1`$. Suppose that $`\mathrm{Fix}_W(\alpha )\mathrm{}`$ yet $`\mathrm{Fix}_W(\alpha )\mathrm{Fix}_W(\beta )=\mathrm{}`$. We may write $`\mathrm{Fix}_W(\alpha )`$ and $`\mathrm{Fix}_W(\beta )`$ as a disjoint union of open intervals, which we shall call $`𝒜`$ and $``$ respectively. Then a finite number $`n`$ of these intervals will cover $`[0,1]`$; let these intervals be $`(a_0,a_1)`$, $`(b_1,b_2)`$, $`(a_2,a_3)`$, $`(b_3,b_4)`$, …, $`(b_{n2},b_{n1})`$, $`(a_{n1},a_n)`$ (so $`n`$ is an odd integer) where the $`(a_{2i},a_{2i+1})`$ are intervals in $`𝒜`$, the $`(b_{2i+1},b_{2i+2})`$ are intervals in $``$, and $`a_0<0<a_1a_2<a_3\mathrm{}a_{n1}<1<a_n`$, $`0<b_1<b_2b_3<b_4<\mathrm{}<b_{n1}<1`$. Note that $`za_0=a_{n1}`$ and $`za_1=a_n`$, and $`a_i\mathrm{Fix}_W(\alpha )`$ and $`b_i\mathrm{Fix}_W(\beta )`$ for all $`i`$. Also $`(b_i,a_i)W\mathrm{}`$ if $`i`$ is odd, and $`(a_i,b_i)W\mathrm{}`$ if $`i`$ is even. To see this let us consider the former case. We have $`b_i\mathrm{Fix}_W(\beta )`$, so certainly $`b_iW`$. Also $`b_i\mathrm{Fix}_W(\alpha )`$, so by replacing $`\alpha `$ with $`\alpha ^1`$ if necessary, we may assume that $`\alpha (b_i)>b_i`$. Then $`\alpha (b_i)W`$ and since $`\alpha (a_i)=a_i`$, it follows that $`\alpha (b_i)(b_i,a_i)`$. Similarly if $`i`$ is even, we can show that $`(a_i,b_i)W\mathrm{}`$. Now choose $`x_i(b_i,a_i)W`$ if $`i`$ is odd ($`1in2`$), and $`x_i(a_i,b_i)W`$ ($`2in1`$) if $`i`$ is even. Finally set $`x_0=z^1x_{n1}`$.
Set $`P_1=(x_0,x_1)(x_2,x_3)\mathrm{}(x_{n3},x_{n2})`$ and $`Q_1=(x_1,x_2)(x_3,x_4)\mathrm{}(x_{n2},x_{n1})`$, and for $`r`$ define
$`z^rP_1`$ $`=(z^rx_0,z^rx_1)(z^rx_2,z^rx_3)\mathrm{}(z^rx_{n3},z^rx_{n2})`$
$`z^rQ_1`$ $`=(z^rx_1,z^rx_2)(z^rx_3,z^rx_4)\mathrm{}(z^rx_{n2},z^rx_{n1}).`$
Now set $`P=_rz^rP_1`$ and $`Q=_rz^rQ_1`$. Observe that $`PQ=\mathrm{}`$. Indeed if $`yPQ`$, then by translating by $`z^r`$ for suitable $`r`$, we may assume that $`y(0,1)`$, and then the result is clear.
If $`i`$ is even, then $`(x_i,x_{i+1})(a_i,a_{i+1})`$. Since $`\mathrm{Fix}_W(\alpha )(a_i,a_{i+1})=\mathrm{}`$, we see that either $`\alpha (x)>x`$ for all $`x(a_i,a_{i+1})W`$, or $`\alpha (x)<x`$ for all $`x(a_i,a_{i+1})W`$; without loss of generality we may assume that $`\alpha (x)>x`$. Now $`a_i,a_{i+1}\mathrm{Fix}_W(\alpha )`$, $`a_i<x_i<b_i`$ and $`b_{i+1}<x_{i+1}<a_{i+1}`$, hence there exists a positive integer $`p_i`$ such that $`\alpha ^r(x_i,x_{i+1})(x_{i+1},a_{i+1})`$ and $`\alpha ^r(x_i,x_{i+1})(a_i,x_i)`$ for all $`r>p_i`$. Let $`p`$ be the maximum of the of the $`p_i`$ ($`0in3`$). Since $`(x_{i+1},a_{i+1})`$, $`(a_i,x_i)Q`$, we see that $`\alpha ^{pr}(x_i,x_{i+1})Q`$ for all $`i`$ and for all $`r0`$, and it follows that $`\alpha ^{pr}PQ`$ for all $`r0`$. Similarly there exists a positive integer $`q`$ such that $`\beta ^{qr}QP`$ for all $`r0`$. It now follows from Klein’s Table Tennis lemma \[5, p. 130\] that $`\alpha ^p,\beta ^q`$ is a free group. ∎
###### Lemma 4.2.
Let $`W`$ be a nonempty closed subset of $``$ and let $`G,Z\mathrm{Aut}(W)`$. Suppose $`G\mathrm{NF}`$, $`Z`$ centralizes $`G`$ and $`\mathrm{Fix}_W(Z)=\mathrm{}`$. Let $`H=\{gG\mathrm{Fix}_W(g)\mathrm{}\}`$. Then $`G^{}HG`$ and $`\mathrm{Fix}_W(F)\mathrm{}`$ for every finitely generated subgroup $`F`$ of $`H`$.
###### Proof.
Let $`\{f_1,\mathrm{},f_n\}`$ be a finite subset of $`H`$, and set $`F=f_1,\mathrm{},f_n`$. The result will follow if we can prove that $`\mathrm{Fix}_W(F)\mathrm{}`$, because then clearly $`HG`$, and $`G/H`$ is abelian by \[1, Theorem 2.1\]. Choose $`a_i\mathrm{Fix}_W(f_i)`$, and then select $`aW`$ such that $`a<a_i`$ for all $`i`$. Now choose $`zZ`$ such that $`za>a_i`$ for all $`i`$.
Suppose the sequence $`\{z^rar>0\}`$ is bounded above, and let $`L`$ be the least upper bound of the sequence. Then $`z^ra_i\mathrm{Fix}_W(f_i)`$ for all $`i`$ and $`L`$ is the least upper bound of the sequence $`\{z^ra_ir>0\}`$. We conclude that $`L\mathrm{Fix}_W(f_i)`$ for all $`i`$ and hence $`L\mathrm{Fix}_W(F)`$. Therefore we may assume that the sequence $`\{z^rar>0\}`$ is not bounded above, and similarly we may assume that the sequence $`\{z^rar<0\}`$ is not bounded below.
Therefore we may assume that $`\mathrm{Fix}_W(z)=\mathrm{}`$. It now follows from Lemma 4.1 that $`\mathrm{Fix}_W(f)\mathrm{}`$ for all $`fF`$. Since $`z`$ commutes with all elements of $`F`$, we see that every element of $`F`$ fixes a point in $`[a,za)`$ and we deduce that the set $`\{fafF\}`$ is bounded above by $`za`$. If $`M`$ is the supremum of the set $`\{fafF\}`$, then $`M\mathrm{Fix}_W(F)`$ and the result is proven. ∎
The statement and proof of the next lemma is just a reformulation of \[19, lemma 3.1\]
###### Lemma 4.3.
Let $`W`$ be a nonempty closed subset of $``$, let $`\alpha ,\beta \mathrm{Aut}(W)`$, let $`a\mathrm{Fix}_W(\alpha )`$, and let $`b\mathrm{Fix}_W(\beta )`$. Suppose that $`a<b`$ and $`\mathrm{Fix}_W(\alpha )(a,b]=\mathrm{}=\mathrm{Fix}_W(\beta )[a,b)`$. Then $`\alpha ,\beta `$ contains a nonabelian free subsemigroup.
###### Proof.
By replacing $`\alpha `$ and or $`\beta `$ with their inverses if necessary, we may assume that $`\alpha (b)<b`$ and $`\beta (a)>a`$. Set $`x=\beta (a)`$ and note that $`a<x<b`$. Since $`\beta (b)=b`$, $`\beta `$ has no fixed points on $`[a,b)`$, and $`\beta (a)>a`$, we see that there exists a positive integer $`n`$ such that $`\beta ^n[a,b)(x,b)`$. Similarly there exists a positive integer $`m`$ such that $`\alpha ^m(a,b](a,x)`$.
We now show that the subsemigroup generated by $`\alpha ^m`$ and $`\beta ^n`$ is free on those generators. Set $`\gamma =\alpha ^m`$ and $`\delta =\beta ^n`$. Suppose to the contrary that two nontrivial distinct finite products $`\pi ,\rho `$ of the form $`\mathrm{}\gamma ^{n_1}\delta ^{n_2}\gamma ^{n_3}\delta ^{n_4}\mathrm{}`$, where the $`n_i`$ are positive integers, yield the same element of $`\mathrm{Aut}(W)`$. By cancelling on the left, we may assume without loss of generality that $`\pi =\gamma \pi _1`$ and $`\rho =\delta \rho _1`$ or 1, where $`\pi _1,\rho _1`$ are also products of the form $`\mathrm{}\gamma ^{n_1}\delta ^{n_2}\gamma ^{n_3}\delta ^{n_4}\mathrm{}`$. Since $`\pi _1x,\rho _1x(a,b)`$, we see that $`\gamma \pi _1x(a,x)`$ and $`\delta \rho _1x`$ or $`1x[x,b)`$. Thus $`\pi x\rho x`$ and we have a contradiction. We deduce that the subsemigroup generated by $`\alpha ^m`$ and $`\beta ^n`$ is free on those generators and the result follows. ∎
The statement and proof of the next lemma is just a reformulation of \[19, lemma 3.6\]
###### Lemma 4.4.
Let $`W`$ be a nonempty closed subset of $``$, let $`n`$ be a positive integer, let $`\alpha _1,\mathrm{},\alpha _n\mathrm{Aut}(W)`$, and let $`G=\alpha _1,\mathrm{},\alpha _n`$. Suppose $`G\mathrm{NS}`$. If $`\mathrm{Fix}_W(\alpha _i)\mathrm{}`$ for all $`i`$, then $`\mathrm{Fix}_W(G)\mathrm{}`$.
###### Proof.
For each $`i`$, we may write $`\mathrm{Fix}_W(\alpha _i)`$ as a disjoint union of open intervals, say $`_jI_{ij}`$, where each $`I_{ij}`$ is an open interval. Then $`I_{ij}`$ for all $`i,j`$, because $`\mathrm{Fix}_W(\alpha _i)\mathrm{}`$. Suppose $`\mathrm{Fix}_W(G)=\mathrm{}`$. If $`I_{ij}I_{kl}`$ and $`(i,j)(k,l)`$ then $`ik`$, so we may choose $`i,j`$ such that $`I_{ij}`$ is not contained in any other open interval. Using the fact that $``$ is connected, we may now choose $`k,l`$ so that $`I_{ij}`$ has nonempty intersection with $`I_{kl}`$, and also does not contain $`I_{kl}`$. Clearly $`ik`$. Write $`I_{ij}I_{kl}=(a,b)`$, and assume without loss of generality that $`aI_{kl}`$ and $`bI_{ij}`$. The result now follows by applying Lemma 4.3 with $`\alpha =\alpha _i`$ and $`\beta =\alpha _k`$. ∎
###### Lemma 4.5.
Let $`G\mathrm{PLT}`$, and let $`A`$ and $`B`$ be finitely generated subgroups of $`G^{}`$. Then there exists $`gG`$ such that $`A^g`$ and $`B`$ centralize each other.
###### Proof.
We may view $`G`$ as a subgroup of the piecewise linear orientation preserving homeomorphisms of $`[0,1]`$ such that $`G`$ fixes no point in $`(0,1)`$. Define $`H=\{gG\text{there exists }ϵ>0\text{ such that }g(t)=t\text{ for all }t[0,ϵ][1ϵ,1]\}`$. Then obviously $`HG`$ and we see that $`HG^{}`$. Therefore if $`A`$ and $`B`$ are finitely generated subgroups of $`G^{}`$, there exists $`0<r<s<1`$ such that $`c`$ is the identity map outside $`(r,s)`$ for all $`cAB`$. Since $`G`$ does not fix any point in $`(0,1)`$, there exists $`gG`$ such that $`gr>s`$. Then $`gAg^1`$ fixes all points outside $`(gr,gs)`$ and the result follows. ∎
###### Lemma 4.6.
Let $`G\mathrm{PLO}`$ be a finitely generated group. Then there exists a series $`G^{}=G_0G_1\mathrm{}G_n=1`$ with the property that $`G_iG`$ and if $`A,B`$ are finitely generated subgroups of $`G_{i1}/G_i`$, then there exists $`gG/G_i`$ such that $`A^g`$ and $`B`$ centralize each other, for $`i=1,\mathrm{},n`$.
###### Proof.
Write $`G=g_1,\mathrm{},g_m`$ and consider $`G`$ as a group of orientation preserving homeomorphisms of $`[0,1]`$. Let $`W^{}`$ denote the complement $`[0,1]W`$ of a subset $`W`$ of $`[0,1]`$. Since $`\mathrm{Fix}_{[0,1]}(g_i)^{}`$ is a finite union of open intervals and $`\mathrm{Fix}_{[0,1]}(G)^{}=_i\mathrm{Fix}_{[0,1]}(g_i)^{}`$, we see that $`\mathrm{Fix}_{[0,1]}(G)^{}`$ is a finite union of open intervals. Therefore $`\mathrm{Fix}_{[0,1]}(G)^{}`$ is a finite union of disjoint open intervals, say
$$(a_1,b_1)(a_2,b_2)\mathrm{}(a_n,b_n)$$
where $`0a_1<b_1a_2<b_2a_3<\mathrm{}a_n<b_n1`$. Set $`H_i=\mathrm{Stab}_G([a_i,b_i])`$, $`G_0=G^{}`$, and $`G_i=G_0H_1\mathrm{}H_i`$ for $`i=1,\mathrm{},n`$. Then $`G^{}=G_0G_1\mathrm{}G_n=1`$ and $`G_iG`$ for all $`i`$. Since $`g[a_i,b_i]=[a_i,b_i]`$ for all $`gG`$, we see that $`G/H_i\mathrm{Aut}([a_i,b_i])`$.
Now suppose $`A/G_i`$ and $`B/G_i`$ are finitely generated subgroups of $`G_{i1}/G_i`$. Then $`AH_i/H_i`$ and $`BH_i/H_i`$ are finitely generated subgroups of $`G^{}H_i/H_i`$, so by Lemma 4.5 there exists $`gG`$ such that $`[A^g,B]`$ (the commutator of $`A^g`$ and $`B`$) is contained in $`H_i`$. Therefore $`[A^g,B]H_iG_{i1}=G_i`$ and the result follows. ∎
###### Lemma 4.7.
Let $`1GQ(\mathrm{PLO})`$ and suppose $`G`$ is finitely generated. Then there exists $`1<HG`$ such that if $`A,B`$ are finitely generated subgroups of $`H`$, then there exists $`gG`$ such that $`A^g`$ and $`B`$ centralize each other.
###### Proof.
Write $`G=P/K`$ where $`P\mathrm{PLO}`$ and $`KP`$. If $`G=Kp_1,\mathrm{},Kp_d`$, then replacing $`P`$ with $`p_1,\mathrm{},p_d`$ and $`K`$ with $`Kp_1,\mathrm{},p_d`$, we may assume that $`P`$ is finitely generated. By Lemma 4.6, there exists a series $`P^{}=P_0P_1\mathrm{}P_n=1`$ such that $`P_iP`$ for all $`i`$ with the property that if $`C/P_i`$ and $`D/P_i`$ are finitely generated subgroups of $`P_{i1}/P_i`$, then there exists $`pP`$ such that $`C^p/P_i`$ and $`D/P_i`$ centralize each other. If $`P^{}K`$, then we may take $`H=G`$ and we are finished. Otherwise we may let $`m`$ be the smallest integer such that $`P_mK`$ and set $`H=P_{m1}K/K`$, a nontrivial normal subgroup of $`G`$. If $`A`$ and $`B`$ are finitely generated subgroups of $`H`$, then there exist finitely generated subgroups $`C`$ and $`D`$ of $`P_{m1}`$ such that $`A=CK/K`$ and $`B=DK/K`$. Then we can find $`pP`$ such that $`C^pP_m/P_m`$ and $`DP_m/P_m`$ centralize each other. If $`g=pK`$, then $`A^g`$ and $`B`$ centralize each other as required. ∎
###### Corollary 4.8.
$`𝒞\mathrm{NF}`$.
###### Proof.
It is easy to see that $`\mathrm{NF}`$ is closed under taking subgroups, quotient groups, group extensions and directed unions. Since $`\mathrm{NS}\mathrm{NF}`$, it will now be sufficient to show that $`𝒫\mathrm{NF}`$. However if $`G`$ is the free group on two generators, $`1HG`$, $`xHH^{}`$ and $`yH^{}1`$, then there is no $`gG`$ such that $`x^g`$ and $`y`$ centralize each other. We deduce from Lemma 4.7 that $`G\mathrm{PLO}`$ and the proof is complete. ∎
###### Lemma 4.9.
Let $`W`$ be a nonempty closed subset of $``$ and let $`G\mathrm{Aut}(W)`$. Suppose $`1G\mathrm{Q}(\mathrm{PLO})`$ and that $`G`$ is finitely generated. Then there exists $`HG`$ such that $`H1`$ and $`\mathrm{Fix}_W(F)\mathrm{}`$ whenever $`F`$ is a finitely generated subgroup of $`H^{}`$.
###### Proof.
Using Lemma 4.7, we can find $`HG`$ such that $`H1`$ with the property that if $`A,B`$ are finitely generated subgroups of $`H`$, then there exists $`gG`$ such that $`A^g`$ and $`B`$ centralize each other. The result is obvious if $`\mathrm{Fix}_W(A)\mathrm{}`$ for all finitely generated subgroups $`A`$ of $`H`$, so we may assume that there is a finitely generated subgroup $`A`$ of $`H`$ such that $`\mathrm{Fix}_W(A)=\mathrm{}`$. Let $`B`$ be any finitely generated subgroup of $`H`$. Then there exists $`gG`$ such that $`A^g`$ and $`B`$ centralize each other. Since $`B\mathrm{NF}`$ by Corollary 4.8, we deduce from Lemma 4.2 that $`\mathrm{Fix}_W(E)\mathrm{}`$ for every finitely generated subgroup $`E`$ of $`B^{}`$. We conclude that $`\mathrm{Fix}_W(F)\mathrm{}`$ for every finitely generated subgroup $`F`$ of $`H^{}`$ as required. ∎
###### Lemma 4.10.
Let $`W`$ be a nonempty closed subset of $``$ and let $`HG\mathrm{Aut}(W)`$ with $`H`$ solvable. Suppose either $`G/H\mathrm{NS}`$, or $`G/H\mathrm{Q}(\mathrm{PLO})`$ and $`G/H`$ is finitely generated. Then either $`G^{\prime \prime }=1`$, or there exists $`FG`$ such that $`F1`$ and $`\mathrm{Fix}_W(E)\mathrm{}`$ whenever $`E`$ is a finitely generated subgroup of $`F`$. Furthermore in the former case there exists $`AG`$ such that $`A`$ and $`G/A`$ are torsion free abelian, and $`\text{C}_G(A)=A`$.
###### Proof.
First suppose $`G`$ has a nontrivial normal abelian subgroup $`A`$. Here we set $`F=\{aA\mathrm{Fix}_W(a)\mathrm{}\}`$. Since $`A\mathrm{NS}`$, we see from Lemma 4.4 that $`FA`$ and $`\mathrm{Fix}_W(E)\mathrm{}`$ for all finitely generated subgroups $`E`$ of $`F`$, so we may assume that $`F=1`$. Let $`C=\text{C}_G(A)`$ and let $`B=\{bC\mathrm{Fix}_W(b)\mathrm{}\}`$. Since $`G\mathrm{NF}`$ by Corollary 4.8, Lemma 4.2 shows that $`B`$ is a normal subgroup of $`G`$ and that $`\mathrm{Fix}_W(E)\mathrm{}`$ for all finitely generated subgroups $`E`$ of $`B`$. Therefore we may assume that $`B=1`$. Using \[16, lemma 4.1\], we see that $`C`$ and $`G/C`$ are abelian and it follows that $`G^{\prime \prime }=1`$. Thus we may assume that $`G`$ has no nontrivial normal abelian subgroup, so in particular $`H=1`$.
We now have two cases to consider, namely $`G\mathrm{NS}`$ and $`G\mathrm{Q}(\mathrm{PLO})`$ and is finitely generated. In the former case the result follows from Lemma 4.4 and \[16, lemma 4.1\], while in the latter case the result follows from Lemma 4.9. ∎
###### Lemma 4.11.
Let $`HGF`$ be groups such that $`F`$ is finitely generated. Assume that $`G/H`$ has a solvable normal subgroup $`K/H`$ such that $`G/K𝒟`$ and is finitely generated. Suppose $`F`$ is left orderable and there exists a left-relatively convex subgroup $`B`$ of $`F`$ such that $`HBF`$. Then there exists a left-relatively convex subgroup $`B_1`$ of $`F`$ such that $`B_1F`$, $`B_1GG`$, and $`G/B_1G`$ has a self centralizing torsion free abelian normal subgroup $`B_2/B_1G`$ such that $`G/B_2`$ is torsion free abelian.
###### Proof.
For each $`XF`$, let $`𝔠X`$ denote the smallest left-relatively convex subgroup of $`F`$ containing $`X`$, and let $`𝒮=\{IGHI`$ and $`𝔠IF\}`$. Then $`𝒮`$ is partially ordered by inclusion. Suppose $`𝒯`$ is a nonempty chain in $`𝒮`$. Then $`_{I𝒯}𝔠I`$ is a left-relatively convex subgroup of $`F`$ by Lemma 2.1, which is not the whole of $`F`$ because $`F`$ is finitely generated and $`𝔠IF`$ for all $`I𝒯`$, consequently $`𝒯`$ is bounded above by $`_{I𝒯}I`$. But $`H𝒮`$ because $`𝔠HBF`$, hence $`𝒮\mathrm{}`$ and we may apply Zorn’s lemma to deduce that $`𝒮`$ has a maximal element $`E`$ say. Set $`B_1=_{gG}(𝔠E)^g`$, which by Lemma 2.1 is a left-relatively convex subgroup of $`F`$, so using the maximality of $`E`$ we see that $`B_1G=E`$ (thus $`B_1=𝔠E`$). If $`G/E`$ has a self centralizing torsion free abelian normal subgroup $`B_2/E`$ such that $`G/B_2`$ is torsion free abelian, then we are finished so we assume that this is not the case.
By Lemma 2.2, there is an order preserving action of $`F`$ on $``$ with kernel $`\mathrm{core}_F(𝔠E)`$ such that $`\mathrm{Stab}_F(0)=𝔠E`$, and $`\mathrm{Stab}_F(v)F`$ for all $`v`$. Replacing $`F`$ with $`F/\mathrm{core}_F(𝔠E)`$ and using Lemma 2.1, we may assume that $`\mathrm{core}_F(𝔠E)=1`$.
Let $`W=\mathrm{Fix}_{}(E)`$. Then $`W`$ is a nonempty closed subset of $``$, and $`G/E`$ is naturally a subgroup of $`\mathrm{Aut}(W)`$. Using the hypotheses of the Lemma, there is a normal solvable subgroup $`K_1/E`$ of $`G/E`$ such that $`G/K_1𝒟`$ and is finitely generated. By Lemma 4.10, there is a nontrivial normal subgroup $`A/E`$ of $`G/E`$ such that $`\mathrm{Fix}_W(C)\mathrm{}`$ whenever $`C`$ is a finitely generated subgroup of $`A/E`$. Write $`A=_iA_i`$ where $`EA_1A_2\mathrm{}`$ and $`A_i/E`$ is finitely generated for all $`i`$ (if $`A/E`$ is finitely generated, we may choose $`A_i=A`$ for all $`i`$), and $`X_i=\mathrm{Fix}_{}(A_i)`$. Then $`X_i\mathrm{}`$ for all $`i`$, and $`\mathrm{Stab}_F(X_1)\mathrm{Stab}_F(X_2)\mathrm{}`$ is an ascending chain of left-relatively convex subgroups of $`F`$ with the property that $`A_i\mathrm{Stab}_F(X_i)F`$ for all $`i`$. Furthermore $`_i\mathrm{Stab}_F(X_i)`$ is a left-relatively convex subgroup by Lemma 2.1, which cannot be $`F`$ itself because $`F`$ is finitely generated. We deduce that $`𝔠AF`$ which contradicts the maximality of $`E`$ and finishes the proof. ∎
###### Theorem 4.12.
Let $`GF1`$ be groups such that $`G𝒞`$ and $`F`$ is finitely generated and left orderable. Then there exists a left-relatively convex subgroup $`B`$ of $`F`$ such that $`BF`$, $`BGG`$, and $`G/BG`$ has a self centralizing torsion free abelian normal subgroup $`A/BG`$ such that $`G/A`$ is torsion free abelian.
###### Proof.
We shall prove the result by transfinite induction on $`G`$, so by Lemma 3.1 choose the least ordinal $`\alpha `$ such that $`G𝒟_\alpha `$ and assume that the result is true whenever $`H𝒟_\beta `$ and $`\beta <\alpha `$. Now $`\alpha `$ cannot be a limit ordinal, and the result is clearly true if $`\alpha =0`$. Therefore we may assume that $`\alpha =\gamma +1`$ for some ordinal $`\gamma `$, and then there exists $`HG`$ such that $`G/H𝒟`$ and $`H\mathrm{L}𝒟_\gamma `$. Using Lemma 3.1, we may write $`H=_iH_i`$ where $`H_1H_2\mathrm{}H`$ and every subgroup of $`H_i`$ is in $`𝒟_\gamma `$ for all $`i`$. For each $`XF`$, let $`𝔠X`$ denote the smallest left-relatively convex subgroup of $`F`$ containing $`X`$.
First consider the case $`G/H`$ is finitely generated. We have an ascending chain of left-relatively convex subgroups $`𝔠(H_1^{\prime \prime })𝔠(H_2^{\prime \prime })\mathrm{}`$, so their union is also a left-relatively convex subgroup by Lemma 2.1 which contains $`H^{\prime \prime }`$. The inductive hypothesis shows that $`𝔠(H_i^{\prime \prime })F`$ for all $`i`$ and since $`F`$ is finitely generated, we deduce that $`𝔠(H^{\prime \prime })F`$. But $`G/H^{\prime \prime }`$ has the solvable normal subgroup $`H/H^{\prime \prime }`$ such that $`G/H𝒟`$, so the result follows from an application of Lemma 4.11.
Finally we need to consider the case $`G/H`$ is not finitely generated. Here we write $`G=_iG_i`$, where $`HG_iG`$ and $`G_i/H`$ is finitely generated for all $`i`$. We now have an ascending chain of left relatively convex subgroups $`𝔠(G_1^{\prime \prime })𝔠(G_2^{\prime \prime })\mathrm{}`$, so their union is also a left-relatively convex subgroup by Lemma 2.1 which contains $`G^{\prime \prime }`$. By the case $`G/H`$ is finitely generated considered in the previous paragraph, we know that $`𝔠(G_i^{\prime \prime })F`$ for all $`i`$ and since $`F`$ is finitely generated, we deduce that $`𝔠(G^{\prime \prime })F`$. Another application of Lemma 4.11 completes the proof. ∎
## 5. Groups of homeomorphisms of the circle
Proof of Corollary 1.3.
###### Proof.
By \[23, lemma 2.3\], we may lift the action of $`G`$ on $`S^1`$ to an action of a group $`H`$ on $``$; specifically $`H`$ is a left orderable group with a central subgroup $`Z`$ such that $`Z`$ and $`H/ZG`$. Note that $`H`$ is finitely generated because $`G`$ is finitely generated. If $`G`$ is finite, then $`H`$ is a torsion free group with an infinite central cyclic subgroup of finite index and it follows that $`H`$. We deduce that $`G`$ is cyclic. Therefore we may assume that $`G`$ is infinite.
By Theorem 1.2 $`H`$ has a normal subgroup $`K`$ such that $`H/K`$. If $`KZ1`$, then $`Z/KZ`$ is finite, consequently $`KZ/K`$ is a finite subgroup of $`H/K`$ and we deduce that $`H/KZ`$. It follows that $`G`$ has an infinite cyclic quotient, so we may assume that $`KZ=1`$.
Note that $`H/KZ`$ is a finite cyclic group. Since $`KZ`$ has finite index in the finitely generated group $`H`$, we see that $`K`$ is finitely generated. Moreover $`K`$ is infinite, so by Theorem 1.2 there exists $`LK`$ such that $`K/L`$. But
$$\frac{KZ}{LZ}\frac{K}{(LZ)K}=\frac{K}{L(ZK)}=\frac{K}{L}$$
and the result follows. ∎
## 6. Free products with amalgamation
###### Lemma 6.1.
Let $`G\mathrm{Aut}()`$. Then there is an action $`\alpha `$ of $`G`$ on $``$ by orientation preserving homeomorphisms with the following properties.
1. If $`cG`$ and $`c(r)>r`$ for all $`r`$, then $`(\alpha c)r<r`$ for all $`r`$.
2. If $`cG`$ and $`c(r)<r`$ for all $`r`$, then $`(\alpha c)r>r`$ for all $`r`$.
###### Proof.
Define an action $`\alpha `$ of $`G`$ on $``$ by $`(\alpha g)r=g(r)`$ for $`gG`$. This action has the required properties. ∎
###### Lemma 6.2.
Let $`G\mathrm{Aut}()`$, let $`H`$ be a left ordered group, let $`1cG`$, let $`C=c`$, and let $`1hH`$. Identify $`C`$ with $`h`$ via the isomorphism $`c^nh^n:Ch`$ for $`n`$. Suppose $`\mathrm{Fix}_{}(c)=\mathrm{}`$. Then $`G_CH`$ is left orderable.
###### Proof.
Write $`H=_iH_i`$, where the $`H_i`$ are finitely generated subgroups containing $`h`$. Then $`G_CH=_iG_CH_i`$, and if each of the $`G_CH_i`$ is left orderable, then so is $`G_CH`$ by \[18, 7.3.2\]. Therefore we may assume that $`H`$ is finitely generated. Using Lemma 2.2, we can view $`H`$ as a subgroup of $`\mathrm{Aut}()`$.
We may write $`\mathrm{Fix}_{}(h)`$ as a countable disjoint union of nonempty open sets, say $`_iP_i`$. On each $`P_i`$, either $`h(x)>x`$ for all $`xP_i`$, or $`h(x)<x`$ for all $`xP_i`$. Using Lemma 6.1, for each $`i`$ there is an action of $`G`$ on $`P_i`$ by orientation preserving homeomorphisms with the property that either $`h(x)>x`$ and $`c(x)>x`$ for all $`x`$, or $`h(x)<x`$ and $`c(x)<x`$ for all $`x`$. Then by \[11, theorem 10\] we may assume that $`h=c`$ on $`P_i`$. We have now defined an action of $`G`$ on $`_iP_i`$, and we extend this to an action $`\alpha `$ on the whole of $``$ by defining $`\alpha g`$ to be the identity on $`\mathrm{Fix}_{}(h)`$ for all $`gG`$. Clearly $`\alpha (G)\mathrm{Aut}()`$ and $`\alpha (G)G`$. Thus we can define a group homomorphism $`\theta :G_CH\mathrm{Aut}()`$ by $`\theta g=\alpha g`$ for $`gG`$ and $`\theta h=h`$ for $`hH`$, because $`\alpha (c^n)=h^n`$ for $`n`$. The result now follows from \[14, theorem 6.2.3\]. ∎
###### Theorem 6.3.
Let $`G`$ be a left ordered group, let $`H`$ be a torsion free nilpotent group, and let $`C`$ be a cyclic group. Then $`G_CH`$ is left orderable.
###### Proof.
If $`C=1`$ then the result follows from \[18, §2.4 on p. 37 and theorem 7.3.2\], so we may assume that $`C`$ is infinite cyclic. We will assume that $`C`$ is a subgroup of $`H`$ and write $`C=c`$, where $`1cH`$. Let $`1gG`$ and identify $`C`$ with $`g`$ via the isomorphism $`c^ng^n`$ for $`n`$. We need to prove that $`G_CH`$ is left orderable.
Write $`H=_iH_i`$, where the $`H_i`$ are finitely generated subgroups containing $`C`$. Then $`G_CH=_iG_CH_i`$, and if each of the $`G_CH_i`$ is left orderable, then so is $`G_CH`$ by \[18, 7.3.2\]. Therefore we may assume that $`H`$ is finitely generated. We shall use induction on the Hirsch length of $`H`$ (so if $`1=H_0H_1\mathrm{}H_n=H`$ is a normal series for $`H`$ with $`H_i/H_{i1}`$ infinite cyclic for all $`i`$, then $`n`$ is the Hirsch length of $`H`$).
First suppose the Hirsch length of $`H`$ is 1. This means that $`H`$ is infinite cyclic, say $`H=h`$ where $`h`$ has infinite order. Then we can view $`H`$ as a subgroup of $`\mathrm{Aut}()`$ by letting $`H`$ act on $``$ according to the rule $`h(r)=r+1`$ for all $`r`$. Then $`\mathrm{Fix}_{}(c)=\mathrm{}`$ and the result follows from Lemma 6.2. Therefore we may assume that the Hirsch length of $`H`$ is at least 2.
Let $`Z`$ be a nontrivial cyclic central subgroup of $`H`$ such that $`H/Z`$ is torsion free. Then $`H/Z`$ is left orderable because $`H/Z`$ is a torsion free nilpotent group \[18, §2.4 on p. 37\].
Suppose $`cZ`$. We have an epimorphism $`G_CHG_{CZ/Z}H/Z`$. Let $`K`$ be the kernel of this map. Then $`KG=1`$ and $`KH=Z`$, so applying \[7, I.7.7\] we see that $`K`$ is free and consequently left orderable. By induction $`G_{CZ/Z}H/Z`$ is left orderable, and so the result follows from \[18, 7.3.2\].
Finally we need to consider the case $`cZ`$. We have an epimorphism $`G_CHH/Z`$. Let $`K`$ be the kernel of this map, and let $`F=G_CH`$. With $`G_CH`$ we have an associated standard tree $`T`$ \[7, I.3.4 definitions\], and $`F`$ acts on this tree. The vertices of $`T`$ are the left cosets $`fG`$ and $`fH`$, and the edges are the left cosets $`fC`$, where $`fF`$. A fundamental $`F`$-transversal \[7, 2.6 proposition\] for $`T`$ consists of the vertices $`G,H`$ and the edge $`C`$. Let $`X`$ be a transversal for $`Z`$ in $`H`$. Then a fundamental $`K`$-transversal $`T_0`$ for $`T`$ consists of the vertices $`xG`$ and $`xH=H`$, and the edges $`xC`$, where $`xX`$. The stabilizers of the vertices of $`T_0`$ are of the form $`G^x`$ and $`Z^x=Z`$, and the stabilizers of the edges are of the form $`C^x=C`$ for $`xX`$. It follows that $`K`$ is the fundamental group of a graph of groups \[7, I.3.4 definitions\] of the following form
where each $`G_i`$ is of the form $`G^x`$ for some $`xX`$ (where $`x`$ depends on $`i`$). We can now define an epimorphism $`\theta :KG_CZ`$ by $`\theta g=x^1gx`$ for $`gG_i`$ and $`\theta z=z`$ for $`zZ`$. The kernel of this map is a free group and hence left orderable. Also $`G_CZ`$ is left orderable by induction. We now apply \[18, theorem 7.3.2\] twice to first deduce that $`K`$ is left orderable, and then $`G_CH`$ is left orderable, as required. ∎
###### Problem 6.4.
Is the free product of two left orderable groups with an amalgamated cyclic subgroup left orderable?
## 7. Examples of left ordered groups which are not locally indicable
If $`G`$ is a group, we shall let $`\mathrm{\Delta }(G)`$ indicate the finite conjugate center of $`G`$, that is $`\{gG\text{C}_G(g)`$ has finite index in $`G\}`$. Let $`n`$ be a positive integer and let $`B_n`$ denote the Braid group on $`n`$ strings with standard generators $`\sigma _1,\mathrm{},\sigma _{n1}`$.
Dehornoy (see also ) has proven that the Braid group $`B_n`$ on $`n`$ strings is left orderable. Therefore $`B_n^{}`$ is also left orderable. It is not difficult to see that $`B_n^{}`$ is a finitely generated perfect group with trivial center for $`n5`$; we shall give a proof of this (probably known) result in Lemma 7.2 below. Thus for $`n5`$, we see that $`B_n^{}`$ is a nontrivial finitely generated left orderable perfect group with trivial center; see the last paragraph of \[2, p. 248\].
###### Lemma 7.1.
Let $`n`$ be a positive integer. Then $`\mathrm{\Delta }(B_n)=\text{Z}(B_n)`$.
###### Proof.
Obviously $`\text{Z}(B_n)\mathrm{\Delta }(B_n)`$. Conversely suppose $`\beta \mathrm{\Delta }(B_n)`$. Then the centralizer of $`\beta `$ in $`B_n`$ has finite index in $`B_n`$, consequently it contains a normal subgroup $`C`$ of finite index $`r`$ in $`B_n`$. Thus $`\sigma _i^rC`$ for all $`i`$, so by \[10, 2.2 theorem\] we see that $`\sigma _i\beta =\beta \sigma _i`$ for all $`i`$. Therefore $`\beta \text{Z}(B_n)`$ and the result is proven. ∎
###### Lemma 7.2.
Let $`n`$ be a positive integer. Then $`B_n^{}`$ is finitely generated and $`\text{Z}(B_n^{})=1`$. Furthermore if $`n5`$, then $`B_n^{\prime \prime }=B_n^{}`$.
###### Proof.
The result is trivial if $`n2`$, so we may assume that $`n3`$. Let $`Z=\text{Z}(B_n)`$. Then \[10, 2.5 corollary\] shows that $`Z=(\sigma _1\mathrm{}\sigma _{n1})^n`$, and we now see from \[13, exercise 7, p. 47\] that $`ZB_n^{}=1`$. Also $`B_n/B_n^{}`$ from \[15, p. 757\] and we deduce that $`B_n^{}Z`$ has finite index in $`B_n`$. Therefore $`B_n^{}`$ is finitely generated and $`\text{Z}(B_n^{})\mathrm{\Delta }(B_n)`$. But $`\mathrm{\Delta }(B_n)=Z`$ by Lemma 7.1 and the first part is proven. Finally if $`n5`$, then $`B_n^{\prime \prime }=B_n^{}`$ from \[15, p. 757\]. ∎
Let $`\stackrel{~}{G}`$ denote the group of piecewise linear homeomorphisms of $``$ which satisfy $`g(x+1)=g(x)+1`$ for all $`g\stackrel{~}{G}`$ and $`x`$, as described in . Thus $`\stackrel{~}{G}`$ is a finitely generated perfect group with infinite cyclic center $`Z`$ generated by the map $`xx+1`$ for $`x`$, and $`\stackrel{~}{G}/Z`$ is a simple group, called $`T`$ in . Then the free product $`\stackrel{~}{G}\stackrel{~}{G}`$ is a finitely generated perfect group with trivial center, and is left orderable by \[18, theorem 7.3.2\].
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# Proper motions of pre-main sequence stars in southern star-forming regions Based on observations made at Valinhos CCD Meridian Circle. Based on measurements made with MAMA automatic measuring machine Table 4 will be only available in electronic form at the CDS via anonymous ftp to cdsarc.u-strasbg.fr (130.79.128.5) or via http://cdsweb.u-strasbg.fr/Abstract.html
## 1 Introduction
Analysis of the motion of pre-main sequence (PMS) stars and of related groups of young stars provide essential tests of star formation models. Different space velocities and velocity gradients of the stellar associations can be derived from the major star formation scenarios, like sequential star formation, star formation by high-velocity clouds, Gould’s Belt models, etc. Proper motion measurements of the members of these associations provide one way to discriminate among these predictions. The PMS stars are supposed to be sufficiently young to be very close to their birthplaces and to have velocities still very similar to the initial ones, so that one can get clear constraints on the birth mechanism.
In this work, we investigated the PMS stars of an extended region of aligned molecular clouds and OB associations that includes the Chamaeleon, Lupus, Upper Scorpius - Ophiuchus and Corona Australis regions. This selected area is specially interesting because the associations are close enough to the Sun (100-150 pc), so that a refined kinematical study can be made.
The HIPPARCOS mission (ESA 1997) provided accurate measurements of positions, parallaxes and proper motions, for OB stars brighter than V=10 mag, allowing to study the kinematics of the same regions (de Zeeuw et al. zeu (1999)). Proper motions of PMS stars associated with these star-forming complexes, based on HIPPARCOS data, were recently obtained (Frink et al. fri (1998), Neuhäuser & Brandner neu (1998), Wichmann et al. wib (1998), Bertout et al. ber (1999)), providing a first comparison of space velocities of different groups of stars. However, since most PMS stars are fainter than the limiting magnitude of HIPPARCOS, the total number of measured PMS stars is still small.
In this work, we present proper motion determination for 213 PMS stars as faint as V$``$16 mag lying in the galactic longitude range $`l=290\mathrm{°}`$ to $`l=360\mathrm{°}`$.
In the following sections, we present the method used to derive proper motions based on different combinations of first and second epoch measurements, we discuss the quality of our results and present a first analysis of this large sample of proper motions obtained for the PMS stars.
## 2 Data
For this work, we collected a list of well known and candidate T Tauri and Herbig Ae/Be (hereinafter HAeBe) stars as exhaustive as possible for the studied regions in the literature (Alcalá ala (1994); Alcalá et al. alb (1995); Bertout et al. ber (1999); Brandner et al. bra (1996); Casanova et al. cas (1995); Feigelson & Kriss fea (1989); Feigelson et al. feb (1993); Gauvin & Strom gau (1992); Gregorio-Hetem et al. gre (1992); Hartigan har (1993); Herbig & Bell heb (1988); Krautter et al. kra (1997); Malfait et al. mal (1998); Marraco & Rydgren mao (1981); Martín et al. mat (1998) , Preibisch et al. pre (1998); Schwartz sca (1977); Thé et al. the (1994); Torres et al. toa (1995); van den Ancker et al. vaa (1997); van den Ancker et al. vab (1998); Walter et al. wab (1997); Wichmann et al. wia (1997); Wilking et al. wil (1992)). With these data, we constructed an input catalogue containing about 680 stars with their approximate positions and magnitudes. In this set of stars, we could identify and measure proper motions of 213 PMS stars as faint as V$``$16 mag. We also present, separately, measurements for 29 stars that were previously considered as PMS, but subsequently classified differently by Covino et al. (cov (1997)) and Wichmann et al. (wic (1999)), who did not confirm their PMS nature. These authors obtained high resolution spectra, and compared the Li line equivalent widths to those of the Pleiades stars, considered by them as a frontier between PMS and non-PMS stars.
The proper motions were obtained from the comparison between several sources including current epoch CCD meridian observations performed at Valinhos Observatory (Viateau et al. via (1999)) and old SERC–J Schmidt plates measured with MAMA automatic measuring machine (Guibert et al. gui (1984)). These data were combined with published positions extracted from the following catalogues: AC2000 (Urban et al. 1998a ), USNO–A2.0 (Monet mon (1998)), HIPPARCOS and Tycho (ESA esa (1997)). One star had its proper motion obtained from a combination of CCD meridian observation and a Digitized Sky Survey<sup>1</sup><sup>1</sup>1The Digitized Sky Survey was produced at the Space Telescope Science Institute under US Government grant NAG W–2166. (DSS) image.
### 2.1 CCD meridian observations
The most recent observations used here were obtained with the Valinhos CCD meridian circle (Dominici et al. dom (1999)) in 1998 and 1999. This instrument is installed in Valinhos at the Abrahão de Moraes Observatory (Latitude $`23\mathrm{°}00\mathrm{}06\mathrm{}`$, Longitude $`+46\mathrm{°}58\mathrm{}03\mathrm{}`$) which belongs to the São Paulo University – Brazil.
The CCD detector has 512x512 square pixels of 19$`\mu m`$ ( 1 pixel = 1.5<sup>′′</sup> square) and works in drift scan mode. In this mode, the telescope is fixed and the electric charges are moved along the columns of the CCD with the same velocity as the transit, which depends on the declination. For $`\delta =0\mathrm{°}`$ the integration time interval is 51 seconds. The observed field has therefore an arbitrary length in right ascension (typically one hour) and is 13 wide in declination. The magnitude limit is about V=16.0 mag. The astrometric and photometric observational precisions depend on the magnitude and in the best interval ($`9.0<`$V$`<14.0`$) are, respectively, approximately $`0.050\mathrm{}`$ in both coordinates and 0.05 magnitudes. At the detection limit, positional measurements are less accurate and the mean square error may reach $`0.100\mathrm{}`$. Observations are carried out with a filter CG495$`+`$BG38 (bandpass from 5200 to 6800 Å), which is wider than Johnson filter. Notwithstanding, the resulting magnitudes are close to the visual standard magnitude system (Dominici et al. dom (1999)). The final positions are obtained by a global reduction procedure (Benevides-Soares & Teixeira ben (1992); Teixeira et al. tei (1992)) using the ACT (Urban et al. 1998b ) as the reference catalogue. The data treatment is made by means of a software package developed and maintained by J.F. Le Campion (Bordeaux Observatory).
In this work, we observed seven strips in the Chamaeleon, Lupus and Upper Scorpius - Ophiuchus regions containing many PMS stars identified in the literature. The length in R.A. of the fields presented in Table 1 was defined to ensure a minimum of 20 reference stars in each strip, necessary for an accurate astrometric reduction, and their coordinates were defined to have the largest number of PMS stars. In this table, we also present the central declination and the number of observations of each strip.
Each star considered here was observed at least three times.
### 2.2 Schmidt plate material
Thirteen SERC–J $`6.5\mathrm{°}\times 6.5\mathrm{°}`$ Schmidt plates have been digitized at the MAMA measuring machine (Guibert et al. gui (1984)), which provides at the present time the most accurate measurements (repeatability of $`0.4\mu m`$). For each plate, a catalogue of ($`x`$,$`y`$), flux and area has been produced for about 1 000 000 objects detected.
Each plate was then astrometrically reduced from ($`x`$$`y`$) to ($`\alpha `$$`\delta `$) with reference stars from the ACT catalogue (Urban et al. 1998b ). The mean residual of these reductions was about 0.25$`\mathrm{}`$ in both coordinates. We give in Table 2 the list of the SERC–J plates used in this work to derive proper motions for the known PMS stars. In this table are shown the numbers of the plates, their observation epochs and their central coordinates.
The Schmidt plates and Valinhos strips distribution in the studied region of the sky are given in Fig. 1, along with the searched PMS stars from our input catalogue.
### 2.3 Other data
For a good determination of proper motions, observations from well separate epochs are needed. Unfortunately, the catalogue with the oldest mean epoch (AC2000, Urban et al. 1998a ), gives positions for stars essentially brighter than V=12 mag. We included these positions in our calculations, when available. For fainter stars (V$`>`$12), proper motions have been determined over a shorter time interval and so with a degraded precision. One of these stars, Sz 108, although observed at the Valinhos observatory, could not be found in the other catalogues. Its position was determined from the DSS (plate identification: 02F8, region: S330), using the IRAF software packages.
We also included many other available astrometric data as from USNO–A2.0 Catalogue (Monet mon (1998)), Tycho and HIPPARCOS positions (ESA esa (1997)) to better constrain the proper motion determination and to provide a larger number of proper motions.
## 3 Proper motion determination
From our initial sample of about 680 stars, we could measure accurate proper motions for 242 of them. These stars have magnitudes within the range 6$`<`$V$`<`$16. Proper motions have been determined only when the time basis was longer than 20 years. In most cases this time basis was longer than 50 years, reaching in some individual cases more than 100 years. About 30 proper motions were determined with a time basis shorter than 50 years. The mean time basis was 80 years.
For the remaining stars, the main reasons of their absence in our final catalogue are that either their magnitudes laid outside the magnitude range covered by our data sources or the request of a time basis longer than 20 years was not fulfilled.
The proper motion calculation was performed in the usual way, via a weighted least squares method (Eqs. 1 to 5).
$`t_0`$ $`=`$ $`{\displaystyle \frac{p_it_i}{p_i}}`$ (1)
$`\alpha _0`$ $`=`$ $`{\displaystyle \frac{\alpha _ip_i}{p_i}}`$ (2)
$`\mu _\alpha `$ $`=`$ $`{\displaystyle \frac{p_i\alpha _i(t_it_0)}{p_i(t_it_0)^2}}`$ (3)
$`\sigma _{\alpha _0}^2`$ $`=`$ $`{\displaystyle \frac{1}{p_i}}`$ (4)
$`\sigma _{\mu _\alpha }^2`$ $`=`$ $`{\displaystyle \frac{1}{p_i(t_it_0)^2}}`$ (5)
where $`p_i=\frac{1}{\sigma _i^2}`$, and $`t_i`$ is the epoch of the position for a given star $`i`$. The same calculation holds for the declination.
We have assumed the following precisions for the various data : $`\sigma =0.25\mathrm{}`$ for AC2000, USNO–A2.0 and SERC–J positions, $`\sigma =0.001\mathrm{}`$ for HIPPARCOS, $`\sigma =0.030\mathrm{}`$ for Tycho and $`\sigma =0.050\mathrm{}`$ for Valinhos positions.
We present the precision of the derived proper motions with various material in Table 3, where $`\mathrm{\Delta }`$t is the mean time basis for the groups of sources given in column one, and N stands for the number of stars whose proper motions were obtained from the combination of these sources.
We give in Tables 4 and 5 the derived mean positions and proper motions for the 213 PMS stars. In these tables, the objects were separated in T Tauri stars (Table 4) and HAeBe stars (Table 5). Table 6 lists the 29 non-PMS ROSAT stars in Chamaeleon and Lupus (see section 2). In addition, a traditional separation by region was adopted.
In the second column of these tables, magnitudes are taken preferably from Tycho catalogue ($`V_T`$). For the Valinhos non-Tycho stars, Valinhos visual magnitudes are used. Otherwise, magnitudes are taken from the literature, as indicated by the references (column 3) explicited at the end of Table 6. For stars RX J1621.4–2332ab and RX J25.3–2402 we give the B magnitude from the AC2000. Magnitudes taken from reference $`[10]`$ are also B magnitudes.
Out of the 213 PMS stars studied here, 101 had no previous determination of proper motion, as far as the consulted literature is concerned, 81 of them had known proper motions from the ACT and 41 from HIPPARCOS. In the last column of Tables 4, 5 and 6, additional references for other available proper motions are provided. The proper motions in common with the ACT and HIPPARCOS catalogues were used to externally evaluate the quality of our results.
A comparison of the obtained proper motions with HIPPARCOS ones is presented in Fig. 2. We can notice the good agreement between both sets of data. The mean dispersion of the differences is 6 mas/yr, which is the best estimate of the external errors of our catalogue, concerning the bright PMS stars. For the fainter stars, no catalogues for comparison were found available.
We also compared our proper motions to the ACT ones in order to evaluate the importance of the addition of SERC–J plates and meridian observations to the material used by ACT. The mean dispersion of differences is 3 mas/yr.
One should stress that the depicted comparison with HIPPARCOS and ACT is not completely independent, but is useful for a first estimation of the coherence of our results.
## 4 Discussion
We present in Fig. 3 a general view of the spatial distribution and proper motions in galactic coordinates of the measured PMS stars, and in Figs. 4, 5, 6 and 7 – upper panels – zooms of the 4 main star-forming regions studied in this work. In all these regions, a dominant orientation of the proper motions towards smaller longitudes can be observed. This is in large part the effect of the reflex solar motion, as can be seen in the $`\mu _b`$ versus $`\mu _l\mathrm{cos}b`$ graphs of the same regions (Figs. 4, 5, 6 and 7 – lower panels). Since the effect of solar motion on the star proper motion depends on the distance of the stars, and the distances of PMS stars are usually poorly known, we present in these figures the reflex solar motion as a function of distance, from 50 to 200 pc. We assumed the basic solar motion, with components U = 9 km/s (directed towards the galactic center), V = 11 km/s, W = 6 km/s. The reflex solar motion depends also on the direction of the stars, and since some of the regions studied here have sizes of several degrees, we present it for two extreme directions in each field (except for Corona Australis, which is a small field).
In all groups, a number of stars are found to have proper motions that lie well outside the average distribution. Notice that some of these stars that can be seen in the ($`l,b`$) maps (for Upper Sco, Chamaeleon and Corona Australis) are absent from the ($`\mu _l\mathrm{cos}(b),\mu _b`$) graphs, due to the scale that we adopted. We consider that these stars have a large probability of being recently formed runaway stars and a possible explanation for them is the disruption of multiple systems. We cannot exclude the possibility of errors, like a wrong identification during the proper motion determination process; however, note that some of the stars with anomalous proper motions are bright and were found in several sources (eg. HD 137727, HD 140637), which reduces the probability of an erroneous identification.
Disregarding the runaway stars, the groups of PMS stars present proper motion dispersions of the order of 10 mas/yr typically. This means that in a period of about one million years, the stars would move apart several degrees in the sky, and would look very dispersed. We must bear in mind, however, that part of this apparent dispersion may not be intrinsic, but due to the fact that the sun is approaching the group. For instance, most of the PMS stars of Upper Scorpius present negative radial velocities, while most of the stars in Chamaeleon present positive radial ones (e.g. Gregorio-Hetem et al. gre (1992); Torres et al. toa (1995); Covino et al. cov (1997)).
Let us now comment on the average proper motion of the stellar groups. We already noticed that the average values are largely explained by the reflex solar motion. However, they are not entirely due to this effect or, in other words, intrinsic average proper motion of the groups are detected. For instance, in the case of Lupus (our largest sample, Fig. 5), the center of mass of the points in the ($`\mu _l\mathrm{cos}(b),\mu _b`$) graph is about ($``$30,$``$9), which suggests mean distance of about 85 pc, if we consider only the reflex solar motion. However, the average distance of 14 stars of this group that have parallaxes from HIPPARCOS is 138 pc. This suggests that the Lupus PMS stars have a mean intrinsic proper motion in the longitude direction of about $`\mu _l\mathrm{cos}(b)`$ = $``$10 mas/yr.
The Chamaeleon group is peculiar, in that it seems to present two distinct kinematic groups, even after excluding the runaway stars. A group with proper motions close to about $`\mu _l\mathrm{cos}(b)`$ = $``$40 mas/yr, $`\mu _b`$ = $``$15 mas/yr, seems to be well behaved. Its observed proper motion could be explained by the the reflex solar motion, if its distance is of the order of 70 pc. And indeed, some of the stars of this group have distances determined by HIPPARCOS, and are not too different from this value (T Cha, 66 pc, RX J1158.5-7754a, 86 pc, RX J1159.7-7601, 92 pc, HD 104237, 116 pc). This group was already discussed by Terranegra et al. (1999); notice, however, that Bertout et al. (ber (1999)) consider that the HIPPARCOS distance of T Cha is incorrect. Another group of stars presents positive values of $`\mu _b`$, considerably different from the reflex solar motion. Among these, only HD97300 has parallax measured by HIPPARCOS (187 pc).
Finally, we remark that in the studied regions, no systematic differences between the proper motions of T Tauri stars and HAeBe stars can be observed. Our results favour the PMS nature of the candidates HAeBe stars included in our list.
A deeper analysis of the proper motion of the groups of PMS stars, in connection with the ages of the subgroups, will be presented in a separate paper, where our results are compared with the models proposed in the literature for the mechanisms that might have triggered the star formation.
###### Acknowledgements.
The authors wish to thank Dr. M. Rapaport for helpful comments and Dr. J. Guibert for his help in finding and measuring the SERC–J plates. We are grateful to Dr. R. Neuhäuser, for his constructive suggestions on this paper. Special thanks are also extended to Mr. W. Monteiro, for his support concerning the observations. A partial financial support from CNPq, FAPESP, PRONEX, CAPES and CNRS is fully acknowledged. This work benefited from SIMBAD database and the Astronomical Data Analysis Center, the last one being operated by the National Astronomical Observatory of Japan.
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# For submission to Monthly Notices Letters Application of Data Compression Methods to the Redshift-space distortions of the PSCz galaxy catalogue
## 1 Introduction
The extraction of cosmological parameters from surveys has entered a new phase, with the advent of very large data sets. But the prospect of accurately determining a wide range of parameters is offset by the difficult task of manipulating these large data sets without loosing important parameter information. In the case of analysing near-all sky redshift surveys we have developed a method based on the harmonic decomposition of the survey into spherical harmonics and radial spherical Bessel functions (Heavens & Taylor 1995, Ballinger, Heavens & Taylor 1995, Tadros et al 1999; hereafter HT, BHT and T99, respectively). The use of a spherical harmonic decomposition allows the accurate analysis of the radial redshift-space distortion effect, without using the small-angle or distant observer approximations, and the natural inclusion of angular and radial window functions. The parameters of interest are the degree of redshift-space distortion, parameterised by<sup>1</sup><sup>1</sup>1Since differently selected galaxies cluster differently, each selection may have its own bias parameter and $`\beta `$. In this paper we shall use $`\beta `$ and $`b`$ to refer exclusively to IRAS selected galaxies.
$$\beta \frac{f(\mathrm{\Omega }_m)}{b}$$
(1)
where $`f(\mathrm{\Omega }_m)=d\mathrm{ln}\delta /d\mathrm{ln}a\mathrm{\Omega }_m^{0.6}`$ (Peebles 1980) is the growth rate of perturbations and $`b`$ is the linear bias parameter defined by
$$P_{\mathrm{gal}}(k)=b^2P_m(k),$$
(2)
and an estimate of the undistorted galaxy power spectrum, $`P_{\mathrm{gal}}(k)`$. This approach combines the spherical harmonic decomposition with a mode-by-mode maximum likelihood analysis, and has been applied to the 1.2Jy survey, a 1:6 subsample of the IRAS galaxy survey (HT, BHT) and the full 1:1 sample, the PSCz (T99). T99 used the most sophisticated analysis to date, and found a distortion parameter of $`\beta =0.58\pm 0.26`$ (marginal error) and an amplitude for the real space galaxy power spectrum at $`k=0.1h\mathrm{Mpc}^1`$ of $`\mathrm{\Delta }_{0.1}=0.42\pm 0.03`$, where $`\mathrm{\Delta }^2(k)=k^3P(k)/2\pi ^2`$. In this paper we employ new data compression methods to improve on this analysis and obtain more accurate parameter determinations. We describe the methods more fully in Ballinger et al (2000).
The limiting factors in our previous analyses were CPU time and stability. To complete the likelihood analysis the data covariance matrix of the full data set must be inverted at each point in parameter space. Small numerical errors can make this procedure unstable; data compression (Section 2) is a great help here, as the high signal-to-noise modes are well-behaved. The CPU factor can also be an issue when one wishes to investigate systematic effects in data sets. This is as much an issue in analysis of the Cosmic Microwave Background (CMB) as it is in galaxy redshift surveys, such as the PSCz, the 2-degree Field (2dF) or the Sloan Digital Sky Survey (SDSS).
The problem of analysing large data sets was addressed by Tegmark, Taylor & Heavens (1997; TTH) who considered the question of what was the optimal transformation of the data for estimating a given parameter, where the model data covariance matrix could be an arbitrary, nonlinear function of the desired parameter. The optimal transformation should have the properties of maximising the information content about the parameter in the minimum number of eigenmodes. By ordering modes the ones with the most information could be selected and the rest discarded. That the data can be ordered this way can be understood if one considers that the data may contain noisy or strongly correlated modes that add little information about the parameter of interest. Choosing the transformed data to have diagonal covariance also decreases the computation time. While the data covariance is only diagonal at one point in parameter space, the removal of correlated and noisy data by trimming produces a numerically stable inversion of the covariance matrix.
This procedure sounds similar to Principle Component Analysis (PCA) or signal-to-noise eigenvalue analysis (SNA; Bond 1995, Vogeley & Szalay 1995), but has important differences. We have previously referred to the procedure as Karhunen-Loève methods, but it is in fact more general, so we shall refer to our method henceforth as Optimal-Mode Analysis (OMA). In this paper, we introduce two new methods for accurate parameter estimation, and refer to the whole method as Generalised Optimal Mode Analysis, or GOMA.
The additional features of GOMA can be split into two parts, dealing with multiparameter estimation and the stable handling of data compression. In multiparameter estimation it is useful to note that for highly-correlated multiparameter estimates (highly-elongated likelihood surfaces, not aligned with parameter axes), the marginal error in the parameters is determined by the length of the longest likelihood principal axis. We therefore want to optimise to keep this length as short as possible. This process is called parameter eigenmode analysis, and was introduced in Ballinger (1997), with some results being presented in Taylor et al. (1997). The second, optional part of GOMA is to split the original data into subsets, optimising each subset, and then combining the best modes together. This procedure can be used hierarchically, to reduce a very large number of modes to a manageable size. This process we refer to as hierarchical data compression. These methods will be detailed in a companion paper (Ballinger et al. 2000).
In this paper we combine our spherical harmonic decomposition, parameter eigenmode analysis and hierarchical data compression methods to analyse the PSCz. We study both nonlinear multiparameter estimation, the redshift-space distortion parameter and the amplitude of power, using hierarchical data compression, from the harmonic modes of the PSCz survey. The increase in analysing power using these methods allows us to increase the number of modes available for study, and hence a corresponding increase in accuracy of our results.
Padmanabhan, Tegmark and Hamilton (1999) have also used the spherical harmonic decomposition to analyse the CFA/SSRS UZC galaxy redshift survey, while Hamilton, Tegmark & Padmanabhan (2000) have applied the analysis to the PSCz redshift survey. In both cases they employ Karhunen-Loève methods to estimate the quadratic band-power in these surveys. In the former survey the band-power was measured in redshift-space, while in the latter analysis of the PSCz they measured the real-space galaxy-galaxy, galaxy-velocity and velocity-velocity power spectra, and estimated $`\beta `$ from a least squares fit to the ratios of the galaxy-velocity to galaxy-galaxy and velocity-velocity to galaxy-galaxy power spectra.
The paper can be summarised as follows. In Section 2 we briefly review the spherical harmonic decomposition, OMA, parameter eigenmodes and hierarchical data compression methods. In Section 3 we perform an maximum likelihood analysis of the PSCz, for the redshift-space distortion parameter and real-space clustering amplitude of galaxies. Our results are presented in Section 4, and conclusions are made in Section 5. We begin by a brief review of our data analysis methods.
## 2 Data Analysis
### 2.1 Spherical Harmonic decomposition
Following HT, BHT, and T99 we can decompose the density field of galaxies in a redshift survey into harmonic modes
$$\rho _{\mathrm{}mn}=c_\mathrm{}nd^3s\rho (𝒔)w(s)j_{\mathrm{}}\left(ks\right)Y_\mathrm{}m^{}\left(\mathrm{\Omega }\right),$$
(3)
where $`Y_\mathrm{}m`$, are spherical harmonics, $`j_{\mathrm{}}`$ are a discrete set of spherical Bessel functions, $`w(s)`$ is an adjustable weighting function and $`𝒔`$ is the redshift-space position variable. The $`c_\mathrm{}n`$ are normalization constants and $`k=k_{ln}`$ are discrete wavenumbers. Each mode was weighted with a Feldmann-Kaiser-Peacock (1994) weight, $`w(k,s)=[1+\varphi (s)P(k)]^1`$, where $`\varphi (s)`$ is the (redshift-dependent) average number density of galaxies in the survey.
The observed coefficients, $`D_{\mathrm{}mn}`$, can be related to the underlying linear density modes, $`\delta _{\mathrm{}mn}`$, by (HT, BHT, T99)
$$𝑫𝝆𝝆_\mathrm{𝟎}=𝑺𝑾(𝚽+\beta 𝑽)𝜹.$$
(4)
The transition matrices $`𝑺`$, $`𝑾`$, $`𝚽`$ and $`𝑽`$ correspond to the effects of small-scale random radial velocity distortions, the angular window function, the radial galaxy selection function and linear redshift space distortions, respectively. The mean field, $`𝝆_\mathrm{𝟎}`$, is nonzero due to the radial selection function and angular window function. Note that there is no matrix multiplication implied between the angular matrix, $`𝑾`$, and the radial matrix, $`𝚽+\beta 𝑽`$, since these are orthogonal.
The small-scale radial velocities can be accurately modelled by multiplying the galaxy power spectrum by a Lorentzian function. This implies that each mode is multiplied by the square-root of a Lorentzian. Inverse Fourier transforming we find that the density field should be convolved with the function
$$S(x)=\frac{2\sqrt{2}}{\sigma _v}K_0\left(\frac{\sqrt{2}}{\sigma _v}x\right),$$
(5)
where $`K_n`$ is an $`n^{th}`$-order modified Bessel function and $`\sigma _v`$ is the 1-d velocity dispersion. (Note that in the analysis of T99, the 3-d velocity dispersion was used incorrectly. Changing to the correct value had little effect. However, as we are pushing our model to higher wavenumbers here, it is more important to have the correct value of the velocity dispersion.) The transition matrix, $`𝑺`$, is found by a spherical harmonic transformation of $`S(xy)`$.
Two immediate advantages of this treatment are accounting for the effects of the monopole and dipole modes. The monopole mode contains information about the mean density of the survey. In previous methods this can bias down the estimated power as the mean is estimated from the survey itself and may not be the true mean (Tadros & Efstathiou 1996). In our treatment the monopole mode can be removed, effectively removing this bias (T99). As the PSCz is not all sky, some aliasing takes place at the few percent level, and we include the effects of this leakage, through $`𝑾`$. The dipole includes contributions from our own motion in the redshift space distortion. Again we can mostly remove this by ignoring the dipole mode, and accounting for aliasing from the angular mask (T99).
The distribution of linear harmonic modes can be described by a multivariate likelihood,
$$2\mathrm{ln}[𝑫|\beta ,P(k)]=𝑫^t𝑪^1𝑫+\mathrm{ln}det𝑪,$$
(6)
where $`𝑪=𝑫𝑫^t`$ is the data covariance matrix. Details for dealing with non-axisymmetric angular window functions in the covariance matrix, as well as further details of the likelihood analysis, are given in T99.
### 2.2 Generalised Optimal Mode Analysis (GOMA)
The technical advance which allows us to reduce the error bars in the determination of $`\beta `$ and the power spectrum is a new optimised form of data compression, which we call generalised optimal mode analysis. Details of the method will appear elsewhere (Ballinger et al. 2000), but we sketch the main ingredients here.
The need for data compression is twofold: first, the speed of analysis generally scales as $`N^3`$, where $`N`$ is the number of modes considered. These modes might be, for example, spherical transform coefficients. Secondly, since the covariance matrix is computed numerically, small numerical errors can lead to a non-positive definite matrix, which makes no physical sense; even a single negative eigenvalue out of several thousand is fatal. Instead of working with the full set of spherical transform coefficients up to some maximum wavenumber, we form orthogonal linear combinations, and use these as the data. GOMA consists of two parts; one is parameter eigenmode analysis, where instead of choosing $`\beta `$ and $`\mathrm{\Delta }_{0.1}`$ as the two parameters to be determined, we introduce a new parameter, which runs along the longest likelihood ridge of Figure 1, and use the data compression methods of TTH to make the likelihood as narrow as possible in this direction. Since the length of this controls the marginal errors in both $`\beta `$ and $`\mathrm{\Delta }_{0.1}`$, the method is very effective. It is worth noting that only OMA will determine the best set of eigenmodes following this procedure, since we now have a linear combinations of parameters which are nonlinear in the covariance matrix.
The second (optional) ingredient is hierarchical data compression. We cannot find the optimal linear combinations of the entire mode set, because of the numerical problems identified above. Instead, we form optimal orthogonal linear combinations of subsets of the modes (in discrete wavenumber ranges), and then combine the best modes to form a new set. This process can be repeated hierarchically to produce a near-optimal set of modes. Note that when mode sets are combined they are not orthogonal and we use the correct correlation properties.
## 3 Application to the PSCz
We have applied our methods to the new IRAS PSCz redshift survey (Saunders et al 1999), with a flux cut of $`0.75`$Jy, and the ultra-conservative mask created by T99, corresponding to optical extinctions of $`A_B=0.75`$mag and excluding $`35\%`$ of the sky. Our flux cut is above the formal limit of of $`0.6`$Jy for the catalogue, but motivated by a systematic effect we found linked to flux in T99. This appeared as a flux-dependency on the amplitude of real space perturbations, where the low-flux galaxies had roughly a factor of two higher clustering amplitude, the cause of which we were unable to identify. However by cutting the catalogue at $`0.75`$Jy the effect could be removed. This cut in flux, along with the more conservative mask left us with around 7000 galaxies.
The survey was put in a sphere of radius $`r_{\mathrm{max}}=200h^1\mathrm{Mpc}`$ and transformed into spherical harmonic modes. We analysed all the modes down to $`k=0.2h\mathrm{Mpc}^1`$, with $`n=114`$ and $`\mathrm{}=236`$ yielding a total of $`4644`$ harmonic modes. Modes in the range $`n=020`$ and $`\mathrm{}=050`$ were used for the convolution matrices. Since we use a higher cut in wavenumber, $`k`$, than in the analysis of T99 we must be cautious about introducing nonlinear modes which go beyond our analysis. The main concern is the effect of “fingers-of-god” contaminating our analysis. The effect of these would be to lower the measured value of $`\beta `$, since its effect is to elongate structure along the line of sight, and lower the clustering amplitude. We have tested our methods using CDM simulations under similar conditions, and find that our correction for radial small-scale velocities, equation (5), is accurate (Ballinger et al 2000).
The $`4664`$ modes were compressed down to $`2278`$, using the hierarchical compression method once, and a new likelihood analysis applied. In T99 around $`1300`$ modes where analysed. In this analysis we have partly made use of the increase in computing power in the time between the two analyses. However the real limiting feature of our previous analysis was numerical instability problems in the covariance matrix. Small errors can produce a covariance matrix which is not positive-definite, which makes no physical sense (a probability distribution in data space which grows exponentially along one principal direction). A great advantage of the compression mechanism is that this numerical problem is completely solved: the best modes for parameter estimation are well-behaved, and the covariance matrix inversion proceeds smoothly. Hence the overall time for computing our results (around 1 weeks CPU time) remained constant.
## 4 Results
### 4.1 Likelihood Analysis
Figure 1 shows the results of our likelihood analysis for the $`\beta \mathrm{\Delta }_{0.1}`$ plane. The larger contours are the results of T99 for 1300 modes. Contours are plotted at intervals of $`0.5`$ in $`\mathrm{ln}`$ from the maximum. The smaller set of contours are for our new analysis for $`4644`$ harmonic modes compressed to $`2278`$. We have used a CDM-type power spectrum with shape parameter $`\mathrm{\Gamma }=0.2`$ to calculate the covariance matrices, leaving the amplitude a free parameter, and in the wavenumber-dependent mode weighting function. We used a value of $`\sigma _v=224\mathrm{km}s^1`$ for the 1-dimensional velocity dispersion of galaxies in the scattering matrices, but have also experimented with $`112\mathrm{km}s^1`$ and $`336\mathrm{km}s^1`$.
The increase in parameter information yields a new, lower value of
$$\beta =0.39\pm 0.12,$$
(7)
where we quote the marginalized errors. This is consistent with our previous result, but with a much smaller error ellipse, around a factor of 3 improvement, which significantly rules out both $`\beta =1`$ and $`\beta =0`$. The amplitude of the real space galaxy power spectrum is
$$\mathrm{\Delta }_{0.1}=0.42\pm 0.02$$
(8)
which is again consistent with the results of T99 with a slightly smaller error. The covariance parameter of $`\beta `$ and $`\mathrm{\Delta }_{0.1}`$ is $`r=0.82`$, estimated from the ellipticity of the error ellipse. We find no evidence that our analysis is being biased by nonlinear effects, since then we would expect a sudden change in both the value of $`\beta `$ and a drop in the amplitude of real-space perturbations, neither of which is seen. We also find that the maximum likelihood values are almost unchanged, moving by less than $`\mathrm{\Delta }\beta 0.1`$, if we change the velocity dispersion to $`112`$ or $`336\mathrm{km}s^1`$, suggesting that nonlinear effects are not significant for this analysis.
As the amplitude of galaxy perturbations is proportional to the bias factor, we can combine $`\beta `$ and the linear galaxy power spectrum to estimate the amplitude of the mass power spectrum,
$$\mathrm{\Delta }_m(k)=\beta \mathrm{\Delta }_{\mathrm{gal}}(k)\mathrm{\Omega }_m^{0.6}.$$
(9)
The real-space power spectrum of the PSCz was estimated by T99, and is plotted in Figure 2. Also shown is the amplitude of mass perturbations for $`\mathrm{\Omega }_m=0.16`$. The point at $`k=0.23h\mathrm{Mpc}^1`$ is the mass amplitude on a scale of $`8h^1\mathrm{Mpc}`$, estimated from the abundance of clusters (Henry & Arnaud 1991, White, Efstathiou & Frenk 1993, Viana & Liddle 1996, Eke, Cole & Frenk 1996),
$$\sigma _8=0.56\mathrm{\Omega }_m^{0.47},$$
(10)
with a conservative error of around $`30\%`$. The mass amplitude from the PSCz at $`k=0.1h\mathrm{Mpc}^1`$ is
$$\mathrm{\Delta }_m(0.1)=(0.16\pm 0.04)\mathrm{\Omega }_m^{0.6}.$$
(11)
In Figure 2 we plot the mass amplitudes for the PSCz and the cluster abundance. There is a clear consistency between these two estimates of the mass amplitudes. Assuming bias is linear on large scales, estimating mass amplitudes from large-scale redshift surveys is simpler, and hence in principle more robust, than the abundance of clusters argument. In addition this can be used to sample the linear spectrum of mass perturbations on a range of scales, rather than being restricted to one scale.
### 4.2 Cosmological Parameters
Although the shape of the real-space, linear galaxy power spectrum can be used to constrain models of structure formation, the range of points in the linear regime is limited. A better test of models is to compare the mass amplitudes from the PSCz against the amplitude estimated from the CMB. To span the scales between different measurements we need to assume a model linear mass power spectrum. We shall assume that a standard CDM-type power spectrum of the form
$$\mathrm{\Delta }_m^2(k)=Q^2(k/H_0)^4T^2(k,\mathrm{\Omega }_m,h),$$
(12)
where we use the CDM transfer function of Bond & Efstathiou (1984), which gives a reasonable fit for a range of CDM models. This provides us with a 3-parameter model dependent on the parameter set $`\mathrm{\Omega }_m`$, $`h`$ and $`Q`$<sup>2</sup><sup>2</sup>2Our parameter Q is equivalent to the parameter $`\delta _H`$, the amplitude of clustering at the present Hubble scale, used elsewhere.. This is an interesting parameter set, in particular since CMB data alone cannot constrain $`\mathrm{\Omega }_m`$ without the addition of LSS data.
If we leave $`h`$ as a free parameter, the fit diverges off to high $`h`$ and $`Q`$ and low $`\mathrm{\Omega }_m`$, preserving the shape, but increasing the clustering amplitude. To avoid this we impose the HST constraint that $`h=0.65\pm 0.12`$. This effectively removes a degree of freedom in the fit. Beyond this we fit the PSCz mass points, the abundance of clusters constraint and the COBE 4-yr normalisation for a flat universe (Bunn & White 1997);
$$Q=(1.94\pm 0.2)\times 10^5\mathrm{\Omega }_m^{0.79}.$$
(13)
We implicitly assume that the universe is spatially flat, as implied by the recent Boomerang (Lange et al 2000) and MAXIMA (Hanany et al. 2000, Balbi et al. 2000) results and that the spectral index is unity, again as implied by CMB results. We also implicitly assume that the contribution to the mass density from baryons and neutrinos is negligible.
Figure 3 shows the $`\chi ^2`$ distribution for the $`\mathrm{\Omega }_mQ`$ plane, with $`h=0.65`$. The minimum $`\chi ^2`$ values are
$$\mathrm{\Omega }_m=0.16\pm 0.03,Q=(8.4\pm 3.8)\times 10^5$$
(14)
where we quote marginal errors. We can also obtain an estimate of the bias parameter for IRAS galaxies:
$$b=0.84\pm 0.28,$$
(15)
and the spectral shape parameter,
$$\mathrm{\Gamma }=0.1\pm 0.03,$$
(16)
which is somewhat lower than the usual value of $`0.2`$ from the shape of the galaxy clustering spectrum alone. This best fit has a $`\chi ^2=2`$ for 3 degrees of freedom, and so lies within the range of acceptable fits. However, while suggestive, this analysis makes more assumptions than one would like. These assumptions can be dropped by a combined analysis of the recent Boomerang and MAXIMA results of the small-angle fluctuations in the CMB, which goes beyond our present analysis.
In Figure 2 we plot the PSCz real-space galaxy power spectrum and real-space mass power spectrum, along with our best-fit model, a $`\mathrm{\Lambda }`$CDM, normalised to the COBE 4-year data. The model fits over two orders of magnitude in scale, but nonlinear and nonlocal bias effects lead to a disagreement on smaller scales, below $`k=0.3h\mathrm{Mpc}^1`$. Recent results from Seljak (2000), Ma & Fry (2000), and Peacock & Smith (2000) suggest that this nonlinear regime can best be understood in terms of a superposition of randomly positioned collapsed clusters, integrated over the cluster mass function. Seljak (2000), and Peacock and Smith (2000) go further to argue that the nonlinear bias function of galaxies arises due to the statistics of the occupation number of galaxies in haloes and show agreement between the $`\mathrm{\Lambda }`$CDM model and the small-scale power-law spectrum of galaxies. Perhaps the analysis of this regime is not as daunting as it once appeared.
## 5 Conclusions
In this paper we have presented a new analysis of the PSCz redshift survey. Using the spherical harmonic analysis of HT, BHT and T99 to decompose the survey, taking into account the effects of linear redshift-space distortions, nonlinear “fingers-of-god” effects, limited sky coverage, and the radial distribution of galaxies in the survey, we have applied a likelihood analysis to a conservative cut of the PSCz survey. The catalogue was cut at $`0.75`$Jy and a conservative mask used to avoid systematic uncertainties at the low-flux end of the catalogue and near the galactic plane. The spherical harmonic analysis of the remaining $`7042`$ galaxies resulted in $`4644`$ harmonic modes. We used a hierarchical data compression to reduce this to $`2278`$ for the final analysis. The compression was applied using a parameter eigenmode approach that compresses information along the line of largest degeneracy in parameter space, thus “squeezing” the uncertainty in this direction.
We used these methods to estimate the redshift-space distortion parameter, $`\beta `$ and the amplitude of the real-space galaxy power spectrum, parameterised by the amplitude, $`\mathrm{\Delta }_{0.1}`$, at a wavenumber $`k=0.1h\mathrm{Mpc}^1`$. Applying the likelihood analysis to wavenumbers below $`k=0.2h\mathrm{Mpc}^1`$, where linear theory will hold and “fingers-of-god” effects can be corrected for, we find
$`\beta `$ $`=`$ $`0.39\pm 0.12,`$ (17)
$`\mathrm{\Delta }_{0.1}`$ $`=`$ $`0.42\pm 0.02`$ (18)
quoting marginalised uncertainties. The distortion parameter is slightly lower than the uncompressed analysis of Tadros et al, with an uncertainty reduced by over a factor of 2. The consistency of these results with that of the earlier analysis of T99 leads us to believe that our analysis has not been heavily contaminated by nonlinear effects, while the conservative cuts in flux and sky coverage avoid contamination by flux systematics. These results also are in agreement with other, independent determinations of the distortion parameter. Comparing the velocity field reconstructed from the PSCz with the ENEAR survey, Nusser at al (2000) find $`\beta =0.5\pm 0.1`$, while Valentine, Saunders & Taylor (2000) find $`\beta =0.5\pm 0.1`$ by comparing the reconstructed dipole with the observed dipole, the bulk flow of galaxies with the MKII survey and the dipole again with the SFI catalogue. Finally, using a least squares fit to the ratios of the galaxy-velocity to galaxy-galaxy and velocity-velocity to galaxy-galaxy power spectra, Hamilton, Tegmark & Padmanabhan find $`\beta =0.41\pm 0.13`$ for the PSCz.
Since in the linear regime the galaxy power spectrum is likely to be proportional to the mass power spectrum we can combine $`\mathrm{\Delta }_{0.1}`$ and $`\beta `$ to find the amplitude of the mass perturbations,
$$\mathrm{\Delta }_m(0.1)=(0.16\pm 0.04)\mathrm{\Omega }_m^{0.6}.$$
(19)
Combined with the constraints from the CMB, and HST observations of the Hubble parameter, and assuming CDM, a flat universe, scale invariant initial mass perturbations and negligible contribution to the total mass-density from baryons and neutrinos we find $`\mathrm{\Omega }_m=0.16\pm 0.03`$, and a bias parameter for IRAS galaxies of $`b=0.84\pm 0.28`$. The minimum $`\chi ^2`$ fit has a value of $`\chi ^2=2`$ for 3 degrees of freedom.
## Acknowledgements
We thanks the PSCz team, and especially Will Saunders, for help in understanding the intricacies of the PSCz catalogue. ANT, WEB and HT thank the PPARC for postdoctoral support. Computations were made using STARLINK facilities.
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# Variations on the theme of Marcinkiewicz’ inequality
## 1 Introduction
Let $`h`$ be a smooth function on $`𝐑`$ with a compact support, and
$$g(x)=\text{p.v.}\frac{1}{\pi }_𝐑\frac{h(t)}{tx}𝑑t$$
be its Hilbert transform. Set $`f=g+ih`$ and introduce the function
$$u_f(z)=_𝐑H(zf(t))𝑑t,H(z)=\mathrm{log}|1z|+\text{Re}(z),$$
called the logarithmic determinant of genus one. It is subharmonic in $`𝐂`$, and its Riesz measure is $`d\mu _f(\zeta )=d\nu _f(\zeta ^1)`$, where $`d\nu _f`$ is the distribution measure of $`f`$:
$$\nu _f(E)=\text{meas}\left(\{t:f(t)E\}\right),E\text{is a borelian subset of}𝐂,$$
and $`\text{meas}(.)`$ stands for the Lebesgue measure on $`𝐑`$. Let
$$\mu _f(r)=\mu _f(\{|z|r\})=\text{meas}\left(\{|f|r^1\}\right)$$
be a (conventional) counting function of $`d\mu _f`$, and let
$`𝔫_f(r)`$ $`=`$ $`\mu _f\left(\{|zir/2|r/2\}\right)+\mu _f\left(\{|z+ir/2|r/2\}\right)`$
$`=`$ $`\mu _f\left(\{|\text{Im}(z^1)|r^1\}\right)=\text{meas}\left(\{|h|r^1\}\right)`$
be its Levin-Tsuji counting function, see , , and \[5, Chapter 1\]. Then the classical estimates of the Hilbert transform can be easily rewritten as upper bounds of $`\mu _f(r)`$ by $`𝔫_f(r)`$.
For example, Marcinkiewicz’ inequality (see \[7, Chapter V\])
$$m_f(\lambda )C\left\{\frac{1}{\lambda ^2}_0^\lambda sm_h(s)𝑑s+\frac{1}{\lambda }_\lambda ^{\mathrm{}}m_h(s)𝑑s\right\},0<\lambda <\mathrm{},$$
(1.1)
where
$$m_f(\lambda )=\text{meas}\left(\{|f|\lambda \}\right)=\nu _f(\{|w|\lambda \})=\mu _f(\lambda ^1),$$
and
$$m_h(\lambda )=\text{meas}\left(\{|h|\lambda \}\right)=\nu _f(\{|\text{Im}w|\lambda \})=𝔫_f(\lambda ^1),$$
reads:
$$\mu _f(r)C\left\{r_0^r\frac{𝔫_f(t)}{t^2}𝑑t+r^2_r^{\mathrm{}}\frac{𝔫_f(t)}{t^3}𝑑t\right\},0<r<\mathrm{}.$$
(1.2)
From this, one readily obtains
$$\mu _f(r)Cr_0^{\mathrm{}}\frac{𝔫_f(t)}{t^2}𝑑t,0<r<\mathrm{},$$
(1.3)
which is equivalent to Kolmogorov’s weak $`L^1`$-type inequality: $`\lambda m_f(\lambda )Ch_{L^1}`$, $`0<\lambda <\mathrm{}`$, and
$$_0^{\mathrm{}}\frac{\mu _f(t)}{t^{p+1}}𝑑tC(p)_0^{\mathrm{}}\frac{𝔫_f(t)}{t^{p+1}}𝑑t,1<p<2,$$
(1.4)
which is equivalent to M. Riesz’ inequality:
$$f_{L^p}C(p)h_{L^p}.$$
(1.5)
The classical proof of inequality (1.1) is based on the interpolation technique which later served as one of the cornerstones of an abstract theory of interpolation of operators in Banach spaces .
A natural question arises: what is special about the subharmonic function $`u_f(z)`$ which makes inequality (1.2) valid? The answer is proposed in : a key fact ensuring this, is the positivity condition:
$$u_f(z)0,z𝐂,$$
(1.6)
which can be easily checked with the help of Green’s formula (see \[3, Lemma 5\]) or by applying of the Cauchy residue theorem (see ).
This leads to a heuristic principle which suggests that
* results about the distribution of the Hilbert transform can be deduced from inequality (1.6) by using methods of the subharmonic function theory.
As will be shown, inequality (1.6) is even too strong, and in many cases it suffices to assume that $`u_f(x)0`$ on $`𝐑`$, or to control the negative part $`u_f^{}=\mathrm{max}(0,u_f)`$ on $`𝐑`$.
The principle shows a path to new results on the Hilbert transform. In , its application leads to a complete description of the distribution of the Hilbert transform of $`L^1`$-functions and measures of finite variations. At the same time, putting known estimates of the Hilbert transform into this setting, we arrive at new interesting questions about the argument-distribution of the Riesz measure in certain classes of subharmonic functions. For example, the proofs of the inequalities of Kolmogorov and M. Riesz found in give new bounds for zeros of polynomials. Positivity condition (1.6) links our work with the theory of uniform algebras and Jensen measures (see ).
In this work we put forward a new approach to the Marcinkiewicz inequality (1.1) (or (1.2)). The methods applied in , are too rigid for this. Here we use a different technique.
Here and later on, we use the following notations:
$`\varphi (s)\psi (s)`$ means that there is a positive numerical constant $`C`$ such that, for each $`s>0`$, $`\varphi (s)C\psi (s)`$;
$`\varphi (s)_\alpha \psi (s)`$ means the same as above but $`C`$ may depend on a parameter $`\alpha `$;
$`H(z)=\mathrm{log}|1z|+\text{Re}(z)`$ is the canonical kernel of genus one;
$`𝐂_\pm `$ are the upper and lower half-planes.
## 2 Main results
Define a subharmonic canonical integral of genus one:
$$u(z)=_𝐂H(z/\zeta )𝑑\mu (\zeta ),$$
(2.1)
where $`d\mu `$ is a non-negative locally finite measure on $`𝐂`$ such that
$$_𝐂\mathrm{min}(\frac{1}{|\zeta |},\frac{1}{|\zeta |^2})𝑑\mu (\zeta )<\mathrm{}.$$
(2.2)
Let
$$\mu (r)=\mu (\{|z|r\})$$
be a (conventional) counting function of the measure $`d\mu `$, and let
$$𝔫(r)=\mu \left(\{|\text{Im}(1/z)|1/r\}\right)=\mu \left(\{|zir/2|r/2\}\right)+\mu \left(\{|z+ir/2|r/2\}\right)$$
be its Levin-Tsuji counting function , (see also \[5, Chapter 1\]).
Let $`M(r,u)=\mathrm{max}_{|z|r}u(z)`$. Then by the Jensen inequality, $`\mu (r)M(r,u)`$. In the opposite direction, a standard estimate of the kernel
$$H(z)\frac{|z|^2}{1+|z|},z𝐂,$$
yields Borel’s estimate
$$M(r,u)r_0^r\frac{\mu (t)}{t^2}𝑑t+r^2_r^{\mathrm{}}\frac{\mu (t)}{t^3}𝑑t.$$
(2.3)
In particular,
$$M(r,u)=\{\begin{array}{cc}o(r),\hfill & r0\hfill \\ & \\ o(r^2),\hfill & r\mathrm{}.\hfill \end{array}$$
Theorem 1. Let $`u(z)`$ be a canonical integral (2.1) of genus one, then
$$M(r,u)r_0^r\frac{𝔫(t)}{t^2}𝑑t+r^2_r^{\mathrm{}}\frac{𝔫(t)}{t^3}𝑑t+r^2_r^{\mathrm{}}\frac{𝔪(t,u)}{t^2}𝑑t,$$
(2.4)
where
$$𝔪(r,u)=\frac{1}{2\pi }_0^{2\pi }u^{}(re^{i\theta }|\mathrm{sin}\theta )|)\frac{d\theta }{\mathrm{sin}^2\theta }$$
is the Tsuji proximity function of $`u`$.
If the function $`u`$ is non-negative in $`𝐂`$, then the proximity function $`𝔪(r,u)`$ vanishes, and applying Jensen’s inequality we arrive at
Corollary 1. Let $`u`$ be a canonical integral of genus one which is non-negative in $`𝐂`$. Then
$$\mu (r)r_0^r\frac{𝔫(t)}{t^2}𝑑t+r^2_r^{\mathrm{}}\frac{𝔫(t)}{t^3}𝑑t.$$
As we explained in the introduction, this result immediately yields the classical Marcinkiewicz inequality (1.1). In this case one can apriori assume that the function $`f`$ is bounded, so that most of the technicalities needed for the proof of Theorem 1 (see Lemmas 2 and 4 below) are redundant, and our proof of Marcinkiewicz’ inequality, being conceptually new, is comparable in length to the classical one.
There is a curious reformulation of Corollary 1. Let $``$ be a measurable space endowed with a locally finite non-negative measure $`dm`$, let $`f`$ be a complex-valued measurable function on $``$, and let $`m_f(\lambda )=m(\{|f|\lambda \})`$ be the distribution function of $`f`$. If
$$_{}\mathrm{min}(|f|,|f|^2)𝑑m<\mathrm{},$$
(2.5)
then we define the logarithmic determinant of $`f`$ of genus one
$$u_f(z)=_{}H(zf(t))𝑑m(t),$$
which is subharmonic in $`𝐂`$, and moreover is represented by a canonical integral of genus one. Applying Corollary 1, we obtain
Corollary 2. If $`f`$ satisfies condition (2.5), and the logarithmic determinant $`u_f`$ is non-negative in $`𝐂`$, then
$$m_f(\lambda )\left\{\frac{1}{\lambda ^2}_0^\lambda sm_{\mathrm{Im}f}(s)𝑑s+\frac{1}{\lambda }_\lambda ^{\mathrm{}}m_{\mathrm{Im}f}(s)𝑑s\right\},0<\lambda <\mathrm{}.$$
In particular,
$$f_{L^p(m)}_p\mathrm{Im}f_{L^p(m)},1<p<2,$$
and
$$m_f(\lambda )\frac{\mathrm{Im}f_{L^1(m)}}{\lambda },0<\lambda <\mathrm{}.$$
Corollary 1 can also be applied to Jensen measures in $`𝐂`$. A compactly supported finite measure $`\sigma `$ in $`𝐂`$ is called a Jensen measure (with respect to the origin) if for an arbitrary subharmonic function $`h`$ in $`𝐂`$
$$h(0)h𝑑\sigma .$$
(2.6)
A simple argument shows that (2.6) then holds true for subharmonic functions in a domain $`G`$ such that $`0G`$ and $`\mathrm{supp}(\sigma )\mathrm{G}`$. For a harmonic function, the equality sign must occur in (2.6). Therefore,
$$\sigma (𝐂)=1𝑑\sigma =1,$$
that is, $`\sigma `$ is a probability measure, and
$$\zeta ^k𝑑\sigma (\zeta )=0,k=1,2,\mathrm{}.$$
(2.7)
Define the potential
$$V_\sigma (z)=\mathrm{log}|1z\zeta |d\sigma (\zeta ).$$
(2.8)
Then, due to (2.6) and (2.7),
$$0V_\sigma (z)\mathrm{log}^+(c|z|),z𝐂,$$
(2.9)
for some $`c>0`$. The opposite is also true: if, for some $`c>0`$, a subharmonic function $`V`$ satisfies (2.9), then it is a potential of a of a Jensen measure \[6, §4\].
Due to condition (2.7), every potential $`V_\sigma `$ of a Jensen measure can be represented by a canonical integral of genus one:
$$V_\sigma (z)=_𝐂H(z/\zeta )𝑑\mu (\zeta ),d\mu (\zeta )=d\sigma (1/\zeta ).$$
Thus, Theorem 1 is applicable to the potential $`V_\sigma `$, and we obtain
Corollary 3. Let $`\sigma `$ be a Jensen measure in $`𝐂`$, $`\sigma (\lambda )=\sigma (|z|\lambda )`$, $`\sigma _I(\lambda )=\sigma (|\mathrm{Imz}|\lambda )`$. Then
$$\sigma (\lambda )\frac{1}{\lambda ^2}_0^\lambda s\sigma _I(s)𝑑s+\frac{1}{\lambda }_\lambda ^{\mathrm{}}\sigma _I(s)𝑑s.$$
The class of Jensen measures is invariant with respect to the holomorphic mappings. More precisely, let $`G`$ be a domain which contains the origin and $`\mathrm{supp}(\sigma )`$, and let $`F`$ be an analytic function in $`G`$, $`F(0)=0`$. Then the push forward $`F_{}\sigma `$ is defined by
$$\varphi 𝑑F_{}\sigma =\varphi F𝑑\sigma ,$$
where $`\varphi `$ is an arbitrary continuous function in $`𝐂`$, and $`\varphi F`$ is a composition of $`\varphi `$ and $`F`$. By the monotone convergence theorem this equation also holds for semicontinuous functions. The measure $`F_{}\sigma `$ automatically has a compact support since $`F`$ is bounded on $`\mathrm{supp}(\sigma )`$. If $`h`$ is subharmonic in $`𝐂`$, then $`hF`$ is subharmonic in $`G`$, and
$$h𝑑F_{}\sigma =hF𝑑\sigma h(F(0))=h(0).$$
Hence, $`F_{}\sigma `$ is a Jensen measure.
Corollary 4. Let $`\sigma `$ be a Jensen measure in $`𝐂`$, and let $`f=g+ih`$ be an analytic function in a domain $`G`$ which contains the origin and $`\mathrm{supp}(\sigma )`$, and $`f(0)=0`$. Let
$$m_{f,\sigma }(\lambda )=\sigma (|f|\lambda ),m_{h,\sigma }(\lambda )=\sigma (|h|\lambda ).$$
Then
$$m_{f,\sigma }(\lambda )\frac{1}{\lambda ^2}_0^\lambda sm_{h,\sigma }(s)𝑑s+\frac{1}{\lambda }_\lambda ^{\mathrm{}}m_{h,\sigma }(s)𝑑s.$$
(2.10)
Corollary 4 probably holds true under a weaker (and more natural) assumption $`g(0)=0`$ rather than $`f(0)=0`$. In that case, using Theorem 2 (see below) one can get an estimate which is slightly weaker than (2.10).
In the next result, we shall not assume that $`u(z)`$ is non-negative in $`𝐂`$ and instead introduce the quantity
$$\delta (r)=𝔫(r)+[u^{}(r)+u^{}(r)]$$
which we keep under control.
We assume that the integrals
$$_0\frac{\delta (t)}{t^2}𝑑t,\text{and}^{\mathrm{}}\frac{\delta (t)}{t^3}(1+\mathrm{log}t)𝑑t$$
(2.11)
are convergent and define
$$\delta ^{}(r)=r_0^r\frac{\delta (t)}{t^2}𝑑t+r^2_r^{\mathrm{}}\frac{\delta (t)}{t^3}\left(1+\mathrm{log}\frac{t}{r}\right)𝑑t$$
(2.12)
The function $`\delta ^{}(r)`$ does not decrease, $`r^2\delta ^{}(r)`$ does not increase, and therefore
$$\delta ^{}(r)\delta ^{}(2r)4\delta ^{}(r),0<r<\mathrm{}.$$
(2.13)
Theorem 2. Let $`u(z)`$ be an arbitrary subharmonic function in $`𝐂`$ represented by a canonical integral of genus one. Then
$$M(r,u)r^2\left[_r^{\mathrm{}}\frac{\sqrt{\delta ^{}(t)}}{t^2}𝑑t\right]^2.$$
(2.14)
Observe, that the RHSs of (2.4) and (2.14) do not depend on the bound for the integral (2.2). Estimate (2.14) is slightly weaker than (2.4); however, it suffices for deriving estimates of M. Riesz and Kolmogorov, as well as of the weak $`(p,\mathrm{})`$-type estimate (see Corollary 6 below).
Fix an arbitrary $`ϵ>0`$. Then by the Cauchy inequality
$`\left[{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{\sqrt{\delta ^{}(t)}}{t^2}}𝑑t\right]^2`$ $`=`$ $`\left[{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{\sqrt{\left(1+\mathrm{log}^{1+ϵ}\frac{t}{r}\right)\delta ^{}(t)}}{t^{3/2}}}{\displaystyle \frac{dt}{t^{1/2}\sqrt{1+\mathrm{log}^{1+ϵ}\frac{t}{r}}}}\right]^2`$
$`_ϵ`$ $`{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{\delta ^{}(t)}{t^3}}\left(1+\mathrm{log}^{1+ϵ}{\displaystyle \frac{t}{r}}\right)𝑑t`$
$`_ϵ`$ $`{\displaystyle \frac{1}{r}}{\displaystyle _0^r}{\displaystyle \frac{\delta (s)}{s^2}}𝑑s+{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{\delta (s)}{s^3}}\left(1+\mathrm{log}^{3+ϵ}{\displaystyle \frac{s}{r}}\right)𝑑s.`$
Thus we get
Corollary 5. For each $`ϵ>0`$,
$$M(r,u)_ϵr_0^r\frac{\delta (t)}{t^2}𝑑t+r^2_r^{\mathrm{}}\frac{\delta (t)}{t^3}\left(1+\mathrm{log}^{3+ϵ}\frac{t}{r}\right)𝑑t.$$
(2.15)
We do not know whether the term $`\mathrm{log}^{3+ϵ}`$ is really needed on the RHS of (2.15). Apparently, our method does not allow us to omit it. Rewriting (2.15) in the form
$$M(r,u)_ϵ_0^1\frac{\delta (rs)}{s^2}𝑑s+_1^{\mathrm{}}\frac{\delta (rs)}{s^3}\left(1+\mathrm{log}^{3+ϵ}s\right)𝑑s,$$
we immediately obtain
Corollary 6. The following inequalities hold for canonical integrals of genus one:
M. Riesz-type estimate:
$$_0^{\mathrm{}}\frac{\mu (r)}{r^{p+1}}𝑑r_p_0^{\mathrm{}}\frac{M(r,u)}{r^{p+1}}𝑑r_p_0^{\mathrm{}}\frac{\delta (r)}{r^{p+1}}𝑑r,1<p<2,$$
(2.16)
weak $`(p,\mathrm{})`$-type estimate:
$$\underset{r(0,\mathrm{})}{sup}\frac{\mu (r)}{r^p}_p\underset{r(0,\mathrm{})}{sup}\frac{M(r,u)}{r^p}_p\underset{r(0,\mathrm{})}{sup}\frac{\delta (r)}{r^p},1<p<2,$$
(2.17)
and Kolmogorov-type estimate:
$$\underset{r(0,\mathrm{})}{sup}\frac{\mu (r)}{r}\underset{r(0,\mathrm{})}{sup}\frac{M(r,u)}{r}_0^{\mathrm{}}\frac{\delta (r)}{r^2}𝑑r.$$
(2.18)
Estimates weaker than (2.16) and (2.18) were obtained in and under additional restrictions which now appear to be redundant. Estimate (2.17) is apparently new.
If we assume that $`d\mu `$ is supported by $`𝐑`$, that is, $`u(z)`$ is harmonic in $`𝐂_\pm `$, then our technique gives a better result:
Let $`u(z)`$ be a canonical integral of genus one of a measure $`d\mu `$ supported by $`𝐑`$. Then, for $`0<r<\mathrm{}`$,
$$M(r,u)r_0^r\frac{u^{}(t)+u^{}(t)}{t^2}𝑑t+r^2_r^{\mathrm{}}\frac{u^{}(t)+u^{}(t)}{t^3}𝑑t.$$
(2.19)
Note, that one cannot replace $`u^{}(x)`$ by $`u^+(x)`$ on the RHS of our estimates. For example, the function
$$u_p(z)=r^p\mathrm{cos}p\left(\frac{\pi }{2}|\theta |\right),1<p<2,$$
is subharmonic in $`𝐂`$, harmonic in the upper and lower half-planes $`𝐂_\pm `$, represented by the canonical integral of genus one of the measure $`d\mu (x)=c_p|x|^{p1}dx`$ ($`c_p>0`$), and non-positive on $`𝐑`$.
There is a corollary to Theorem 2 which is parallel to Corollary 2. Let $``$ be a measurable space endowed with a locally finite non-negative measure $`dm`$, and let $`f:𝐑^{n+1}`$, $`n1`$, be a measurable function such that
$$_{}\mathrm{min}(f,f^2)𝑑t<\mathrm{},$$
(2.20)
where $`||.||`$ stands for the $`n+1`$-dimensional Euclidean norm. We start to enumerate the coordinates in $`𝐑^{n+1}`$ with $`j=0`$, and denote by $`e_0`$ the vector in $`𝐑^{n+1}`$ with the zeroth coordinate equal one, and other coordinates vanishing. Let $`f_j(t)`$ be the $`j`$-th coordinate function of $`f(t)`$, and $`\widehat{f}(t)=\left\{_{j=1}^nf_j^2(t)\right\}^{1/2}`$. We define the logarithmic determinant
$`v_f(x)`$ $`=`$ $`{\displaystyle _{}}\left[\mathrm{log}e_0xf(t)+xf_0(t)\right]𝑑m(t)`$
$`=`$ $`{\displaystyle _{}}\left[\mathrm{log}\sqrt{12xf_0(t)+x^2f^2}+xf_0(t)\right]𝑑m(t),x𝐑,`$
where the integral converges due to assumption (2.20). Then, if the function $`v_f(x)`$ is non-negative on $`𝐑`$, we may estimate its distribution function $`m_f(\lambda )=m(\{f\lambda \})`$ by the distribution function $`m_{\widehat{f}}=m(\{\widehat{f}\lambda \})`$ of $`\widehat{f}`$.
For this, observe that
$$v_f(x)=_{}H(xf_C(t))𝑑m(t),x𝐑,$$
where $`f_C`$ is a “complex-valued surrogate” of $`f`$: $`f_C=f_0+i\widehat{f}`$. That is, $`v_f`$ has a subharmonic continuation from $`𝐑`$ to $`𝐂`$ by a canonical integral of genus one
$$u_{f_C}(z)=_{}H(zf_C(t))𝑑m(t),z𝐂.$$
Next, observe that $`m_f(\lambda )=m(\{f_0^2+\widehat{f}^2\lambda ^2\})=m_{f_C}(\lambda )`$, and $`m_{\widehat{f}}(\lambda )=m_{\mathrm{Im}f_C}(\lambda )`$ for $`0<\lambda <\mathrm{}`$. Hence, Theorem 2 is applicable in this situation. For simplicity, we restrict ourselves to the case when $`v_f`$ is non-negative on the real axis.
Corollary 7. Let $`f`$ satisfy condition (2.20), and let the logarithmic determinant $`v_f`$ be non-negative on the real axis. Then, for $`0<\lambda <\mathrm{}`$ and $`ϵ>0`$,
$$m_f(\lambda )_ϵ\frac{1}{\lambda ^2}_0^\lambda s\left(1+\mathrm{log}^{3+ϵ}\frac{\lambda }{s}\right)m_{\widehat{f}}(s)𝑑s+\frac{1}{\lambda }_\lambda ^{\mathrm{}}m_{\widehat{f}}(s)𝑑s.$$
In particular,
$$f_{L^p(m)}_p\widehat{f}_{L^p(m)},1<p<2,$$
and
$$m_f(\lambda )\frac{\widehat{f}_{L^1(m)}}{\lambda }.$$
This corollary may be of some interest in view of the results of Aleksandrov and Kargaev .
Our third result pertains to a more general class of subharmonic functions represented by a generalized canonical integral of genus one. It gives a Kolmogorov-type estimate which can be applied to a wider class of functions than (2.18):
Theorem 3. Let $`d\mu `$ be a non-negative locally finite measure on $`𝐂`$ such that
$$_{\{|\zeta |1\}}\frac{d\mu (\zeta )}{|\zeta |^2}<\mathrm{},$$
and let there exist a finite principal value integral
$$\underset{\epsilon 0}{lim}_{\{\epsilon |\zeta |1\}}\frac{d\mu (\zeta )}{\zeta }.$$
Let
$$u(z)=\underset{\epsilon 0}{lim}_{|\zeta |>\epsilon }H(z/\zeta )𝑑\mu (\zeta ),$$
then
$$\underset{0<r<\mathrm{}}{sup}\frac{M(r,u)}{r}_0^{\mathrm{}}\frac{\delta (t)}{t^2}𝑑t+\underset{r0}{lim\; sup}\frac{\mu (r)}{r}.$$
(2.21)
It is easy to see that if the integral (2.2) converges at the origin, then the upper limit on the RHS of (2.21) vanishes, and in this case (2.21) coincides with (2.18).
In fact, our proof yields a stronger result
$$_{\mathrm{}}^{\mathrm{}}\frac{u^+(t)}{t^2}𝑑t+\underset{r\mathrm{}}{lim\; sup}\frac{M(r,u)}{r}_0^{\mathrm{}}\frac{\delta (t)}{t^2}𝑑t+\underset{r0}{lim\; sup}\frac{\mu (r)}{r},$$
(2.22)
which gives control over the positive harmonic majorants of $`u`$ in the upper and lower half-planes. Applying a known technique of functions of Cartwright class , , one can extract from (2.22) information about the asymptotic regularity of $`u`$ and $`\mu `$ at infinity and near the origin.
Notice, that one can reformulate Theorem 3 in the spirit of Corollaries 2 and 7. We leave this to the reader.
## 3 Auxiliary Lemmas
We shall need several known facts about harmonic and subharmonic functions.
Lemma 1. Let $`v`$ be a subharmonic function in the angle $`S=\{z:\mathrm{\hspace{0.17em}0}<\mathrm{arg}z<\alpha \}`$, $`0<\alpha <2\pi `$, let
$$\underset{z\zeta ,zS}{lim\; sup}v^+(z)\mathrm{\Phi }(|\zeta |),\zeta S;$$
(3.1)
and let
$$_0^\alpha v^+(re^{i\theta })\mathrm{sin}\left(\frac{\pi }{\alpha }\theta \right)𝑑\theta =o(r^{\pi /\alpha }),r\mathrm{}.$$
(3.2)
Then, for $`z=re^{i\theta }S`$,
$$v(re^{i\theta })\mathrm{sin}\left(\frac{\pi }{\alpha }\theta \right)_\alpha r^{\pi /\alpha }_0^r\mathrm{\Phi }(t)t^{\pi /\alpha 1}𝑑t+r^{\pi /\alpha }_r^{\mathrm{}}\frac{\mathrm{\Phi }(t)}{t^{\pi /\alpha +1}}𝑑t.$$
(3.3)
If the majorant $`\mathrm{\Phi }(t)`$ does not decrease, then the factor $`\mathrm{sin}(\pi \theta /\alpha )`$ on the LHS of (3.3) can be omitted.
Proof: The general case is easily reduced to the special case when $`S=𝐂_+`$, so that, without loss of generality, we assume that $`\alpha =\pi `$. First, we show that $`v(z)`$ is majorized by the Poisson integral of $`\mathrm{\Phi }(|t|)`$, and then we estimate this integral.
Denote by $`h_R(z)`$ a harmonic function in the semi-disk $`\{\text{Im}z>0,|z|<R\}`$ with boundary values $`h_R(t)=\mathrm{\Phi }(|t|)`$, $`R<t<R`$, and $`h_R(Re^{i\theta })=v^+(Re^{i\theta })`$, $`0<\theta <\pi `$. Applying the Poisson-Nevanlinna representation in this semi-disk (see \[5, Chapter 1, Theorem 2.3\], \[10, Section 24.3\]), we obtain for $`z=re^{i\theta }`$, $`r<R`$,
$$v(z)h_R(z)=_R^R\mathrm{\Phi }(|t|)K_1(z,t)𝑑t+_0^\pi v^+(Re^{i\varphi })K_2(z,Re^{i\varphi })𝑑\varphi ,$$
(3.4)
where
$$K_1(z,t)=\frac{r\mathrm{sin}\theta }{\pi }\left\{\frac{1}{|zt|^2}\frac{R^2}{|R^2zt|^2}\right\},$$
(3.5)
$$K_2(z,Re^{i\varphi })=\frac{1}{2\pi }\frac{4Rr(R^2r^2)\mathrm{sin}\varphi \mathrm{sin}\theta }{(R^2+r^22Rr\mathrm{cos}(\varphi \theta ))(R^2+r^22Rr\mathrm{cos}(\varphi +\theta ))},$$
(3.6)
By condition (3.2), the second integral on the RHS of (3.4) tends to $`0`$ as $`R\mathrm{}`$. Therefore, letting $`R\mathrm{}`$ in (3.4), we obtain
$$v(z)\frac{r\mathrm{sin}\theta }{\pi }_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{\Phi }(|t|)}{|zt|^2}𝑑t.$$
(3.7)
Making use of straightforward estimates of the Poisson kernel, we get
$$v(z)\frac{2}{\pi r\mathrm{sin}\theta }_0^{2r}\mathrm{\Phi }(t)𝑑t+\frac{4r}{\pi }_{2r}^{\mathrm{}}\frac{\mathrm{\Phi }(t)}{t^2}𝑑t,$$
and estimate (3.3) follows.
If the majorant $`\mathrm{\Phi }(t)`$ does not decrease, then we modify the previous argument:
$$v(z)\frac{4}{r}_0^{r/2}\mathrm{\Phi }(t)dt+\mathrm{\Phi }(2r)+4r_{2r}^{\mathrm{}}\frac{\mathrm{\Phi }(t)}{t^2}dt\frac{1}{r}_0^r\mathrm{\Phi }(t)dt++r_r^{\mathrm{}}\frac{\mathrm{\Phi }(t)}{t^2}dt,$$
completing the proof. $`\mathrm{}`$
The next lemma asserts that under certain conditions the Carleman integral formula \[10, Lecture 24\], \[5, Chapter 1\] holds without remainder.
Lemma 2. Let $`v(z)`$ be a subharmonic function on $`D_R=\{z\overline{𝐂}_+:|z|R\}`$ which satisfies conditions
$$_0^\pi v^+(re^{i\theta })\mathrm{sin}\theta d\theta =o(r),r0,$$
(3.8)
and
$$_0\frac{\delta (t)}{t^2}𝑑t<\mathrm{}.$$
(3.9)
Then
$`{\displaystyle \frac{1}{2\pi }}{\displaystyle _R^R}v(t)\left({\displaystyle \frac{1}{t^2}}{\displaystyle \frac{1}{R^2}}\right)𝑑t`$ $`+`$ $`{\displaystyle \frac{1}{\pi R}}{\displaystyle _0^\pi }v(Re^{i\varphi })\mathrm{sin}\varphi d\varphi `$ (3.10)
$`=`$ $`{\displaystyle _{D_R}}\left({\displaystyle \frac{1}{|\zeta |^2}}{\displaystyle \frac{1}{R^2}}\right)\mathrm{Im}\zeta 𝑑\mu (\zeta ),`$
where the first integral on the LHS is absolutely convergent.
Proof: We start with the Nevanlinna representation
$`{\displaystyle _R^R}v^+(t)K_1(z,t)𝑑t`$ $`+`$ $`{\displaystyle _0^\pi }v^+(Re^{i\varphi })K_2(z,Re^{i\varphi })𝑑\varphi `$
$`=`$ $`v(z)+{\displaystyle _R^R}v^{}(t)K_1(z,t)𝑑t+{\displaystyle _0^\pi }v^{}(Re^{i\varphi })K_2(z,Re^{i\varphi })𝑑\varphi `$
$`+{\displaystyle _{D_R}}K_3(z,\zeta )𝑑\mu (\zeta ),`$
where the kernels $`K_1`$ and $`K_2`$ were defined by (3.5) and (3.6), and
$$K_3(z,\zeta )=\mathrm{log}\left|\frac{z\overline{\zeta }}{z\zeta }\frac{R^2z\overline{\zeta }}{R^2z\zeta }\right|.$$
(3.11)
We multiply both the left and right hand sides of the Nevanlinna representation by $`r^1\mathrm{sin}\theta `$, integrate it with respect to $`\theta `$ from $`0`$ to $`\pi `$ and change the integration order in all terms. We shall use the formulas:
$`{\displaystyle \frac{1}{r}}{\displaystyle _0^\pi }K_1(re^{i\theta },t)\mathrm{sin}\theta d\theta ={\displaystyle \frac{1}{2}}\left[\mathrm{min}\left({\displaystyle \frac{1}{t^2}}{\displaystyle \frac{1}{r^2}}\right){\displaystyle \frac{1}{R^2}}\right],`$ (3.12)
$`{\displaystyle \frac{1}{r}}{\displaystyle _0^\pi }K_2(re^{i\theta },Re^{i\varphi })\mathrm{sin}\theta d\theta ={\displaystyle \frac{1}{R}}\mathrm{sin}\varphi ,`$ (3.13)
and
$`{\displaystyle \frac{1}{r}}{\displaystyle _0^\pi }K_3(re^{i\theta },\zeta )\mathrm{sin}\theta d\theta =\pi \text{Im}\zeta \left[\mathrm{min}({\displaystyle \frac{1}{|\zeta |^2}},{\displaystyle \frac{1}{r^2}}){\displaystyle \frac{1}{R^2}}\right].`$ (3.14)
Observe that the RHS of relations (3.12)-(3.14) are non-decreasing functions of $`r^1`$. Therefore, making the limit transition $`r0`$, and using the monotone convergence theorem and condition (3.9) of the lemma, we get
$`{\displaystyle \frac{1}{2}}{\displaystyle _R^R}v^+(t)\left({\displaystyle \frac{1}{t^2}}{\displaystyle \frac{1}{R^2}}\right)𝑑t`$ $`+`$ $`{\displaystyle \frac{1}{R}}{\displaystyle _0^\pi }v^+(Re^{i\varphi })\mathrm{sin}\varphi d\varphi `$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _R^R}v^{}(t)\left({\displaystyle \frac{1}{t^2}}{\displaystyle \frac{1}{R^2}}\right)𝑑t+{\displaystyle \frac{1}{R}}{\displaystyle _0^\pi }v^{}(Re^{i\varphi })\mathrm{sin}\varphi d\varphi `$
$`+\pi {\displaystyle _{D_R}}\left({\displaystyle \frac{1}{|\zeta |^2}}{\displaystyle \frac{1}{R^2}}\right)\text{Im}\zeta 𝑑\mu (\zeta ).`$
The first and third integrals on the RHS are finite due to condition (3.9). This completes the proof. $`\mathrm{}`$
Remark. Condition (3.8) holds true for canonical integrals of genus one defined in (2.1).
Indeed, if $`u(z)`$ is such an integral, then due to (2.4)
$$_0^{2\pi }u^+(re^{i\theta })𝑑\theta =o(r),r0.$$
Since $`u(0)=0`$, this yields
$`{\displaystyle _0^{2\pi }}|u(re^{i\theta })|𝑑\theta `$ $`=`$ $`2{\displaystyle _0^{2\pi }}u^+(re^{i\theta })𝑑\theta {\displaystyle _0^{2\pi }}u(re^{i\theta })𝑑\theta `$
$`=`$ $`2{\displaystyle _0^{2\pi }}u^+(re^{i\theta })𝑑\theta =o(r),r0.`$
The third lemma was proved in (cf. \[10, Lecture 26\]). Its proof uses the Nevanlinna representation for the semi-disk.
Lemma 3. Let $`v(z)`$ be a function which is harmonic in $`𝐂_+`$, subharmonic in $`\overline{𝐂}_+`$, and satisfies conditions (3.8) and (3.9) of Lemma 2. Then, for $`z=re^{i\theta }𝐂_+`$,
$$v(re^{i\theta })\mathrm{sin}\theta \frac{1}{\pi }_0^\pi v^{}(2re^{i\phi })\mathrm{sin}\phi d\phi +\frac{r}{2\pi }_{2r}^{2r}\frac{v^{}(t)}{t^2}𝑑t.$$
(3.15)
The next lemma is a version of the Levin integral formula (, \[5, Chapter 1\]) without a remainder.
Lemma 4. Let $`v`$ be a subharmonic function in $`𝐂`$ such that $`v(z)`$ and $`v(\overline{z})`$ satisfy conditions (3.8) and (3.9) of Lemma 2. Then
$$\frac{1}{2\pi }_0^{2\pi }v(Re^{i\theta }|\mathrm{sin}\theta |)\frac{d\theta }{R\mathrm{sin}^2\theta }=_0^R\frac{𝔫(t)}{t^2}𝑑t,$$
(3.16)
where $`𝔫(t)`$ is the Levin-Tsuji counting function, and the integral on the LHS is absolutely convergent.
Proof: It suffices to prove that
$$\frac{1}{2\pi }_0^\pi v(Re^{i\theta }\mathrm{sin}\theta )\frac{d\theta }{R\mathrm{sin}^2\theta }=_{\left|\mathrm{Im}\frac{1}{\zeta }\right|>\frac{1}{R}}\left[\left|\mathrm{Im}\frac{1}{\zeta }\right|\frac{1}{R}\right]𝑑\mu (\zeta ).$$
(3.17)
Then (3.16) follows by adding to (3.17) a similar formula for the integral from $`\pi `$ to $`2\pi `$.
First, we prove that the integral on the LHS of relation (3.16) is absolutely convergent. Making use of notations introduced in (3.5), (6.3) and (3.11), observe that the Nevanlinna formula implies that
$$|v(z)|_R^R|v(t)|K_1(z,t)𝑑t+_0^\pi |v(Re^{i\varphi })|K_2(z,Re^{i\varphi })𝑑\varphi +_{D_R}K_3(z,\zeta )𝑑\mu (\zeta ).$$
We set $`z=Re^{i\theta }\mathrm{sin}\theta `$, multiply the formula by $`(R\mathrm{sin}^2\theta )^1`$, integrate it with respect to $`\varphi `$ from $`0`$ to $`\pi `$, and change the integration order in all terms. We shall use the following relations:
$$_0^\pi K_1(Re^{i\theta }\mathrm{sin}\theta ,t)\frac{d\theta }{R\mathrm{sin}^2\theta }=\frac{1}{t^2}\frac{1}{R^2},$$
$$_0^\pi K_2(Re^{i\theta }\mathrm{sin}\theta ,Re^{i\varphi })\frac{d\theta }{R\mathrm{sin}^2\theta }=\frac{2}{R}\mathrm{sin}\varphi ,$$
and
$$_0^\pi K_3(Re^{i\theta }\mathrm{sin}\theta ,\zeta )\frac{d\theta }{R\mathrm{sin}^2\theta }=2\pi \left[\mathrm{min}(\left|\text{Im}\frac{1}{\zeta }\right|,\frac{1}{R})\frac{\text{Im}\zeta }{R^2}\right].$$
Using these relations, we verify that
$`{\displaystyle _0^\pi }|v(Re^{i\theta }\mathrm{sin}\theta )|{\displaystyle \frac{d\theta }{R\mathrm{sin}^2\theta }}`$ $``$ $`{\displaystyle _R^R}|v(t)|\left({\displaystyle \frac{1}{t^2}}{\displaystyle \frac{1}{R^2}}\right)𝑑t`$
$`+`$ $`{\displaystyle \frac{2}{R}}{\displaystyle _0^\pi }|v(Re^{i\varphi })|\mathrm{sin}\varphi d\varphi +2\pi {\displaystyle _{D_R}}\left|\text{Im}{\displaystyle \frac{1}{\zeta }}\right|𝑑\mu (\zeta ).`$
The first integral on the RHS is finite due to Lemma 2, and the third is finite due to condition (3.9). That is, the integral on the LHS of (3.16) is absolutely convergent.
Now, we write the Nevanlinna formula in the form
$$v(z)=_R^Rv(t)K_1(z,t)𝑑t+_0^\pi v(Re^{i\varphi })K_2(z,Re^{i\varphi })𝑑\varphi _{D_R}K_3(z,\zeta )𝑑\mu (\zeta ).$$
Again, we set here $`z=Re^{i\theta }\mathrm{sin}\theta `$, multiply by $`(R\mathrm{sin}^2\theta )^1`$, integrate with respect to $`\theta `$ from $`0`$ to $`\pi `$ and change the integration order in all terms. We can do this since we already know that the integrals with $`|v|`$ instead of $`v`$ are finite. As a result, we obtain the equation
$`{\displaystyle _0^\pi }v(Re^{i\theta }\mathrm{sin}\theta ){\displaystyle \frac{d\theta }{R\mathrm{sin}^2\theta }}`$ $`=`$ $`{\displaystyle _R^R}v(t)\left({\displaystyle \frac{1}{t^2}}{\displaystyle \frac{1}{R^2}}\right)𝑑t+{\displaystyle \frac{2}{R}}{\displaystyle _0^\pi }v(Re^{i\varphi })\mathrm{sin}\varphi d\varphi `$
$`2\pi {\displaystyle _{D_R}}\left[\mathrm{min}(\left|\text{Im}{\displaystyle \frac{1}{\zeta }}\right|,{\displaystyle \frac{1}{R}}){\displaystyle \frac{\text{Im}\zeta }{R^2}}\right]𝑑\mu (\zeta ).`$
(3.18)
Taking into account (3.10), we get
$`{\displaystyle _0^\pi }v(Re^{i\varphi }\mathrm{sin}\varphi ){\displaystyle \frac{d\varphi }{R\mathrm{sin}^2\varphi }}`$ $`=`$ $`2\pi {\displaystyle _{D_R}}\left[{\displaystyle \frac{1}{|\zeta |^2}}{\displaystyle \frac{1}{R^2}}\right]\text{Im}\zeta 𝑑\mu (\zeta )`$
$`2\pi {\displaystyle _{D_R}}\left[\mathrm{min}(\left|\text{Im}{\displaystyle \frac{1}{\zeta }}\right|,{\displaystyle \frac{1}{R}}){\displaystyle \frac{\text{Im}\zeta }{R^2}}\right]𝑑\mu (\zeta )`$
$`=`$ $`2\pi {\displaystyle _{|\mathrm{Im}\frac{1}{\zeta }|\frac{1}{R}}}\left[\left|\text{Im}{\displaystyle \frac{1}{\zeta }}\right|{\displaystyle \frac{1}{R}}\right]𝑑\mu (\zeta ).`$
Then (3.17) follows and the proof is complete. $`\mathrm{}`$
In other words, in the assumptions of Lemma 4, the first fundamental theorem for Tsuji characteristics holds without a remainder term:
$$𝔗(r,u)=𝔪(r,u)+_0^r\frac{𝔫(t)}{t^2}𝑑t,0<r<\mathrm{},$$
(3.19)
where
$$𝔗(r,u)=\frac{1}{2\pi }_0^{2\pi }u^+(re^{i\theta }|\mathrm{sin}\theta )|)\frac{d\theta }{\mathrm{sin}^2\theta },$$
and
$$𝔪(r,u)=\frac{1}{2\pi }_0^{2\pi }u^{}(re^{i\theta }|\mathrm{sin}\theta )|)\frac{d\theta }{\mathrm{sin}^2\theta }.$$
The last lemma was proved in a slightly different setting in (see also \[5, Lemma 5.2, Chapter 6\]):
Lemma 5. Let $`u(z)`$ be a subharmonic function in $`𝐂`$, and let
$$T(r,u)=\frac{1}{2\pi }_0^{2\pi }u^+(re^{i\theta })𝑑\theta $$
be its Nevanlinna characteristic function. Then, for $`0<R<\mathrm{}`$,
$$_R^{\mathrm{}}\frac{T(r,u)}{r^3}𝑑r_R^{\mathrm{}}\frac{𝔗(r,u)}{r^2}𝑑r.$$
(3.20)
## 4 Proof of Theorem 1
Using monotonicity of $`T(r,u)`$, Lemma 5, and then Lemma 4, we obtain
$`{\displaystyle \frac{T(R,u)}{R^2}}`$ $``$ $`2{\displaystyle _R^{\mathrm{}}}{\displaystyle \frac{T(r,u)}{r^3}}𝑑r`$
$`\stackrel{(\text{3.20})}{}`$ $`2{\displaystyle _R^{\mathrm{}}}{\displaystyle \frac{𝔗(r,u)}{r^2}}𝑑r`$
$`\stackrel{(\text{3.19})}{=}`$ $`2{\displaystyle _R^{\mathrm{}}}{\displaystyle \frac{dr}{r^2}}\left({\displaystyle _0^r}{\displaystyle \frac{𝔫(t)}{t^2}}𝑑t+𝔪(r,u)\right)`$
$`=`$ $`{\displaystyle \frac{2}{R}}{\displaystyle _0^R}{\displaystyle \frac{𝔫(t)}{t^2}}𝑑t+2{\displaystyle _R^{\mathrm{}}}{\displaystyle \frac{𝔫(t)}{t^3}}𝑑t+2{\displaystyle _R^{\mathrm{}}}{\displaystyle \frac{𝔪(t,u)}{t^2}}𝑑t.`$
Then the inequality $`M(r,u)3T(2r,u)`$ completes the proof. $`\mathrm{}`$
## 5 Proof of Theorem 2
We split the proof into several parts. Without loss of generality, we assume convergence of the integrals
$$_0\frac{\delta (t)}{t^2}𝑑t\text{and}^{\mathrm{}}\frac{\delta (t)}{t^3}\mathrm{log}tdt.$$
We define a measure $`\mu _1`$, $`\mathrm{supp}(\mu _1)\overline{𝐂}_{}`$, by reflecting at the real axis the part of the measure $`\mu `$ which lies in the upper half-plane. Formally,
$$\mu _1(E)=\mu (E\overline{𝐂}_{})+\mu (E^{}\overline{𝐂}_{}),$$
where $`E𝐂`$ is a borelian set, and $`E^{}=\{z:\overline{z}E\}`$. Then the measure $`\mu _1`$ also satisfies condition (2.1) and we denote by $`u_1(z)`$ its canonical integral of genus one. Observe that $`u_1(t)=u(t)`$, so that $`\delta (t,u_1)=\delta (t,u)`$, $`t𝐑`$.
### 5.1 Estimate of $`u_1^{}(iy)`$, $`y>0`$.
We have
$`H(iy/\zeta )`$ $`=`$ $`\mathrm{log}\left|1+y\text{Im}{\displaystyle \frac{1}{\zeta }}iy\text{Re}{\displaystyle \frac{1}{\zeta }}\right|y\text{Im}{\displaystyle \frac{1}{\zeta }}`$
$``$ $`\mathrm{log}\left|1+y\text{Im}{\displaystyle \frac{1}{\zeta }}\right|y\text{Im}{\displaystyle \frac{1}{\zeta }}.`$
Since the RHS is non-positive for $`y>0`$ and $`\zeta \overline{𝐂}_{}`$,
$`u_1^{}(iy)`$ $``$ $`{\displaystyle _{\overline{𝐂}_{}}}\left[\mathrm{log}\left|1+y\text{Im}{\displaystyle \frac{1}{\zeta }}\right|y\text{Im}{\displaystyle \frac{1}{\zeta }}\right]𝑑\mu (\zeta )`$ (5.1)
$`=`$ $`{\displaystyle _0^{\mathrm{}}}\left[\mathrm{log}\left(1+{\displaystyle \frac{y}{t}}\right){\displaystyle \frac{y}{t}}\right]𝑑𝔫(t)`$
$`=`$ $`y^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{𝔫(t)}{t^2(t+y)}}𝑑t`$
$``$ $`y{\displaystyle _0^y}{\displaystyle \frac{𝔫(t)}{t^2}}𝑑t+y^2{\displaystyle _y^{\mathrm{}}}{\displaystyle \frac{𝔫(t)}{t^3}}𝑑t.`$
### 5.2 Estimates of $`u_1^+(re^{i\theta })`$, $`0<\theta <\pi `$.
Using harmonicity of the function $`u_1`$ in the upper half-plane, we transform the lower bound for $`u_1`$ into the upper bound. We shall show that
$$u_1^+(re^{i\theta })\mathrm{sin}\theta \delta ^{}(r)0<r<\mathrm{},0<\theta <\pi ,$$
(5.2)
where $`\delta ^{}(r)`$ is defined by (2.12).
Consider the function $`u_1(z)`$ and apply Lemma 1 to the angles $`\{0<\mathrm{arg}z<\pi /2\}`$ and $`\{\pi /2<\mathrm{arg}z<\pi \}`$ with
$$\mathrm{\Phi }(r)=[u_1^{}(r)+u_1^{}(r)]+r_0^r\frac{𝔫(t)}{t^2}𝑑t+r^2_r^{\mathrm{}}\frac{𝔫(t)}{t^3}𝑑t.$$
Condition (3.1) holds due to estimate (5.1), and condition (3.2) holds due to estimate (2) combined with Jensen’s inequality:
$$_0^\pi u_1^{}(re^{i\theta })𝑑\theta _0^{2\pi }u_1^+(re^{i\theta })𝑑\theta M(r,u_1)=o(r^2),r\mathrm{}.$$
Therefore,
$`u_1(re^{i\theta })|\mathrm{sin}2\theta |`$ $``$ $`{\displaystyle \frac{1}{r^2}}{\displaystyle _0^r}\mathrm{\Phi }(t)t𝑑t+r^2{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{\mathrm{\Phi }(t)}{t^3}}𝑑t`$ (5.3)
$``$ $`{\displaystyle \frac{1}{r^2}}{\displaystyle _0^r}[u_1^{}(t)+u_1^{}(t)]t𝑑t+r^2{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{u_1^{}(t)+u_1^{}(t)}{t^3}}𝑑t`$
$`+r{\displaystyle _0^r}{\displaystyle \frac{𝔫(s)}{s^2}}𝑑s+r^2{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{𝔫(s)}{s^3}}\left(1+\mathrm{log}{\displaystyle \frac{s}{r}}\right)𝑑s`$
$``$ $`\delta ^{}(r).`$
Observe that the factor $`|\mathrm{sin}2\theta |`$ on the LHS of (5.3) can be replaced by $`\mathrm{sin}\theta `$. This follows from inspection of the proof of Lemma 1 (since on the imaginary axis the function $`u(iy)`$ has an increasing majorant). Alternatively, one may again apply Lemma 1 to a small angle around the imaginary axis, say in $`\{|\theta \pi /2|<\pi /8\}`$. That is, we have
$$u_1(re^{i\theta })\mathrm{sin}\theta \delta ^{}(r).$$
(5.4)
Using Lemma 3 we obtain
$$u_1^+(re^{i\theta })\mathrm{sin}\theta _0^\pi u_1^{}(2re^{i\varphi })\mathrm{sin}\varphi d\varphi +r_0^{2r}\frac{u_1^{}(t)+u_1^{}(t)}{t^2}𝑑t\delta ^{}(r),$$
proving estimate (5.2).
### 5.3 Estimate of $`u^+(re^{i\theta })`$, $`\theta 0,\pi `$.
Here we prove that, for an arbitrary $`\eta >0`$,
$`u^+(re^{i\theta }){\displaystyle \frac{\delta ^{}(r)}{\eta |\mathrm{sin}\theta |}}+{\displaystyle \frac{\eta r^2}{\mathrm{sin}^2\theta }}{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{M(t,u)}{t^3}}𝑑t.`$ (5.5)
For this, we shall need several upper bounds for the difference
$$D=D(z,\zeta )=H(z/\zeta )H(z/\overline{\zeta })=\mathrm{log}\left|\frac{1z/\zeta }{1z/\overline{\zeta }}\right|+\mathrm{Re}\left[z\left(\frac{1}{\zeta }\frac{1}{\overline{\zeta }}\right)\right],$$
when $`z,\zeta \overline{𝐂}_+`$.
First,
$$D=\mathrm{log}\left|\frac{z\zeta }{z\overline{\zeta }}\right|+2\mathrm{I}\mathrm{m}z\left|\mathrm{Im}\frac{1}{\zeta }\right|2|z|\left|\mathrm{Im}\frac{1}{\zeta }\right|.$$
(5.6)
We shall use this estimate when $`|z|\left|\mathrm{Im}\frac{1}{\zeta }\right|1`$.
Next, let $`t=|z|/|\zeta |`$, $`\theta =\mathrm{arg}(z)`$, $`\varphi =\mathrm{arg}(\zeta )`$. Then
$`D`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}\left[1{\displaystyle \frac{4t\mathrm{sin}\theta \mathrm{sin}\varphi }{|1te^{i(\theta +\varphi )}|^2}}\right]+2t\mathrm{sin}\theta \mathrm{sin}\varphi `$ (5.7)
$``$ $`{\displaystyle \frac{2t\mathrm{sin}\theta \mathrm{sin}\varphi }{|1te^{i(\theta +\varphi )}|^2}}+2t\mathrm{sin}\theta \mathrm{sin}\varphi `$
$`=`$ $`2t\mathrm{sin}\theta \mathrm{sin}\varphi {\displaystyle \frac{2t\mathrm{cos}(\theta +\varphi )+t^2}{|1te^{i(\theta +\varphi )}|^2}}`$
$``$ $`t\mathrm{sin}\theta \mathrm{sin}\varphi {\displaystyle \frac{\mathrm{max}(t,t^2)}{|1te^{i(\theta +\varphi )}|^2}}.`$
If $`t1/2`$, then
$$|1te^{i(\theta +\varphi )}|^21,$$
and we obtain
$$Dt^2\mathrm{sin}\theta \mathrm{sin}\varphi \eta t^2+\eta ^1t^2\mathrm{sin}^2\varphi =\eta \frac{|z|^2}{|\zeta |^2}+\frac{|z|^2}{\eta }\left|\mathrm{Im}\frac{1}{\zeta }\right|^2,$$
(5.8)
with an arbitrary $`\eta >0`$.
If $`t1/2`$, then
$$|1te^{i(\theta +\varphi )}|^2t^2\mathrm{sin}^2\theta ,$$
so that (5.7) gives us
$$Dt\frac{\mathrm{sin}\varphi }{\mathrm{sin}\theta }\frac{\eta }{\mathrm{sin}^2\theta }+\frac{t^2}{\eta }\mathrm{sin}^2\varphi =\frac{\eta }{\mathrm{sin}^2\theta }+\frac{|z|^2}{\eta }\left|\mathrm{Im}\frac{1}{\zeta }\right|^2,$$
(5.9)
again, with an arbitrary positive $`\eta `$. We shall use the bounds (5.8) and (5.9) when $`|z|\left|\mathrm{Im}\frac{1}{\zeta }\right|1`$.
Now, for $`z𝐂_+`$, $`r=|z|`$, we have
$`u(z)u_1(z)`$ $`=`$ $`{\displaystyle _{𝐂_+}}D(z,\zeta )𝑑\mu (\zeta )`$
$``$ $`\left({\displaystyle _{\left|\mathrm{Im}\frac{1}{\zeta }\right|\frac{1}{r}}}+{\displaystyle _{\left|\mathrm{Im}\frac{1}{\zeta }\right|\frac{1}{r},|\zeta |2r}}+{\displaystyle _{\left|\mathrm{Im}\frac{1}{\zeta }\right|\frac{1}{r},|\zeta |2r}}\right)D(z,\zeta )d\mu (\zeta )`$
$``$ $`r{\displaystyle _0^r}{\displaystyle \frac{d𝔫(t)}{t}}+{\displaystyle \frac{r^2}{\eta }}{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{d𝔫(t)}{t^2}}+{\displaystyle \frac{\eta }{\mathrm{sin}^2\theta }}{\displaystyle _0^r}𝑑\mu (t)+\eta r^2{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{d\mu (t)}{t^2}}`$
$``$ $`{\displaystyle \frac{\delta (r)}{\eta }}+{\displaystyle \frac{\eta r^2}{\mathrm{sin}^2\theta }}{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{M(t,u)}{t^3}}𝑑t.`$
Then, using estimate (5.2) for $`u_1^+(z)`$ in the upper half-plane, we obtain estimate (5.5) for $`0<\theta <\pi `$. The same argument applies for the lower half-plane, and the proof of (5.5) is complete.
### 5.4 Integral inequality for $`M(r,u)`$.
Here we prove the integral inequality
$$M(r,u)\sqrt{\delta ^{}(r)r^2_r^{\mathrm{}}\frac{M(t,u)}{t^3}𝑑t}.$$
(5.10)
First, we improve estimate (5.5) near the real axis. Consider the function $`u(z)`$ in the angles $`\{|\mathrm{arg}z|\pi /6\}`$ and $`\{|\mathrm{arg}z\pi |\pi /6\}`$. On the boundary of these angles,
$$u(re^{i\theta })\mathrm{\Phi }(r),\theta =\pm \frac{\pi }{6},\pi \pm \frac{\pi }{6},$$
where
$$\mathrm{\Phi }(r)=\eta ^1\delta ^{}(r)+\eta r^2_r^{\mathrm{}}\frac{M(t,u)}{t^3}𝑑t.$$
Applying Lemma 1 to $`u(z)`$ in these angles, we obtain for $`|\theta |\pi /8`$ and $`|\pi \theta |\pi /8`$,
$$u(re^{i\theta })r^3_0^r\mathrm{\Phi }(t)t^2𝑑t+r^3_r^{\mathrm{}}\frac{\mathrm{\Phi }(t)}{t^4}𝑑t\mathrm{\Phi }(r).$$
The second inequality follows since the function $`\mathrm{\Phi }(r)`$ does not decrease, and the function $`r^2\mathrm{\Phi }(r)`$ does not increase.
Thus, for $`0<r<\mathrm{}`$,
$$M(r,u)\mathrm{\Phi }(r)=\eta ^1\delta ^{}(r)+\eta r^2_r^{\mathrm{}}\frac{M(t,u)}{t^3}𝑑t.$$
Choosing
$$\eta =\sqrt{\delta ^{}(r)}:\sqrt{r^2_r^{\mathrm{}}\frac{M(t,u)}{t^3}𝑑t},$$
we obtain inequality (5.10).
### 5.5 Solution of the integral inequality (5.10).
We set
$$M_1(r)=_r^{\mathrm{}}\frac{M(t,u)}{t^3}𝑑t.$$
Then
$$M(r,u)=r^3M_1^{}(r),$$
and inequality (5.10) takes the form
$$M_1^{}(r)r^2\sqrt{\delta ^{}(r)M_1(r)}$$
or
$$\frac{d\sqrt{M_1(r)}}{dr}\frac{\sqrt{\delta ^{}(r)}}{r^2}.$$
Integrating this inequality from $`\mathrm{}`$ to $`r`$, we obtain
$$M_1(r)\left[_r^{\mathrm{}}\frac{\sqrt{\delta ^{}(t)}}{t^2}𝑑t\right]^2.$$
On the other hand, since $`M(r,u)`$ does not decrease,
$$M_1(r)M(r,u)_r^{\mathrm{}}\frac{dt}{t^3}=\frac{M(r,u)}{2r^2}.$$
Therefore,
$$M(r,u)r^2M_1(r)r^2\left[_r^{\mathrm{}}\frac{\sqrt{\delta ^{}(t)}}{t^2}𝑑t\right]^2,$$
completing the proof of Theorem 2. $`\mathrm{}`$
## 6 Proof of Theorem 3
We divide the proof into 4 parts. Set
$$B:=\underset{r0}{lim\; sup}\frac{\mu (r)}{r},$$
$$C:=_0^{\mathrm{}}\frac{\delta (t)}{t^2}𝑑t.$$
Without loss of generality, we assume that both values $`B`$ and $`C`$ are finite.
First, we shall prove the theorem under the additional assumption
$$\mathrm{supp}(\mu )\overline{𝐂}_{},$$
(6.1)
and till Section 6.4 we assume that the function $`u(z)`$ is harmonic in $`𝐂_+`$.
### 6.1 The function $`u(z)`$ has nonnegative harmonic majorants in $`𝐂_\pm `$.
Consider the function
$$U(z):=u(z)\frac{y}{\pi }_{\mathrm{}}^{\mathrm{}}\frac{u^{}(t)dt}{(tx)^2+y^2}.$$
This function is harmonic in $`𝐂_+`$ and $`U(x)0,x𝐑`$. Moreover, for $`y>0`$,
$`U(iy)u(iy)`$ $`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle _{|\zeta |\epsilon ,\zeta \overline{𝐂}_{}}}\left[\mathrm{log}\left|1{\displaystyle \frac{iy}{\zeta }}\right|+\text{Re}{\displaystyle \frac{iy}{\zeta }}\right]𝑑\mu (\zeta )`$
$``$ $`\underset{\epsilon 0}{lim}{\displaystyle _{|\zeta |\epsilon ,\zeta \overline{𝐂}_{}}}\text{Re}{\displaystyle \frac{iy}{\zeta }}𝑑\mu (\zeta )`$
$`=`$ $`y{\displaystyle _{\overline{𝐂}_{}}}\text{Im}{\displaystyle \frac{1}{\zeta }}𝑑\mu (\zeta )Cy.`$
By the Poisson-Nevanlinna representation of harmonic functions in the semi-disk $`D_R`$ (cf. Section 3), we have
$`U(z)`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^\pi }U(Re^{i\varphi })K_2(z,Re^{i\varphi })𝑑\varphi `$ (6.2)
$``$ $`{\displaystyle \frac{2Rr(R+r)}{\pi (Rr)^3}}{\displaystyle _0^\pi }u^{}(Re^{i\varphi })𝑑\varphi `$
$``$ $`{\displaystyle \frac{4Rr(R+r)}{(Rr)^3}}T(R,u),`$
where $`T(R,u)`$ is the Nevanlinna characteristic of $`u`$.
Note that, for any $`\delta >0`$, the function $`u`$ can be represented in the form
$`u(z)`$ $`=`$ $`{\displaystyle _{|\zeta |>\delta }}\left[\mathrm{log}\left|1{\displaystyle \frac{z}{\zeta }}\right|+\text{Re}{\displaystyle \frac{z}{\zeta }}\right]𝑑\mu (\zeta )`$ (6.3)
$`+{\displaystyle _{|\zeta |\delta }}\mathrm{log}\left|1{\displaystyle \frac{z}{\zeta }}\right|d\mu (\zeta )+\text{Re}\left(z{\displaystyle _{|\zeta |<\delta }}{\displaystyle \frac{d\mu (\zeta )}{\zeta }}\right)`$
$`=:`$ $`u^\delta (z)+v^\delta (z)+\text{Re}\left(z{\displaystyle _{|\zeta |<\delta }}{\displaystyle \frac{d\mu (\zeta )}{\zeta }}\right).`$
The well-known Borel estimates
$$\underset{|z|r}{\mathrm{max}}u^\delta (z)=o(|z|^2),\underset{|z|r}{\mathrm{max}}v^\delta (z)=o(|z|),z\mathrm{},$$
imply that $`T(R,u)=o(R^2),R\mathrm{}`$. Therefore, by setting $`R=2r`$ in (6.2), we get
$$U^+(z)=o(|z|^2),z\mathrm{},\text{Im}z>0.$$
Applying the Phragmén-Lindelöf principle in the angles $`\{0<\mathrm{arg}z<\pi /2\}`$ and $`\{\pi /2<\mathrm{arg}z<\pi \}`$, we conclude that
$$U(z)Cy,z=x+iy𝐂_+;$$
(6.4)
i.e. $`U(z)+Cy`$ is a nonnegative harmonic function in $`𝐂_+`$. Since $`u(z)U(z)+Cy`$, the function $`u(z)`$ also has a nonnegative harmonic majorant in $`𝐂_+`$.
For $`z𝐂_{}`$, we write
$`u(z)u(\overline{z})`$ $`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle _{|\zeta |\epsilon ,\zeta 𝐂_{}}}\mathrm{log}\left|{\displaystyle \frac{1z/\zeta }{1z/\overline{\zeta }}}\right|+\text{Re}\left[z\left({\displaystyle \frac{1}{\zeta }}{\displaystyle \frac{1}{\overline{\zeta }}}\right)\right]d\mu (\zeta )`$ (6.5)
$``$ $`{\displaystyle _𝐂_{}}\text{Re}\left[z2i\text{Im}{\displaystyle \frac{1}{\zeta }}\right]𝑑\mu (\zeta )2C|y|.`$
Because $`u(\overline{z})`$ has a nonnegative harmonic majorant in $`𝐂_{}`$, we get the desired conclusion.
### 6.2 Estimate of $`u(z)`$ near the origin.
Set
$$I(r):=\frac{1}{r}_0^\pi u(re^{i\theta })\mathrm{sin}\theta d\theta .$$
Let us prove that
$$\underset{r0}{lim\; sup}I(r)B.$$
(6.6)
For any given $`\epsilon >0`$, choose a positive $`\delta <\epsilon `$ such that
$$\mu (r)<(B+\epsilon )r,\text{for}0<r<\delta .$$
Let us represent $`u`$ by the formula (6.3) with this $`\delta `$. Since
$$u^\delta (z)=O(|z|^2),z0,$$
we have
$$\underset{r0}{lim\; sup}I(r)\underset{r0}{lim\; sup}I^\delta (r)+\left|_{|\zeta |<\delta }\frac{d\mu (\zeta )}{\zeta }\right|,$$
(6.7)
where
$$I^\delta (r)=\frac{1}{r}_0^\pi v^\delta (re^{i\theta })\mathrm{sin}\theta d\theta .$$
It suffices to show that
$$I^\delta (r)B+\epsilon +_{|\zeta |<\delta }\text{Im}\frac{1}{\zeta }𝑑\mu (\zeta ),0<r<\delta .$$
(6.8)
Indeed, if (6.8) is valid, then substituting it into (6.7), we get
$$\underset{r0}{lim\; sup}I(r)B+\epsilon +_{|\zeta |<\delta }\text{Im}\frac{1}{\zeta }𝑑\mu (\zeta )+\left|_{|\zeta |<\delta }\frac{d\mu (\zeta )}{\zeta }\right|.$$
Taking the limit as $`\epsilon 0`$ (then $`\delta 0`$ as well), we obtain (6.6).
To prove (6.8), we set for $`|z|=r,\mathrm{\hspace{0.33em}0}<r<\delta `$:
$$v^\delta (z)=_{r<|\zeta |<\delta }+_{|\zeta |<r}=:v_1^\delta (z)+v_2^\delta (z),$$
and
$$I_j^\delta (r):=\frac{1}{r}_0^\pi v_j^\delta (re^{i\theta })\mathrm{sin}\theta d\theta ,j=1,2.$$
Note that
$`{\displaystyle _0^\pi }\left|\mathrm{log}\left|1{\displaystyle \frac{re^{i\theta }}{\zeta }}\right|\right|\mathrm{sin}\theta d\theta `$ $``$ $`2{\displaystyle _0^{2\pi }}\mathrm{log}^+\left|1{\displaystyle \frac{re^{i\theta }}{\zeta }}\right|d\theta {\displaystyle _0^{2\pi }}\mathrm{log}\left|1{\displaystyle \frac{re^{i\theta }}{\zeta }}\right|d\theta `$
$``$ $`4\pi \mathrm{log}\left(1+{\displaystyle \frac{\delta }{|\zeta |}}\right)+2\pi \mathrm{log}^+{\displaystyle \frac{r}{|\zeta |}}.`$
This estimate will allow us to change the integration order in the double integrals that arise when estimating $`I_j^\delta (r),j=1,2,`$ below.
Write
$$I_1^\delta (r)=\frac{1}{r}_{r<|\zeta |<\delta }𝑑\mu (\zeta )_0^\pi \mathrm{log}\left|1\frac{re^{i\theta }}{\zeta }\right|\mathrm{sin}\theta d\theta .$$
For $`r<|\zeta |`$, we have
$`{\displaystyle _0^\pi }\mathrm{log}\left|1{\displaystyle \frac{re^{i\theta }}{\zeta }}\right|\mathrm{sin}\theta d\theta `$ $`=`$ $`\text{Re}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{r^k}{k\zeta ^k}}{\displaystyle _0^\pi }e^{ik\theta }\mathrm{sin}\theta d\theta `$
$`=`$ $`{\displaystyle \frac{\pi r}{2}}\text{Im}{\displaystyle \frac{1}{\zeta }}+\text{Re}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{r^{2m}}{m(4m^21)\zeta ^{2m}}}.`$
Hence
$$I_1^\delta (r)\frac{\pi }{2}_{|\zeta |<\delta }\text{Im}\frac{1}{\zeta }𝑑\mu (\zeta )+\underset{m=1}{\overset{\mathrm{}}{}}\frac{r^{2m1}}{m(4m^21)}_{r<|\zeta |<\delta }\frac{d\mu (\zeta )}{|\zeta |^{2m}}.$$
Since
$`{\displaystyle _{r<|\zeta |<\delta }}{\displaystyle \frac{d\mu (\zeta )}{|\zeta |^{2m}}}`$ $`=`$ $`{\displaystyle _r^\delta }{\displaystyle \frac{d\mu (t)}{t^{2m}}}`$
$``$ $`{\displaystyle \frac{\mu (\delta )}{\delta ^{2m}}}+2m{\displaystyle _r^\delta }{\displaystyle \frac{\mu (t)}{t^{2m+1}}}𝑑t`$
$``$ $`(B+\epsilon )\delta ^{2m+1}+2(B+\epsilon )r^{2m+1}<3(B+\epsilon )r^{2m+1},`$
we get
$$I_1^\delta (r)\frac{\pi }{2}_{|\zeta |<\delta }\text{Im}\frac{1}{\zeta }𝑑\mu (\zeta )+3(B+\epsilon )\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m(4m^21)}.$$
(6.9)
Further,
$$I_2^\delta (r)=\frac{1}{r}_{|\zeta |<r}𝑑\mu (\zeta )_0^\pi \mathrm{log}\left|1\frac{re^{i\theta }}{\zeta }\right|\mathrm{sin}\theta d\theta .$$
For $`|\zeta |r`$, we have
$`{\displaystyle _0^\pi }\mathrm{log}\left|1{\displaystyle \frac{re^{i\theta }}{\zeta }}\right|\mathrm{sin}\theta d\theta `$ $`=`$ $`2\mathrm{log}{\displaystyle \frac{r}{|\zeta |}}+{\displaystyle _0^\pi }\mathrm{log}\left|1{\displaystyle \frac{\overline{\zeta }e^{i\theta }}{r}}\right|\mathrm{sin}\theta d\theta `$
$`=`$ $`2\mathrm{log}{\displaystyle \frac{r}{|\zeta |}}\text{Re}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\overline{\zeta }^k}{kr^k}}{\displaystyle _0^\pi }e^{ik\theta }\mathrm{sin}\theta d\theta `$
$`=`$ $`2\mathrm{log}{\displaystyle \frac{r}{|\zeta |}}{\displaystyle \frac{\pi }{2r}}\text{Im}\zeta +\text{Re}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\overline{\zeta }^{2m}}{m(4m^21)r^{2m}}}.`$
Whence
$`I_2^\delta (r)`$ $``$ $`{\displaystyle \frac{1}{r}}{\displaystyle _{|\zeta |<r}}\left[2\mathrm{log}{\displaystyle \frac{r}{|\zeta |}}+{\displaystyle \frac{\pi }{2}}+{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m(4m^21)}}\right]𝑑\mu (\zeta )`$ (6.10)
$``$ $`{\displaystyle \frac{1}{r}}{\displaystyle _0^r}\mathrm{log}{\displaystyle \frac{r}{t}}d\mu (t)+{\displaystyle \frac{\mu (r)}{r}}`$
$``$ $`B+\epsilon `$
Since $`I^\delta =I_1^\delta +I_2^\delta `$, the desired inequality (6.8) follows from (6.9) and (6.10).
### 6.3 Estimate of $`u^+(z)`$ on the real and imaginary axes.
Let us prove that
$$_{\mathrm{}}^{\mathrm{}}\frac{u^+(t)}{t^2}𝑑t+\underset{y+\mathrm{}}{lim\; sup}\frac{u^+(iy)}{y}B+C.$$
(6.11)
Since $`u`$ has a nonnegative harmonic majorant in $`𝐂_+`$, we have
$$_{\mathrm{}}^{\mathrm{}}\frac{|u(t)|dt}{1+t^2}<\mathrm{},$$
and $`u`$ admits the Poisson representation
$$u(re^{i\phi })=\frac{r\mathrm{sin}\phi }{\pi }_{\mathrm{}}^{\mathrm{}}\frac{u(t)dt}{r^2+t^22rt\mathrm{cos}\phi }+kr\mathrm{sin}\phi ,$$
(6.12)
where
$$k=\underset{y+\mathrm{}}{lim\; sup}\frac{u(iy)}{y}\mathrm{}.$$
Note that inequality (6.4) implies
$$u(iy)\frac{y}{\pi }_{\mathrm{}}^{\mathrm{}}\frac{u^{}(t)dt}{t^2+y^2}Cy=o(y)Cy,y\mathrm{};$$
i.e. $`kC`$.
Multiplying (6.12) by $`\mathrm{sin}\phi `$, integrating against $`\phi `$ from $`0`$ to $`\pi `$, and taking into account that
$$_0^\pi \frac{\mathrm{sin}^2\phi d\phi }{r^2+t^22rt\mathrm{cos}\phi }=\frac{\pi }{2}\mathrm{min}(\frac{1}{r^2},\frac{1}{t^2}),$$
we get
$$_0^\pi u(re^{i\phi })\mathrm{sin}\phi d\phi =\frac{r}{2}_{\mathrm{}}^{\mathrm{}}u(t)\mathrm{min}(\frac{1}{r^2},\frac{1}{t^2})𝑑t+\frac{k\pi r}{2}.$$
Hence
$$_{\mathrm{}}^{\mathrm{}}u^+(t)\mathrm{min}(\frac{1}{r^2},\frac{1}{t^2})𝑑t+k^+\pi =_{\mathrm{}}^{\mathrm{}}u^{}(t)\mathrm{min}(\frac{1}{r^2},\frac{1}{t^2})𝑑t+k^{}\pi +2I(r).$$
Letting $`r0`$, we obtain by the monotone convergence theorem
$$_{\mathrm{}}^{\mathrm{}}\frac{u^+(t)}{t^2}𝑑t+k^+\pi =_{\mathrm{}}^{\mathrm{}}\frac{u^{}(t)}{t^2}𝑑t+k^{}\pi +2\underset{r0}{lim}I(r)$$
(it turns out that the last limit exists). Taking into account that $`k^{}C`$ and using (6.6), we get (6.11).
### 6.4 Concluding steps.
From (6.4) and (6.11) we obtain
$$\underset{y\mathrm{}}{lim\; sup}\frac{u^+(iy)}{|y|}\underset{y+\mathrm{}}{lim\; sup}\frac{u^+(iy)+2Cy}{y}B+C.$$
Since $`u`$ has nonnegative harmonic majorants in both upper and lower half-planes, the following inequality holds in the whole plane:
$$u(z)\frac{|y|}{\pi }_{\mathrm{}}^{\mathrm{}}\frac{u^+(t)dt}{(tx)^2+y^2}+(B+C)|y|,z=x+iy.$$
(6.13)
The assertion of Theorem 3 can be obtained from this inequality and (6.11) by applying a known argument (cf. , ). First, one applies (6.13) and (6.11) to get the upper bound for $`u(z)`$ in the angles $`\{|\mathrm{arg}z\pm \pi /2|\pi /4\}`$, and then, using the Phragmén-Lindelöf principle, one obtains the upper bound for $`u(z)`$ in the complementary angles. This gives $`u(z)(B+C)|z|`$, and completes the proof of estimate (2.21) for the special case (6.1).
Now, let $`\mu `$ be an arbitrary measure in $`𝐂`$ satisfying conditions of Theorem 3 and having finite value $`C`$. As in Section 5, we define the measure $`\mu _1`$, $`\mathrm{supp}\mu _1\overline{𝐂}_{}`$, by reflecting with respect to the real axis the part of $`\mu `$ which charges $`𝐂_+`$. Since
$$_𝐂_{}\mathrm{Im}\frac{1}{\zeta }𝑑\mu _1(\zeta )C<\mathrm{},$$
the measure $`\mu _1`$ also satisfies the conditions of Theorem 3, and we can define the corresponding generalized canonical integral $`u_1(z)`$ of this measure. Then a straightforward estimate (cf. 6.5) shows that for $`z\overline{𝐂}_+`$
$$u(z)u_1(z)+2|y|_𝐂\left|\mathrm{Im}\frac{1}{\zeta }\right|𝑑\mu (\zeta )(B+C)|z|.$$
The same estimate holds in the lower half-plane, and the general case of Theorem 3 follows. $`\mathrm{}`$
Vladimir Matsaev: School of Mathematical Sciences, Tel-Aviv University,
Ramat-Aviv, 69978, Israel
matsaev@math.tau.ac.il
Iossif Ostrovskii: Department of Mathematics, Bilkent University,
06533 Bilkent, Ankara, Turkey
iossif@fen.bilkent.edu.tr
and Verkin Institute for Low Temperature Physics and Engineering,
310164 Kharkov, Ukraine
ostrovskii@ilt.kharkov.ua
Mikhail Sodin: School of Mathematical Sciences, Tel-Aviv University,
Ramat-Aviv, 69978, Israel
sodin@math.tau.ac.il
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# A GRAVITATIONAL LENS SURVEY WITH THE PLANCK SURVEYOR
## 1 Introduction
The Planck Surveyor CMB imaging satellite$`^\mathrm{?}`$ will provide a mm/sub-mm-wave map of the whole sky with a resolution of about 5 arcmin down to a detection limit of order 100 mJy at wavelengths of 350, 550 and 850 $`\mu `$m. These wavelengths are expected to be close to the peak of the redshifted thermal dust emission spectrum of distant galaxies. The effect of redshifting the dust spectrum peak into these observing bands is to provide access to the distant Universe, with little contamination from bright low-redshift galaxies. This remarkable, helpful $`K`$-correction is confirmed observationally to be effective in the sub-mm waveband, and is an almost unique feature of the waveband, although a similar effect may be at work in the hard X-ray surveys now being carried out. A priori, the counts of sub-mm galaxies at the faintest depths probed by Planck are expected to be very steep, an expectation which seems to be confirmed by the results of recent ground-based sub-mm-wave surveys.$`^\mathrm{?}`$
Very steep sub-mm counts generate a strong positive magnification bias,$`^{\mathrm{?},\mathrm{?}}`$ and thus boost the fraction of lensed galaxies in a flux-limited sample. This enhancement to the surface density of lensed galaxies will hopefully be exploited in the Planck survey to allow a very large catalogue of mainly unknown strongly lensed galaxies to be compiled.
## 2 Sensitivity and confusion noise in a Planck survey
The expected sensitivities of the Planck-HFI instrument in its three highest frequency channels, in which the largest number of distant dusty galaxies and lenses will be detected,$`^\mathrm{?}`$ are listed in Table 1, along with the predicted levels of source confusion noise.
It is important to understand the significance of source confusion noise introduced into the Planck images due to the varying number and flux densities of unresolved dusty galaxies in the observing beam. At 5 arcmin across, the Planck beam is very much larger than the mean separation of $`L^{}`$ galaxies, and also about 4 times greater in area than the 5-arcmin<sup>2</sup> field of view of the SCUBA camera at the JCMT, from which most observational information about the sub-mm source population has thus far been derived.
In Fig. 1 the results of many independent simulations of sampling a random unclustered distribution of sources on the sky with the observing beam in the relevant Planck-HFI channels are compared with the estimates of instrumental noise. Confusion noise due to Galactic cirrus predicted by IRAS-based studies,$`^\mathrm{?}`$ and of simple estimates based on the surface density of galaxies,$`^{\mathrm{?},\mathrm{?}}`$ are compared in Table 1. All these calculations of the confusion signal due to extragalactic sources are made assuming a model for the underlying surface density of sources on the sky that is at the top end of the range of possible values, the predictions are therefore likely to be conservatively high, and thus pessimistic for making reliable detections. Extragalactic confusion noise is expected to dominate both the instrumental and Galactic noise in these bands.
In Fig. 1, the results of the simulations are enveloped by solid curves, which represent log-normal distributions providing good representations of the results. The value of $`\sigma `$ in the log-normal distribution ($`\mathrm{exp}[\mathrm{ln}x\overline{x}]^2/2\sigma ^2`$) at all three frequencies is close to 0.2.
Unfortunately, the actual level of confusion noise that will be contributed to the Planck image depends on the uncertain surface density of galaxies that are brighter than those selected in existing sub-mm surveys (several tens of mJy at 850 $`\mu `$m/350 GHz). The results of future wider-field ground-based sub-mm surveys, and experience gained from the results of long-duration balloon-borne mm- and sub-mm-wave CMB experiments, including the existing BOOMERanG data and forthcoming data from TOPHAT will be useful in this respect. Recent simulations of extraction algorithms,$`^\mathrm{?}`$ in which the surface density of galaxies is assumed to be a factor of 2-3 times lower than that assumed here,$`^\mathrm{?}`$ suggest that 350-GHz point sources with flux densities greater than of order 75 mJy could be extracted from the Planck all-sky map.
## 3 Expected source density and numbers of detections
The two very different counts of unlensed galaxies shown in Fig. 2 should provide an envelope to the maximum range of possible values of the sub-mm counts at flux densities at which galaxies will be detected using Planck. The two underlying models$`^{\mathrm{?},\mathrm{?}}`$ predict very different bright sub-mm counts, but are both consistent with the observed far-infrared and sub-mm-wave counts and background radiation intensity, and with what little is currently known about the redshift distribution of sub-mm-selected galaxies. Our knowledge of the population will continue to improve until the launch of Planck. Some of the facilities that will contribute to this knowledge are listed elsewhere.$`^\mathrm{?}`$
Current observational limits to the population of sub-mm galaxies are reasonably well determined at flux densities less than about 10 mJy at 350 GHz from the results of several independent surveys: the surface density of galaxies brighter than 4 mJy at 850 $`\mu `$m/350 GHz is about 2000 deg<sup>-2</sup>. These surveys have mainly been carried out using the SCUBA camera at the JCMT, but results are also now coming in from the MAMBO 1.25-mm bolometer array camera at the IRAM 30-m telescope.$`^\mathrm{?}`$ These are flux densities considerably fainter than Planck will probe, and so far, these instruments have mapped only very small regions of sky (several hundred square arcminutes). The counts of objects brighter than about 10 mJy has been only weakly constrained.
Information is also available about the counts at the very brightest flux densities, based on the results of a targeted SCUBA survey of low-redshift galaxies selected from the IRAS catalogue.$`^\mathrm{?}`$ This survey was used to define a luminosity function of local IRAS galaxies in the sub-mm, which can be used to impose a lower limit to the bright counts at flux densities brighter than the detection limit in the all-sky Planck survey: this lower limit is about 10 sources on the sky brighter than 1 Jy at 850 $`\mu `$m/350 GHz.
The probability of lensing by foreground galaxies out to the high redshifts, at which a significant fraction of the sources detected in the Planck survey are expected to lie, is reasonably well-known, certainly to within a factor of a few. Specific predictions of this probability, tailored to observations in the sub-mm waveband, are discussed elsewhere.$`^{\mathrm{?},\mathrm{?}}`$ The predicted density of lenses on the sky should also be quite well determined at the flux densities that will be probed by Planck. This is because the intrinsic flux densities of these objects, after correcting for the effect of lensing, are likely to be of order 10 mJy. At this flux density the surface density of galaxies is reasonably well constrained by SCUBA observations.$`^\mathrm{?}`$ Predicted numbers of lensed objects$`^\mathrm{?}`$ are shown by the dot-dashed lines in Fig. 2.
A significant potential uncertainty remains in these results, however, because the maximum magnification that can be produced by a lens is limitted by the size of the source, being smaller for larger sources. In these calculations, we have assumed that background high-redshift sources are less than about 10 kpc in size, so that a maximum magnification of several tens can be produced. This assumption is supported by high-resolution interferometric observations of the size of the continuum emitting region in low-redshift ultraluminous galaxies, which is found to be much smaller – several 100 pc across.$`^\mathrm{?}`$ However, there are indications that at least some very luminous dusty high-redshift galaxies display CO and dust emission on scales greater than 10 kpc.$`^{\mathrm{?},\mathrm{?}}`$ Unfortunately, sub-arcsecond resolution images are required in order to measure the sizes of high-redshift dusty galaxies smaller than 10 kpc, and the necessary interferometric observations are difficult and time-consuming. Hence, as we do not yet know the size distribution of these distant ultraluminous galaxies, we cannot be certain that their sizes are typically small enough for large lensing magnifications to be possible. Despite this caveat, individual examples of high-redshift sub-mm objects lensed by foreground galaxies with magnifications of several tens are known,$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ and so it is likely that Planck will be able to detect a significant number.
Based on an extrapolation of the log-normal distribution of pixel values predicted by the conservative simulations of confusion noise shown in Fig. 1, it is possible to predict the number of pixels that are likely to exceed a certain flux density due to the effects of confusion over the all-sky Planck survey. This expected count of spurious detections, which is shown by the dotted lines in Fig. 2, can then be compared with the expected count of lensed and unlensed galaxies. Reliable detections should be possible if these counts lie safely above the dotted lines at any flux density. Hence, it is likely that lensed galaxies brighter than about 250, 500 and 1500 mJy at 353, 545 and 857 GHz respectively could be detected using Planck, even in this pessimistic model of the confusion noise. At these limits, of order 100 lensed sources are expected over the whole sky. If the counts of bright unlensed sources are less than assumed here,$`^\mathrm{?}`$ then the confusion noise will be reduced, and a greater number will be detected. The counts shown in Fig. 2 allow these different scenarios to be investigated directly. It is important to note that much better estimates of confusion noise will be available well in advance of the launch of Planck.
The emission from both lensed and unlensed distant galaxies will pass unattenuated through the Galactic plane in the Planck bands; however, the ability to distinguish them against a bright and structured Galactic foreground is likely to be limited. The level of Galactic confusion (see Table 1) depends on the Galactic background intensity $`B_0`$ as $`B_0^{1.5}`$.$`^\mathrm{?}`$ The results suggest that the sensitivity of the Planck survey can only be exploited fully to find extragalactic sources in regions of the sky where the 100-$`\mu `$m surface brightness from the Milky Way is less than about 5 MJy sr<sup>-1</sup>, a condition which is satisfied over most of the sky.$`^\mathrm{?}`$
## 4 Follow-up observations
From the Planck all-sky images alone, point sources will be detectable, but no information will be available about their morphology, including whether or not they display arc and multiple image geometries characteristic of lensing. For this diagnosis, sub-arcsecond mm/sub-mm images using the high-resolution extremely-sensitive interferometer array ALMA,$`^\mathrm{?}`$ or very deep radio images using the VLA will be required. Accurate positions for Planck sources will first be required: these could be obtained using either ALMA or the 3.5-m ESA cornerstone FIRST telescope.$`^\mathrm{?}`$ Note that it is very unlikely for the foreground lens galaxy to be bright in the sub-mm waveband,$`^\mathrm{?}`$ as this would require the foreground lens itself to be a very luminous infrared galaxy, and the space density of such sources is rather low. Hence, only the lensed images will be detected at bright flux levels in an ALMA continuum image, free from the blending and masking effects of emission from the foreground lens. Since the first discussion of a Planck lens survey,$`^\mathrm{?}`$ the specifications of ALMA have been significantly refined.$`^\mathrm{?}`$ At sub-mm flux densities of order 100 mJy, ALMA will be able to detect galaxies in integrations lasting much less than a second, and to make high-quality images in only several minutes. ALMA will make it easy to confirm lensing features and obtaining good quality images of the galaxies detected in a Planck survey.
The form of the mid-infrared SED of the galaxies detected using Planck can be determined using the SPIRE and PACS instruments aboard FIRST, allowing the redshift/dust temperature of the detected galaxies to be determined. Direct spectroscopic redshifts for the lensed galaxies should come quite easily from molecular emission line data taken using ALMA or from observations of features in the redshifted mid-infrared spectrum taken using FIRST-PACS.$`^\mathrm{?}`$ Redshifts for the lensing galaxies should be readily obtained from the frequencies of CO absorption lines imposed on the bright dust continuum emission of the magnified background galaxies due to the interstellar medium in the foreground lens, a technique which has already been developed and demonstrated using existing mm-wave interferometers.$`^\mathrm{?}`$
Other relevant information can be provided by the FIRST radio survey at the VLA, which covers over 5000 deg<sup>2</sup> to a depth of about 1 mJy at 1.4 GHz. The high-resolution radio images from FIRST should detect a significant fraction of the Planck detected galaxies with sub-mm flux densities of several hundred mJy, based on our current knowledge of the SEDs of confirmed high-redshift SCUBA galaxies. For example, the 25-mJy 850-$`\mu `$m sub-mm source SMM J02399$``$0136 has a 1.4-GHz flux of 520 $`\mu `$Jy.$`^\mathrm{?}`$ The brighter cousins of these objects in the Planck survey might typically appear at the faintest levels in the FIRST/VLA survey. The 25-mJy 60-$`\mu `$m all-sky image from the IRIS/ASTRO-F satellite$`^\mathrm{?}`$ will also be useful for detecting Planck sources.
## 5 Summary
The Planck Surveyor CMB imaging mission should detect many thousands of high redshift dusty galaxies, a significant fraction of which are expected to be gravitationally lensed. In order to make more detailed predictions it will be necessary to better quantify the abundance of sub-mm-selected galaxies with flux densities of order 100 mJy. By combining the Planck results with observations made using FIRST and ALMA a large catalogue of lenses should be identified, providing a unique sample of objects for further study. The properties of the individual galaxies in the catalogue will be useful for obtaining information about both the geometry of the Universe and the evolution of distant dusty galaxies.
## Acknowledgements
The author, Raymond and Beverly Sackler Foundation Research Fellow, gratefully acknowledges generous support from the Foundation as part of their Deep Sky Initiative Programme at the IoA. Attendance at the meeting was made possible by the European Union TMR programme, which supports European research in gravitational lensing via contract number ERBFMRXCT98-0172 (Lensnet). This work was carried out within the Cambridge Planck Analysis Centre (CPAC), supported by PPARC. I thank Vicki Barnard, Dave Frayer, Ben Metcalf, Daniel Mortlock and members of the FIRST-Planck extragalactic science team for useful discussions and Kate Quirk for comments on the manuscript.
## References
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# SUPERBOX – An Efficient Code for Collisionless Galactic Dynamics
## 1 Introduction
In direct-summation N-body methods, the computational time $`t_{CPU}`$ scales with the square of the particle number $`N_\mathrm{p}`$, $`t_{CPU}N_\mathrm{p}^2`$. This relation inevitably leads to bottlenecks when computing large-N galaxy models. While it is still impossible to model galaxies star by star over sensible time-laps, alternative algorithms have been developed over a number of years which ease the computer requirements at the expense of the exact $`N_\mathrm{p}^2`$ summation, yet provide an adequate treatment of Newtonian galactic dynamics. Examples of such techniques are the tree method (Barnes & Hut 1986, Davé, Dubinski & Hernquist 1997), and the multi-grid or nested grid particle mesh-codes like Pandora (Villumsen 1989), Hydra (Pearce & Couchman 1997), Superbox, or combinations of the methods like the adaptive refinement tree (ART) (Kravtsov, Klypin & Khokhlov. 1997), which can cope with very inhomogeneous matter distributions and have high spatial resolution at the density maxima.
While the hierarchical tree-method is independent of the geometry of the problem, its main disadvantage is the limitation in particle number and the softening needed to avoid undesired two-body relaxation. In contrast, the main advantage of a particle-mesh technique is the ability to use very high particle numbers (up to several millions on desk-top computers), because the CPU-time scales only linearly with the number of particles, and with $`N_{\mathrm{gc}}\mathrm{log}N_{\mathrm{gc}}`$, where $`N_{\mathrm{gc}}`$ is the number of grid-cells. High particle numbers keep the statistical noise low. The disadvantage of a grid-based method, however, is the dependency on the geometry of the grid. But with increasing spatial resolution, the geometry of the cells becomes less important. The only limitation in resolving phenomena is the size of one grid-cell. A general discussion and comparison of various numerical schemes will be found in Sellwood (1987), and – in the context of direct $`N`$-body simulations – in Spurzem (1999).
The basic idea behind Superbox, as originally conceived by R. Bien (Bien et al. 1991), is to increase the resolution only at places where it is necessary, while simultaneously keeping the computational overhead as small as possible by using a fixed (i.e. not adaptive) nested grid-architecture. Accuracy is improved by the use of nested high-resolution sub-grids and a linear force interpolation to the exact position of the particle inside a cell. Two higher-resolution sub-grids are introduced: the medium-resolution grid contains an entire galaxy, and the high-resolution grid treats its core. Both grids stay focused on the galaxy, moving through the ’local universe’ that is contained in the coarse outermost grid.
Superbox has already been used successfully in a variety of investigations, namely in the study of the high velocity encounter of the two similar early-type galaxies NGC 4782/4783 by Madejski & Bien (1993), in the survey of encounters between elliptical galaxies by Wassmer et al. (1993), and in the research on tidal disruption of satellite dwarf galaxies by Kroupa (1997) and Klessen & Kroupa (1998). A version of Superbox exists that includes sticky particles as a simple model of the dynamics of cold molecular gas clouds (Fellhauer 1996).
In this article we provide an implementation and benchmarks for the Superbox algorithm. In Section 2, we lay out the mathematics behind the algorithm and give details of how the potential is mapped on the different grids. We then move on to test the code: Section 3 deals with the conservation of energy and angular momentum, while Section 4 discusses relaxation. Section 5 contains profiling information about storage and CPU-time requirements. Comparison with a direct-summation $`N_\mathrm{p}`$-body code is made in Section 6. In Section 7 we show – as an example – the good agreement of satellite decay with Chandrasekhar’s formula for dynamical friction, and close with conclusions in Section 8.
## 2 Method
Detailed discussions of the particle-mesh (PM) technique can be found in Eastwood & Brownrigg (1978), Hockney & Eastwood (1981) and Sellwood (1987).
Briefly, the array of mass densities, $`\varrho _{ijk}`$, is derived by counting the number of particles in each Cartesian grid cell $`i,j,k`$, i.e. by using the simple particle-in-cell (NGP = “Nearest Grid Point”) algorithm. All particles belonging to one galaxy $`N_{\mathrm{p},\mathrm{gal}}`$ have the same mass, $`m=M_{\mathrm{gal}}/N_{\mathrm{p},\mathrm{gal}}`$. An alternative would be to use a cloud-in-cell technique, which is known to improve momentum and energy conservation, in which the mass of the particle is distributed proportionally among the neighbouring cells. This is, however, not necessary, since the number of particles is large.
Poisson’s equation is solved for this density array to get the grid-based potential, $`\mathrm{\Phi }_{ijk}`$, which becomes,
$`\mathrm{\Phi }_{ijk}`$ $`=`$ $`G{\displaystyle \underset{a,b,c=0}{\overset{N1}{}}}\varrho _{abc}H_{ai,bj,ck},i,j,k=0,\mathrm{},N1,`$ (1)
where $`N`$ denotes the number of grid-cells per dimension ($`N^3=N_{\mathrm{gc}}`$), and $`H_{ijk}`$ is some Green’s function. To avoid this $`N_{\mathrm{gc}}^2`$ procedure, the discrete Fast Fourier Transform (FFT) is used, for which $`N=2^K`$, $`K>0`$ being an integer. The stationary Green’s function is Fourier transformed once at the beginning of the calculation, and only the density array is transformed at each time-step:
$`\widehat{\varrho }_{abc}`$ $`=`$ $`{\displaystyle \underset{i,j,k=0}{\overset{N1}{}}}\varrho _{ijk}\mathrm{exp}\left(\sqrt{1}{\displaystyle \frac{2\pi }{N}}\left(ai+bj+ck\right)\right),`$ (2)
$`\widehat{H}_{abc}`$ $`=`$ $`{\displaystyle \underset{i,j,k=0}{\overset{N1}{}}}H_{ijk}\mathrm{exp}\left(\sqrt{1}{\displaystyle \frac{2\pi }{N}}\left(ai+bj+ck\right)\right).`$
The two resulting arrays are multiplied cell by cell and transformed back to get the grid-based potential,
$`\mathrm{\Phi }_{ijk}`$ $`=`$ $`{\displaystyle \frac{G}{N^3}}{\displaystyle \underset{a,b,c=0}{\overset{N1}{}}}\widehat{\varrho }_{abc}\widehat{H}_{abc}\mathrm{exp}\left(\sqrt{1}{\displaystyle \frac{2\pi }{N}}\left(ai+bj+ck\right)\right).`$ (3)
The potential is differentiated numerically to estimate the accelerations for each particle, each of which is pushed forward on its orbit using the leapfrog scheme (see Fig. 1 and sections below).
### 2.1 Green’s function and the FFT
The simplest Green’s function is adopted, in which the side of a grid cell has unit length:
$`H_{ijk}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{i^2+j^2+k^2}}},i,j,k=0,1,..,N`$ (4)
$`H_{000}`$ $`=`$ $`4/3.`$
The value of $`H_{000}`$ is somewhat arbitrary, and $`H_{000}=H_{100}`$ is often used (see e.g. Sellwood 1987). Numerical tests performed with Superbox show that for high particle numbers, $`H_{000}=4/3`$ yields a slightly slower drift in total energy with time (see also Section 3) than $`H_{000}=1`$ or 1.5, while for low particle numbers (compared to the number of grid-cells) $`H_{000}=1`$ is the better choice. The reason for this is that for high particle numbers the collective force dominates over the particle–particle (pairwise) force. On the other hand, for low particle numbers (only one or a few particles are together in the same cell) pairwise forces are more dominant, and it is better to exclude ’self-gravity’ (i.e. a particle feels its own force inside its cell), with $`H_{000}=1`$. This is consistent with analytical arguments by D. Pfenniger (private communication). He shows that in one dimension, $`H_0=\mathrm{ln}4`$ minimises the power at the grid wave-number $`n/2`$ for the kernel (compare with Eq. 3)
$`H_i`$ $`=`$ $`1/|i|,0<\left|i\right|n/2,`$ (5)
$`H_i`$ $`=`$ $`\mathrm{ln}4,i=0.`$
He finds that the power is exactly cancelled at the wavenumber $`n/2`$, and for finite $`n`$ can also be cancelled for a value of $`H_0`$ close to $`\mathrm{ln}4`$. For $`H_0>\mathrm{ln}4`$ the power at high frequencies increases, while for $`H_0<\mathrm{ln}4`$ some wavenumbers $`<n/2`$ are filtered out. In this sense, $`H_0=\mathrm{ln}44/3`$ optimises the kernel. However, the 3D-problem is slightly different since forces instead of the potential are used. The artificial wavenumbers induced by the grid, that are to be minimised, are then much more numerous. This problem would need further addressing.
The FFT-algorithm gives the exact solution of the grid-based potential for a periodical system. For the exact solution of an isolated system, which is what we are interested in, the size of the density-array has to be doubled ($`2N`$), filling all inactive grid cells with zero density, and extending the Green’s function in the empty regions in the following way:
$`H_{2ni,j,k}`$ $`=`$ $`H_{2ni,2nj,k}=H_{2ni,j,2nk}=H_{2ni,2nj,2nk}`$
$`=`$ $`H_{i,2nj,k}=H_{i,2nj,2nk}=H_{i,j,2nk}=H_{i,j,k}.`$
This provides the isolated solution of the potential in the simulated area between $`i,j,k=0`$ and $`N1`$. In the inactive part the results are unphysical. To keep storage as low as possible, only a $`2N\times 2N\times N`$-array is used for transforming the densities, and a $`(N+1)\times (N+1)\times (N+1)`$-array is used for the Green’s function. For a detailed discussion see Eastwood & Brownrigg (1978) and also Hockney & Eastwood (1981).
The FFT-routine incorporated in Superbox is a simple one-dimensional FFT and is taken from Werner & Schabach (1979) and Press et al. (1986). It is fast and makes the code portable and not machine specific. The low-storage algorithm to extend the FFT to 3 dimensions, to obtain the 3-D-potential, is taken from Hohl (1970). The performance of Superbox can be increased by incorporating machine-optimised FFT routines.
### 2.2 Acceleration and orbit integration
In this section we explain the force calculation performed in Superbox, and compare it with standard methods. We follow closely the notation of Hockney & Eastwood (1981) and Couchman (1999) to document the differences.
With the mass density as an example we show, in a formal way, how its values on a 3D mesh are obtained. Using the $`\delta `$-functional for the particle shape function, a particle distribution $`n(𝐱)`$ is given by
$`n(𝐱)`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}\delta ^3(𝐱𝐱_\alpha ),`$ (7)
where $`\alpha `$ is a summation index running over all particles, and $`𝐱_\alpha =(x_\alpha ,y_\alpha ,z_\alpha )`$ is the position vector of particle $`\alpha `$. As smoothing kernel we select here
$`W(𝐱,𝚫𝐱)`$ $`=`$ $`\mathrm{\Pi }\left({\displaystyle \frac{x}{\mathrm{\Delta }x}}\right)\mathrm{\Pi }\left({\displaystyle \frac{y}{\mathrm{\Delta }y}}\right)\mathrm{\Pi }\left({\displaystyle \frac{z}{\mathrm{\Delta }z}}\right).`$
$`𝚫𝐱`$ denotes a vector of smoothing lengths $`(\mathrm{\Delta }x,\mathrm{\Delta }y,\mathrm{\Delta }z)`$ in three coordinate directions. $`W`$ is a three dimensional generalisation of the standard top-hat function
$`\mathrm{\Pi }(\xi )`$ $`=`$ $`\{\begin{array}{cc}0,\hfill & |\xi |>\frac{1}{2}\text{,}\hfill \\ \frac{1}{2},\hfill & |\xi |=\frac{1}{2}\text{,}\hfill \\ 1,\hfill & |\xi |<\frac{1}{2}\text{ .}\hfill \end{array}`$
Hence the smoothed mass-density field is
$`\varrho (𝐱)`$ $`=`$ $`mWn=m{\displaystyle W(𝐱𝐱^{},𝚫𝐱)n(𝐱^{})\mathrm{d}^3𝐱^{}},`$ (8)
where $`m`$ denotes the particle mass, and the $``$-operator a convolution as defined by the equation above. Mathematically we get from $`\varrho (𝐱)`$ a so-called mesh sampled functional $`\varrho ^{}(𝐱)`$ on a three dimensional mesh by
$`\varrho ^{}(𝐱)`$ $`=`$ $`(𝐱)\varrho (𝐱),`$ (9)
using the three-dimensional comb or sampling function defined as
$`(𝐱)`$ $`=`$ $`{\displaystyle \underset{i,j,k=0}{\overset{N}{}}}\delta (xx_{ijk})\delta (yy_{ijk})\delta (zz_{ijk}),`$
where $`i,j,k`$ are indices of grid cells, whose centres are located at the vector points $`𝐱_{ijk}=\{x_{ijk},y_{ijk},z_{ijk}\}`$. From the functional $`\varrho ^{}(𝐱)`$ a finite set of discrete mesh-sampled density values $`\varrho _\mathrm{m}^{}=\{\varrho _{ijk},i,j,k=0,1,\mathrm{},N\}`$ is obtained, which are defined at the centres of the grid cells. For each cell centre or mesh point they are computed by integrating $`\varrho ^{}`$ over an arbitrarily shaped and sized volume containing just this and only this mesh point. The standard discrete FFT procedure (see above) is applied to the set of numbers $`\varrho _\mathrm{m}^{}`$ in order to determine the grid-based potential $`\mathrm{\Phi }_\mathrm{m}^{}=\{\mathrm{\Phi }_{ijk},i,j,k=0,1,\mathrm{},N\}`$. $`\mathrm{\Phi }_\mathrm{m}^{}`$ is a set of numbers resulting from the FFT and represents the mesh sampled functional $`\mathrm{\Phi }^{}(𝐱)=(𝐱)\mathrm{\Phi }(𝐱)`$, where $`\mathrm{\Phi }(𝐱)`$ is a true smooth gravitational potential. The set of values $`\mathrm{\Phi }_\mathrm{m}^{}`$ is obtained from $`\mathrm{\Phi }^{}`$ by integrating over volumes containing single mesh points, just in the same way as explained above for $`\varrho `$. We define the one-dimensional two-point difference operator for $`x`$ direction in the standard way (Hockney & Eastwood 1981, p. 163),
$`D_x(x,y,z,\mathrm{\Delta }x)`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Delta }x}}\left(\delta (x+\mathrm{\Delta }x)\delta (x\mathrm{\Delta }x)\right)\delta (y)\delta (z).`$
To keep the notation clearer we remain for the moment in one dimension. A first order difference approximation to the x-component of the acceleration vector is given by
$`a_x{}_{}{}^{(1)}(x,y,z,\mathrm{\Delta }x)`$ $`=`$ $`D_x\mathrm{\Phi }(x)={\displaystyle \frac{\mathrm{\Phi }(x+\mathrm{\Delta }x,y,z)\mathrm{\Phi }(x\mathrm{\Delta }x,y,z)}{2\mathrm{\Delta }x}}.`$ (10)
Sampling this on the mesh yields
$`a_x^{}{}_{}{}^{(1)}(x,y,z,\mathrm{\Delta }x)`$ $`=`$ $`(𝐱)a_x{}_{}{}^{(1)}(x,y,z)`$ (11)
$`=`$ $`(𝐱){\displaystyle \frac{\mathrm{\Phi }(x+\mathrm{\Delta }x,y,z)\mathrm{\Phi }(x\mathrm{\Delta }x,y,z)}{2\mathrm{\Delta }x}}.`$
If this is integrated over volumes containing a single mesh point, again exactly as described above for the cases of $`\varrho `$ and $`\mathrm{\Phi }`$, a set of accelerations $`a_{ijk}^{(1)}`$ defined at the mesh point centres results:
$`a_{x,ijk}^{(1)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_{i+1,jk}\mathrm{\Phi }_{i1,jk}}{2\mathrm{\Delta }x}};`$ (12)
here we have assumed that $`2\mathrm{\Delta }x=x_{i+1,jk}x_{i1,jk}`$ is selected so as to match with the mesh point distances. In our notation the set of values $`\{a_{x,ijk}{}_{}{}^{(1)},i,j,k=0,1,\mathrm{},N\}`$ is denoted by $`a_{\mathrm{m},\mathrm{x}}^{}^{(1)}`$. The corresponding vector components for the acceleration in $`y`$ and $`z`$ direction are analogous. In the standard particle-mesh schemes, the accelerations between the mesh centres are obtained by using the same smoothing kernel as for the density, i.e.
$`a_x{}_{}{}^{(1)}(x,y,z,\mathrm{\Delta }x)=Wa_x^{}{}_{}{}^{(1)}(x,y,z)=a_{x,ijk}{}_{}{}^{(1)}`$ $`\mathrm{for}`$ $`xx_{ijk}<{\displaystyle \frac{\mathrm{\Delta }x}{2}}`$
$`\mathrm{and}`$ $`yy_{ijk}<{\displaystyle \frac{\mathrm{\Delta }y}{2}}`$
$`\mathrm{and}`$ $`zz_{ijk}<{\displaystyle \frac{\mathrm{\Delta }z}{2}}.`$
The use of the $`\delta `$-functional for the particle shape and the top-hat function for smoothing gives the rather crude standard NGP scheme, where the acceleration on a particle is constant in each cell and has a discontinuity at the cell boundaries. A possible improvement regarding energy and momentum conservation is to use a so-called cloud-in-cell or CIC scheme, where the shape function of particles is a top-hat function, and the assignment function is a triangle function (see Hockney & Eastwood 1981). However, as one can also see in the cited book, such schemes suffer from force anisotropies, despite having very good conservation properties (i.e. the magnitude and direction of the force error depends on whether one goes parallel to the mesh coordinates or not).
Therefore a non-standard scheme is applied in Superbox. It is NGP in nature, provides a very good energy and angular momentum conservation and a significant improvement for force anisotropies. Note that the rather involved mathematical formalism above will later yield significant hindsight what is special with Superbox. The notation follows closely Hockney & Eastwood’s and Couchman’s.
We begin with a Taylor expansion of the acceleration (for example the $`x`$-component) centred on $`(x,y,z)`$, the position of the centre of cell $`i,j,k`$:
$`a_x{}_{}{}^{(2)}(x+dx,y+dy,z+dz)`$ $`=`$ $`a_x{}_{}{}^{(2)}(x,y,z)+`$
$`{\displaystyle \frac{a_x}{x}}(x,y,z)dx`$ $`+`$ $`{\displaystyle \frac{a_x}{y}}(x,y,z)dy+{\displaystyle \frac{a_x}{z}}(x,y,z)dz`$
$`+`$ $`𝒪(\mathrm{𝐝𝐱}^2),`$
where $`𝒪(\mathrm{𝐝𝐱}^2)`$ denotes any higher order terms in $`dx`$, $`dy`$, $`dz`$ or combinations thereof. The displacements $`dx`$, $`dy`$ and $`dz`$ should not be linked or confused with the mesh spacings $`\mathrm{\Delta }x`$, $`\mathrm{\Delta }y`$, $`\mathrm{\Delta }z`$. They are free and should provide an interpolated expression for $`a_x`$ at any point inside a given mesh cell, e.g. a particle’s position. We then use the difference approximations
$`a_x{}_{}{}^{(2)}(x,y,z,\mathrm{\Delta }x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }}{x}}(x,y,z,\mathrm{\Delta }x)=D_x\mathrm{\Phi },`$
$`{\displaystyle \frac{a_x}{x}}(x,y,z,\mathrm{\Delta }x,\mathrm{\Delta }y,\mathrm{\Delta }z)`$ $`=`$ $`{\displaystyle \frac{^2\mathrm{\Phi }}{x^2}}(x,y,z,\mathrm{\Delta }x,\mathrm{\Delta }y,\mathrm{\Delta }z)D_{xx}\mathrm{\Phi },`$
$`{\displaystyle \frac{a_x}{y}}(x,y,z,\mathrm{\Delta }x,\mathrm{\Delta }y,\mathrm{\Delta }z)`$ $`=`$ $`{\displaystyle \frac{^2\mathrm{\Phi }}{xy}}(x,y,z,\mathrm{\Delta }x,\mathrm{\Delta }y,\mathrm{\Delta }z)D_{xy}\mathrm{\Phi },`$
$`{\displaystyle \frac{a_x}{z}}(x,y,z,\mathrm{\Delta }x,\mathrm{\Delta }y,\mathrm{\Delta }z)`$ $`=`$ $`{\displaystyle \frac{^2\mathrm{\Phi }}{xz}}(x,y,z,\mathrm{\Delta }x,\mathrm{\Delta }y,\mathrm{\Delta }z)D_{xz}\mathrm{\Phi }.`$ (14)
We may now generate a mesh sampled $`a_x^{}{}_{}{}^{(2)}=a_x^{(2)}`$ in the standard way and integrate this over any volume containing one grid cell centre to obtain a set of acceleration values $`a_{\mathrm{m},\mathrm{x}}^{}^{(2)}`$ given by the following expressions:
$`a_{ijk,x}{}_{}{}^{(2)}(dx,dy,dz)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_{i+1,jk}\mathrm{\Phi }_{i1,jk}}{2\mathrm{\Delta }x}}`$
$`+`$ $`{\displaystyle \frac{\mathrm{\Phi }_{i+1,jk}+\mathrm{\Phi }_{i1,jk}2\mathrm{\Phi }_{ijk}}{(\mathrm{\Delta }x)^2}}dx`$
$`+`$ $`{\displaystyle \frac{\mathrm{\Phi }_{i+1,j+1,k}\mathrm{\Phi }_{i1,j+1,k}+\mathrm{\Phi }_{i1,j1,k}\mathrm{\Phi }_{i+1,j1,k}}{4\mathrm{\Delta }x\mathrm{\Delta }y}}dy`$
$`+`$ $`{\displaystyle \frac{\mathrm{\Phi }_{i+1,j,k+1}\mathrm{\Phi }_{i1,j,k+1}+\mathrm{\Phi }_{i1,j,k1}\mathrm{\Phi }_{i+1,j,k1}}{4\mathrm{\Delta }x\mathrm{\Delta }z}}dz`$
Note that we generally take $`\mathrm{\Delta }x=\mathrm{\Delta }y=\mathrm{\Delta }z=1`$ i.e. the cell-length is assumed to be equal along the three axes and unity. The accelerations in $`y`$\- and $`z`$-direction are calculated analogously.
At first glance the interpolation scheme Eq. 2.2 for the acceleration seems to be closely related to the procedure used in a standard cloud-in-cell (CIC) scheme. In one dimension, CIC and Superbox methods can be more easily compared to each other. For $`0dx<\mathrm{\Delta }x/2`$ one has
$`\mathrm{CIC}:a(x+dx)`$ $`=`$ $`a_i{\displaystyle \frac{\mathrm{\Delta }xdx}{\mathrm{\Delta }x}}+a_{i+1}{\displaystyle \frac{dx}{\mathrm{\Delta }x}}`$
$`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_{i+1}\mathrm{\Phi }_{i1}}{2\mathrm{\Delta }x}}+{\displaystyle \frac{\mathrm{\Phi }_{i+2}\mathrm{\Phi }_{i+1}\mathrm{\Phi }_i+\mathrm{\Phi }_{i1}}{2(\mathrm{\Delta }x)^2}}dx,`$
$`\mathrm{𝖲𝗎𝗉𝖾𝗋𝖻𝗈𝗑}:a(x+dx)`$ $`=`$ $`a_i+{\displaystyle \frac{\mathrm{d}a}{\mathrm{d}x}}|_idx`$
$`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_{i+1}\mathrm{\Phi }_{i1}}{2\mathrm{\Delta }x}}+{\displaystyle \frac{\mathrm{\Phi }_{i+1}+\mathrm{\Phi }_{i1}2\mathrm{\Phi }_i}{(\mathrm{\Delta }x)^2}}dx,`$
where $`i`$ denotes the cell-index. Clearly this is different, because CIC uses accelerations in two neighbouring cells, while Superbox takes acceleration and its derivative only in one cell. Therefore we propose to consider Superbox as an NGP type scheme, although a second order interpolated acceleration is used, because nowhere the acceleration of the neighbouring cell is entering into the equations. Though of NGP nature, $`a_x`$ in Superbox behaves steadily when crossing cell boundaries along the $`x`$-axis (and $`a_y`$, $`a_z`$ correspondingly on their respective axes). In contrast to CIC our scheme is not steady in any other direction. However, the errors induced by that are much smaller than in a standard NGP scheme due to the above higher order interpolation, and are perfectly tolerable as will be shown throughout this paper.
The salient feature of Superbox regarding the force anisotropies can be understood by considering the force at a distance $`r=\left|\mathrm{𝐝𝐫}\right|`$ from a particle located at the centre of a cell. To lowest order the Superbox force-error $`\epsilon `$ is then
$`\stackrel{}{\epsilon }`$ $`=`$ $`{\displaystyle \frac{\mathrm{𝐝𝐫}}{r^3}},`$ (16)
where $`\mathrm{𝐝𝐫}=(dx,dy,dz)`$ is the displacement vector from the cell centre, and we have neglected all terms of second order in $`\mathrm{\Delta }x`$, $`\mathrm{\Delta }y`$, or $`\mathrm{\Delta }z`$ in a Taylor series expansion of the potential differences. In other words to lowest order the Superbox force error points to the exact centre of the cell. Thus, compared with standard NGP schemes, we get a very precise force-calculation in Superbox, as shown in Figs. 2 and 3.
The orbits of the particles are integrated forward in time using the leapfrog scheme. For example, for the $`x`$-components of the velocity, $`v_x`$, and position, $`x`$, vectors of particle $`l`$,
$`v_{x,l}^{n+1/2}`$ $`=`$ $`v_{x,l}^{n1/2}+a_{x,l}^n\mathrm{\Delta }t`$ (17)
$`x_l^{n+1}`$ $`=`$ $`x_l^n+v_{x,l}^{n+1/2}\mathrm{\Delta }t,`$
where $`n`$ denotes the $`n`$th time-step and $`\mathrm{\Delta }t`$ is the length of the integration step.
### 2.3 The grids
For each galaxy, 5 grids with 3 different resolutions are used. This is made possible by invoking the additivity of the potential (Fig. 4).
The five grids are as follows:
* Grid 1 is the high-resolution grid which resolves the centre of the galaxy. It has a length of $`2\times R_{\mathrm{core}}`$ in one dimension. In evaluating the densities, all particles of the galaxy within $`rR_{\mathrm{core}}`$ are stored in this grid.
* Grid 2 has an intermediate resolution to resolve the galaxy as a whole. The length is $`2\times R_{\mathrm{out}}`$, but only particles with $`rR_{\mathrm{core}}`$ are stored here, i.e. the same particles as are also stored in grid 1.
* Grid 3 has the same size and resolution as grid 2, but it only contains particles with $`R_{\mathrm{core}}<rR_{\mathrm{out}}`$.
* Grid 4 has the size of the whole simulation area (i.e. ’local universe’ with $`2\times R_{\mathrm{system}}`$), and has the lowest resolution. It is fixed. Only particles of the galaxy with $`rR_{\mathrm{out}}`$ are stored in grid 4.
* Grid 5 has the same size and resolution as grid 4. This grid treats the escaping particles of a galaxy, and contains all particles with $`r>R_{\mathrm{out}}`$.
Grids 1 to 3 are focused on a common centre of the galaxy and move with it through the ’local universe’, as detailed below. All grids have the same number of cells per dimension, $`N`$, for all galaxies. The boundary condition, requiring two empty cells with $`\varrho =0`$ at each boundary, is open and non-periodic, thus providing an isolated system. This however means that only $`N4`$ active cells per dimension are used.
To keep the storage requirement low, all galaxies are treated consecutively in the same grid-arrays, whereby the particles belonging to different galaxies can have different masses. Each of the 5 grids has its associated potential $`\mathrm{\Phi }_\mathrm{i},i=1,2,\mathrm{},5`$ computed by the PM technique from the particles of one galaxy located as described above. The accelerations are obtained additively from the 5 potentials of each galaxy in turn in the following way:
$`\mathrm{\Phi }(r)`$ $`=`$ $`\left[\theta (R_{\mathrm{core}}r)\mathrm{\Phi }_1+\theta (rR_{\mathrm{core}})\mathrm{\Phi }_2+\mathrm{\Phi }_3\right]\theta (R_{\mathrm{out}}r)`$
$`+`$ $`\theta (rR_{\mathrm{out}})\mathrm{\Phi }_4+\mathrm{\Phi }_5,`$
$`\mathrm{\Phi }(R_{\mathrm{core}})`$ $`=`$ $`\mathrm{\Phi }_1+\mathrm{\Phi }_3+\mathrm{\Phi }_5`$
$`\mathrm{\Phi }(R_{\mathrm{out}})`$ $`=`$ $`\mathrm{\Phi }_2+\mathrm{\Phi }_3+\mathrm{\Phi }_5`$
where $`\theta (\xi )=1`$ for $`\xi >0`$ and $`\theta (\xi )=0`$ otherwise. This means:
* For a particle in the range $`rR_{\mathrm{core}}`$, the potentials of grids 1, 3 and 5 are used to calculate the acceleration.
* For a particle with $`R_{\mathrm{core}}<rR_{\mathrm{out}}`$, the potentials of grids 2, 3 and 5 are combined.
* And finally, if $`r>R_{\mathrm{out}}`$ then the acceleration is calculated from the potentials of grids 4 and 5.
* A particle with $`r>R_{\mathrm{system}}`$ is removed from the computation.
Due to the additivity of the potential (and hence its derivatives, the accelerations) the velocity changes originating from the potentials of each of the galaxies can be separately updated and accumulated in the first of the leap-frog formulas Eq. 17. The final result does not depend on the order by which the galaxies are taken into account and it could be computed even in parallel, if a final accumulation takes place. After all velocity changes have been applied in all galaxies, the positions of the particles are finally updated (see Fig. 1).
As long as the galaxies are well separated, they feel only the low-resolution potentials of the outer grids. But as the galaxies approach each other their high-resolution grids overlap, leading to a high-resolution force calculation during the interaction.
#### 2.3.1 Grid tracking
Two alternative schemes to position and track the inner and middle grids can be used. The most useful scheme centres the grids on the density maximum of each galaxy at each step. The position of the density maximum is found by constructing a sphere of neighbours centred on the densest region, in which the centre of mass is computed. This is performed iteratively. The other option is to centre the grids during run-time on the position of the centre of mass of each galaxy using all its particles remaining in the computation.
#### 2.3.2 Edge-effects
It is shown in Fig. 4 that only spherical regions of the cubic grids contain particles (except for Grid 5). Particles with eccentric orbits can cross the border of two grids, thus being subject to forces resolved differently. No interpolation of the forces is done at the grid-boundaries. This keeps the code fast and slim, but the grid-sizes have to be chosen properly in advance to minimise the boundary discontinuities. As shown in Section 4, this leads to some additional but negligible relaxation-effects, because the derived total potential has insignificant discontinuities at the grid-boundaries (Wassmer 1992). The best way to avoid these edge-effects is to place the grid-boundaries at ’places’ where the slope of the potential is not steep.
Fig. 5 shows that, as long as the grid-sizes are chosen properly, there are no serious effects at the grid-boundaries. Only in the example shown in the bottom-right panel the innermost grid was chosen too small, leading to a low particle number per grid-cell. Performing such computations with $`H_{000}=4/3`$ leads to spurious results, as explained in Section 2.1. Nevertheless the profile is smooth at the grid-boundary.
#### 2.3.3 Model units
Superbox employs model units, with the gravitational constant in model units being $`G_{\mathrm{mod}}=2`$ owing to historical reasons. A flag specifies if the user wants to input and output data in physical units (i.e. kpc, M, km/s, Myr). Scaling to and from physical units is achieved via the formulae:
$`\left({\displaystyle \frac{L_{\mathrm{phys}}}{L_{\mathrm{mod}}}}\right)^3`$ $`=`$ $`{\displaystyle \frac{G_{\mathrm{phys}}}{G_{\mathrm{mod}}}}{\displaystyle \frac{M_{\mathrm{phys}}}{M_{\mathrm{mod}}}}\left({\displaystyle \frac{T_{\mathrm{phys}}}{T_{\mathrm{mod}}}}\right)^2,`$ (19)
$`{\displaystyle \frac{V_{\mathrm{phys}}}{V_{\mathrm{mod}}}}`$ $`=`$ $`{\displaystyle \frac{L_{\mathrm{phys}}}{T_{\mathrm{phys}}}}{\displaystyle \frac{T_{\mathrm{mod}}}{L_{\mathrm{mod}}}},`$
where $`L`$ is the length, $`M`$ the mass, $`T`$ the time, $`V`$ the velocity and $`G`$ the gravitational constant. The model length unit is taken to be $`L_{\mathrm{mod}}=1`$ for the length of $`R_{\mathrm{core},1}`$ of galaxy 1. The corresponding physical length scale, $`L_{\mathrm{phys}}`$, is taken from the input data that specifies the length of this grid. The physical mass of the first galaxy, $`M_{\mathrm{phys},1}`$, is taken to be unity in model units, $`M_{\mathrm{mod},1}=1`$. The integration step-size, $`T_{\mathrm{phys}}`$, is specified by the user. The step-size in model units, $`T_{\mathrm{mod}}`$, follows from Eq. 19.
Converting factors for each grid of all galaxies are calculated to get the forces and accelerations, due to the FFT and the differentiation of the potentials requiring a cell-length of unity.
## 3 Conservation of Energy and Angular Momentum
In order to check how well energy and angular momentum are conserved, calculations with isolated Plummer-spheres in virial equilibrium with different particle numbers, time-steps and grid-resolutions are presented here. To have a non-vanishing angular momentum vector in $`z`$-direction for the isolated galaxy model, all particle orbits are taken to have the same direction in the $`xy`$-plane. This procedure does not change any other properties of the Plummer-model, but makes it possible to check for relative changes in angular momentum.
### 3.1 Total Energy
Calculating the correct potential energy of a stellar system requires summing all two-body interactions between all particles. This is quite impossible for $`N_\mathrm{p}\mathrm{¿}\mathrm{}10^6`$. By using the available grid-based approximation and adding up only the mean potential values of the cells for each particle, a useful estimate of the total potential energy is obtained. If, on the other hand, the deviation from the centre of the cell and contributions from the neighbouring cells were taken into account to obtain a more accurate estimate, then CPU-time would have to be used for an arguably purposeless endeavour. For simplicity, in Superbox the total potential energy is computed without any interpolation inside grid cells (as it was used for the acceleration); only the potential value of the particle’s cells are taken into account to compute the potential energy. The total potential energy at the $`n`$th time-step is thus
$`E_{\mathrm{pot}}^n`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N_\mathrm{p}}{}}}m_i\mathrm{\Phi }^n(x_i,y_i,z_i),`$ (20)
where $`\mathrm{\Phi }^n`$ is the sum of the potentials of all grids and galaxies used to calculate the force on particle $`i`$ at the $`n`$-th time-step, or in other words, all potentials which contribute to the leap-frog update of the velocity in Eq. 17. $`m_i`$ is the mass of the particle (equal for all particles belonging to the same galaxy).
The total kinetic energy is also calculated only approximately during runtime. This comes about because in the leapfrog integration-scheme the velocities are half a time-step behind the positions. Within our integration scheme the time-centred interpolation for the kinetic energy would be
$`E_{\mathrm{kin}}^n`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N_\mathrm{p}}{}}}m_i{\displaystyle \underset{j=1}{\overset{3}{}}}\left({\displaystyle \frac{v_{i,j}^{n+1/2}+v_{i,j}^{n1/2}}{2}}\right)^2,`$ (21)
where $`v_{i,j}^n`$ is the $`j`$-component of the velocity vector of particle $`i`$ at time-step $`n`$. In Superbox the philosophy is to keep the amount of storage as small as possible, so that only the current velocities are stored. To keep the error in $`E_{\mathrm{kin}}`$ as small as possible, particle kinetic energies are calculated before and after the velocity vectors are updated in one time-step. From these two values the arithmetic average is taken:
$`E_{\mathrm{kin}}^n{\displaystyle \frac{1}{4}}{\displaystyle \underset{i=1}{\overset{N_\mathrm{p}}{}}}m_i{\displaystyle \underset{j=1}{\overset{3}{}}}\left((v_{i,j}^{n1/2})^2+(v_{i,j}^{n+1/2})^2\right).`$ (22)
This still implies an additional error of the order of $`(\mathrm{\Delta }v)^2`$, but keeps storage low and the code fast.
To quantify the performance of Superbox in terms of conservation of total energy, $`E_{\mathrm{tot}}(t)=E_{\mathrm{pot}}(t)+E_{\mathrm{kin}}(t)`$, calculations with different time-steps, particle number and grid resolution are done. We consider the relative change $`\mathrm{\Delta }E_{\mathrm{tot}}(t)/E_{\mathrm{tot}}(t_0)`$, where $`\mathrm{\Delta }E_{\mathrm{tot}}(t)=E_{\mathrm{tot}}(t)E_{\mathrm{tot}}(t_0)`$, and $`t_0`$ is the reference time. It is chosen to be zero when the adjustment to equilibrium after setting up the discrete rendition of the galaxy is accomplished. Max$`(\mathrm{\Delta }E_{\mathrm{tot}}(t)/E_{\mathrm{tot}}(0))`$ denotes the largest deviation in energy that occurs during a computation (see Fig. 6).
Table 1 shows that deviations in $`E_{\mathrm{tot}}(t)`$ are less than 0.5 per cent, as long as the time-step, $`\mathrm{\Delta }t\mathrm{¡}\mathrm{}T_{\mathrm{cr}}/10`$, where $`T_{\mathrm{cr}}`$ is the half-mass crossing-time of the Plummer-model. We consider $`\mathrm{\Delta }tT_{\mathrm{cr}}/50`$ to be a safe choice. In theory (see Hockney & Eastwood 1981) the error of the leapfrog integrator should decrease as $`(\mathrm{\Delta }t)^2`$ for $`\mathrm{\Delta }t0`$, but as is evident from Table 1, other error sources become prominent at a sufficiently small time step.
Fig. 6 shows the time evolution of the change in total energy for a particular calculation. A linear drift fits the data well, which holds true in computations extending over 50 crossing-times. The slope of the energy-drift is 0.002 per cent per crossing-time. The dispersion in energy is artificial due to the crude calculation of potential energy. Over a Hubble time, which corresponds to $`150\times T_{\mathrm{cr}}`$ if the Plummer model galaxy is assumed to have $`T_{\mathrm{cr}}=10^8`$ yr, the resulting energy error is 0.3 per cent.
This linear drift is intrinsic to a PM code (see Hockney & Eastwood, their fig. 9.4), and its strength depends mainly on the choice of $`H_{000}`$ in the Green’s function. $`H_{000}`$ determines the strength of the gravitational interaction of particles in the same cell, and also the self-gravitation of a particle (see Section 2.1). Additional errors arise through the limited grid-resolution and the approximations adopted in evaluating the energies. This gives an additional oscillation around the value of the linear drift seen in Fig. 6.
Increasing the number of particles per galaxy to $`N_\mathrm{p}>10^5`$ leads to little further improvement in max$`(\mathrm{\Delta }E_{\mathrm{tot}}/E_{\mathrm{tot}})`$ with $`N=32`$. The decrease of the error levels off if the mean number of particles per cell becomes sufficiently large. However, changing the grid resolution from $`N=32`$ to 64 cells per dimension gives an improvement by a factor of at least 4, with the decrease in error still progressing for $`N_\mathrm{p}>10^6`$ (Fig. 7). This is due to the estimate of the true potential by our grid-based potential, and thus of the potential energy, improving with increasing $`N`$. It also shows that increasing the number of grid-cells while keeping the number of particles low (with $`H_{000}=4/3`$) does not improve energy conservation (e.g. $`N=64`$, $`N_\mathrm{p}=10^4`$).
### 3.2 Angular Momentum
To calculate the absolute value of the total angular momentum of the galaxy, $`L_{\mathrm{tot}}`$, the same technique as for the kinetic energy is applied. That is, the arithmetic mean of the two velocity values (before and after the velocity update) is calculated for each velocity component to obtain an estimate of the velocity vector that is in-phase with the position vector.
The time evolution of $`\mathrm{\Delta }L_{\mathrm{tot}}(t)/L_{\mathrm{tot}}(t)`$ for one particular calculation is shown in Fig. 8. As with $`E_{\mathrm{tot}}`$, computations over $`50T_{\mathrm{cr}}`$ show insignificant deviations from the linear drift. The slope of the drift is about 0.003 per cent per crossing-time.
Conservation of $`L_{\mathrm{tot}}`$ does not depend on $`\mathrm{\Delta }t`$ as long as the time-step is small enough ( $`\mathrm{\Delta }tT_{\mathrm{cr}}/50`$ as for energy in Section 3.1). But it is highly dependent on the grid-resolution. Changing from $`N=32`$ to 64 grids per dimension improves the change in angular momentum by at least a factor of 10. From Fig. 9 it can be seen that max$`(\mathrm{\Delta }L_{\mathrm{tot}}/L_{\mathrm{tot}})`$ is quite independent of $`N_\mathrm{p}`$.
## 4 Relaxation
Galaxies have two-body relaxation times of the order of $`10^7`$ Gyr, and show little or no relaxation over a Hubble time even in their inner parts. Therefore, a programme which models galaxies has to be collision-less. Since no particle-mesh code can be entirely free of relaxation, we have to check on which time-scales Superbox provides reliable results.
Following Standish & Aksnes (1969), the following experiment is performed: A number of equal-mass and fixed particles is distributed homogeneously inside a sphere of radius $`R_{\mathrm{sph}}`$ with a constant density distribution. A second group of particles is distributed on the surface of this sphere. They are allowed to move through the centre and leave the sphere on the opposite side. To make sure they leave the sphere, they are given an initial non-zero radial velocity component towards the centre of the sphere. The points where the moving particles leave the sphere are noted and the calculation is stopped after all moving particles have left the sphere. For every particle, its deflection angle $`\alpha `$ is obtained from
$`\mathrm{cos}\alpha `$ $`=`$ $`{\displaystyle \frac{\stackrel{}{R}_I\stackrel{}{R}_F}{|\stackrel{}{R}_I||\stackrel{}{R}_F|}},`$ (23)
where $`\stackrel{}{R}_I`$ and $`\stackrel{}{R}_F`$ are the vectors from the centre of the sphere to the initial and final position of the particle (see Fig. 10). The mean deflection angle, $`\overline{\alpha }`$, of all particles is computed and the relaxation time, $`T_{\mathrm{rel}}`$, is calculated from
$`T_{\mathrm{rel}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}90}{\mathrm{sin}\overline{\alpha }}}\overline{T}_{\mathrm{cr}}.`$ (24)
Here $`\overline{T}_{\mathrm{cr}}`$ is the mean time the moving particles require to travel through the sphere.
The number of moving particles is chosen to be $`10^5`$, which gives a sufficiently accurate estimate for $`T_{\mathrm{rel}}`$. The number of fixed particles, $`N_{\mathrm{fix}}`$, and the size of the inner mesh, $`R_{\mathrm{core}}`$, are varied. Grids with $`32`$ grid-cells per dimension are used. For a certain combination of $`N_{\mathrm{fix}}`$ and $`R_{\mathrm{core}}`$, 100 calculations with different random number seeds (both for the fixed and moving particles) are performed. In all cases, $`R_{\mathrm{sph}}<R_{\mathrm{out}}`$ (see also Fig. 4).
Fig. 11 shows the ratio $`T_{\mathrm{rel}}/\overline{T}_{\mathrm{cr}}`$ as a function of $`N_{\mathrm{fix}}`$. If one grid contains the whole sphere ($`R_{\mathrm{core}}=1.2R_{\mathrm{sph}}`$) a linear dependence of $`T_{\mathrm{rel}}/\overline{T}_{\mathrm{cr}}`$ on $`N_{\mathrm{fix}}`$ fits the results very well. Such a dependency is predicted for large particle numbers by standard relaxation theory (Spitzer & Hart 1971; Binney & Tremaine 1987). For $`N_{\mathrm{fix}}\mathrm{¿}\mathrm{}10^5`$, $`T_{\mathrm{rel}}/\overline{T}_{\mathrm{cr}}>10^3`$. Since $`T_{\mathrm{cr}}10^8`$ yr for a typical galaxy, the relaxation time is one order of magnitude larger than a Hubble time. Hence galaxies can be modelled with Superbox provided $`N_\mathrm{p}>10^5`$ particles per galaxy.
These results are not substantially changed if the particles have to cross a grid boundary. As an example, the case $`R_{\mathrm{core}}=0.5R_{\mathrm{sph}}`$ is shown in Fig. 11. Compared to the one-grid case, $`T_{\mathrm{rel}}`$ is reduced by about 20 per cent nearly independently of $`N_{\mathrm{fix}}`$. Hence, grid boundaries do not decrease $`T_{\mathrm{rel}}`$ significantly. The relaxation time also drops if the number of grid-cells is increased: Due to the better resolution, forces from particles in adjacent cells become larger, increasing the relaxation. As an example the case with 128 cells per dimension is shown in Fig. 11. The relaxation times are a factor of $`2.3`$ smaller compared to the case with $`N=32`$. Assuming that all particles within a cell are located at the cell centre, and treating the deflection of the moving particles as a pure N-body problem, we may obtain an analytical estimate of the relaxation in Superbox. A consideration similar to the one described in Binney & Tremaine (1987) p. 187 ff shows that the relaxation time should depend in the following way on $`N_{\mathrm{fix}}`$ and $`N_{\mathrm{gc}}`$:
$`T_{\mathrm{rel}}`$ $`=`$ $`\kappa {\displaystyle \frac{N_{\mathrm{fix}}}{\mathrm{ln}(\gamma N_{\mathrm{gc}})}}T_{\mathrm{cr}}`$ (25)
with $`\gamma =1/3`$. From the experiments described in this section we obtain a value for the constant of proportionality $`\kappa =0.05`$. Together with this value, Eq. 25 gives an estimate of the relaxation time of Superbox. From this result we can quantify, for a given grid resolution and computational time, the minimum number of particles that should be used to exclude any unwanted relaxation effects. It is again obvious that calculations with a high number of grid-cells should be done only with a sufficiently large particle number to ensure $`T_{\mathrm{rel}}>T_{\mathrm{Hubble}}`$.
## 5 Memory and CPU-Time Requirements
### 5.1 Memory
The memory requirement of Superbox scales both linearly with the particle number, $`N_\mathrm{p}`$, and with the number of grid-cells, $`N_{\mathrm{gc}}=N^3`$. For the particle array, 24 byte per particle (4 bytes per phase-space variable) are needed. The total amount of memory for different particle numbers and grid-resolutions is shown in Fig. 12.
The memory required for a Superbox calculation can be estimated from
$`\mathrm{memory}_{\mathrm{grid}}`$ $`=`$ $`\left(5N^3+N(2N)^2+(N+1)^3\right)\times 4\mathrm{byte},`$ (26)
$`\mathrm{memory}_{\mathrm{particles}}`$ $`=`$ $`\left(8N_\mathrm{p}+\mathrm{overhead}\right)\times 4\mathrm{byte}.`$ (27)
The first term on the right-hand side of Eq. 26 comes from the 5 grids (Section 2.3). The second term comes from the array needed for the FFT, and the last term comes from the array required to store the Fourier-transformed Green’s function. The memory usage of the grids is independent of the number of galaxies, since they are treated consecutively in the same grid-arrays. The ’overhead’ contains the arrays of the Superbox parameters (e.g. grid sizes), the centre of mass and density, and output data such as energies, angular momenta and Lagrangian radii, for each galaxy. Increasing the number of galaxies, while keeping the total particle and grid number constant, reduces the relative overhead contribution (Eq. 27) slightly due to these large arrays scaling only with the number of particles per galaxy.
### 5.2 CPU-time
The time-step cycle of Superbox is divided into three main routines (Fig. 1). Firstly, in Getrho the mass density of the grids is calculated. Secondly, the FFT routine computes the potential on the grids. Thirdly, the Pusher routine contains the force calculation, the position and velocity updating, and collection of the output data. These three routines need about 99 per cent of the total CPU time. Fig. 13 shows the amount of CPU-time per step, $`t_{\mathrm{CPU}}`$, needed for the different routines. All data are derived on a Pentium II 400MHz processor with Linux as the operating system and the GNU–g77 Fortran Compiler (with options $``$O3 $``$m486) for a model single galaxy.
For the Getrho and Pusher routines, $`t_{\mathrm{CPU}}N_\mathrm{p}`$, while for the FFT routine, $`t_{\mathrm{CPU}}N_{\mathrm{gc}}\mathrm{log}_{10}N_{\mathrm{gc}}`$ ($`N_{\mathrm{gc}}=N^3`$; Fig. 14), which is by far the dominant contribution. The resultant CPU time per step for galaxy $`i`$ is:
$`t_{\mathrm{CPU},i}`$ $`=`$ $`\alpha N_{p,i}+\beta \times 5N_{\mathrm{gc}}\mathrm{log}_{10}N_{\mathrm{gc}}.`$ (28)
As another example, for a Pentium 200MHz MMX, $`t_{\mathrm{CPU}}20`$ sec/step for $`N_\mathrm{p}=10^6`$ particles in one galaxy and $`N_{\mathrm{gc}}=64^3`$ grids.
The CPU-time needed for a computation with $`N_{\mathrm{gal}}`$ galaxies is derived by adding the individual times,
$`t_{\mathrm{CPU},\mathrm{tot}}={\displaystyle \underset{i=1}{\overset{N_{\mathrm{gal}}}{}}}t_{\mathrm{CPU},i}.`$ (29)
## 6 Comparison with other Codes
Demonstration of the conservation of energy and angular momentum is a necessary but not sufficient condition for validating any particle-mesh code. Additional validation of the code comes from an inter-comparison of results obtained with entirely different numerical techniques. Such a comparison is available in the study of the tidal dissolution of small satellite galaxies orbiting in an extended Galactic dark halo (Klessen & Kroupa 1998). Such an application of Superbox is rather extreme in that the jump in resolution from the middle grid (0.29 kpc/cell length) containing the small satellite to the outermost grid (25 kpc/cell length) containing the ’local universe’ is very large.
In that study, two Superbox computations are compared with calculations done using a direct but softened N-body integrator on the special hardware device named Grape 3 connected with a SUN-Ultra 20. The Superbox calculation uses $`N_{\mathrm{gc}}=32^3`$ grids, has $`3\times 10^5`$ satellite particles, and treats the live Galactic dark halo with $`10^6`$ particles. The Grape computation uses $`1.3\times 10^5`$ satellite particles with the Galactic dark halo being an analytic isothermal potential. Both calculations proceed for many thousands of time steps, corresponding to a time interval of many Gyr, until beyond disruption of the satellite. The results are in nice agreement, yielding essentially the same behaviour of the satellite, apart from expected (small) differences in exact disruption time. Of interest is also a comparison in CPU times required for the calculations, keeping in mind the very different number of particles used: with Superbox it takes 0.57 days on an IBM RISC/6000 workstation to compute 1 Gyr using $`1.3\times 10^6`$ particles, and it takes 0.36 days with Grape using $`1.3\times 10^5`$ particles (with an $`𝒪(N_\mathrm{p}^2)`$ scaling).
In the above calculations, the satellite mass is too small for the orbit to be affected by dynamical friction. A comparison of the orbital decay through dynamical friction of a massive satellite is of interest with analytical estimates, because this allows an altogether independent verification if Superbox treats collective behaviour on larger scales correctly. This is discussed in Section 7.
## 7 Dynamical Friction
An object moving through a homogeneous distribution of much lighter masses induces a trailing over-density. This over-density attracts the object which, as a result, is decelerated (Chandrasekhar 1943).
The orbital decay of a satellite galaxy or globular cluster orbiting within a larger galaxy is an astrophysical process relevant, for example, to the rate at which a satellite population is depleted and to the built-up of galactic bulges. In the Local Group, dynamical friction affects the orbit of the Magellanic Clouds, which is important for tracing their origin (see Westerlund 1997 and Kroupa & Bastian 1997 for summaries). It may also change the orbit of the Sagittarius dwarf spheroidal galaxy (Ibata & Lewis 1998; Gomez-Flechoso et al. 1999) and other nearby companions, depending on their masses. Understanding and proper application of dynamical friction is thus a very important issue.
The question if Chandrasekhar’s formula (Eq. 30 below) can be applied to orbits within finite and inhomogeneous density distributions has been the issue of a significant debate throughout the 1970’s and the 1980’s. The result is that for satellite masses that are relatively small compared to the large galaxy, the formula is a good approximation (Bontekoe & van Albada 1987; Zaritsky & White 1988; Velazquez & White 1999). Massive satellites induce non-local perturbations in the larger galaxy that are not taken account of in the derivation of Eq. 30, and in the collision of two galaxies of comparable mass analytical estimates become intractable, and the self-consistent numerical experiment must be resorted to (e.g. Madejski & Bien 1993). In order to do this we must, however, have reason to trust the numerical results.
To demonstrate the reliability of our results, we compare the orbital evolution of a satellite galaxy orbiting in an extended dark halo computed with the fully self-consistent Superbox-code, with the evolution resulting from the by now well-established Chandrasekhar approximation (Eq. 30).
### 7.1 Calculations with Superbox
For the numerical experiment a compact spherical satellite galaxy is injected into an extended isothermal dark parent halo. The numerical rendition of both systems is described in detail in Kroupa (1997), and only a short description is given here.
The parent halo is assumed to extend to $`R_\mathrm{c}=250`$ kpc with a core radius $`\gamma =5`$ kpc and a total mass $`M_{\mathrm{halo}}=2.85\times 10^{12}M_{}`$, corresponding to a circular velocity $`V_\mathrm{c}=220`$ km/s. A particle takes 588 Myr to cross the 33 per cent mass diameter. The inner, middle and outer grids are centred on the density maximum and have edge lengths of 50, 188 and 700 kpc, respectively. The model is allowed to relax to virial equilibrium, after which state it has a slightly more compact configuration with a circular velocity at a radius of 150 kpc of 245km/s, corresponding to a mass within that radius of about $`2.1\times 10^{12}M_{}`$. Due to a mild radial orbit instability the halo is also slightly prolate.
Plummer density distributions are taken for the satellites, each with a Plummer radius $`R_{\mathrm{pl}}=3`$ kpc and a cutoff radius of 15 kpc. The innermost and middle grids are centred on the satellite’s density maximum and have edge lengths of 10 and 40 kpc, respectively. The outer grid is the same as for the halo. Three satellite masses are used, $`M_{\mathrm{sat}}=1\times 10^{10},7\times 10^{10}\mathrm{and}\mathrm{\hspace{0.17em}5}\times 10^{11}M_{}`$. The respective crossing times of the 33 per cent mass diameter are $`t_{\mathrm{cr},33}=84`$, 32 and 12 Myr.
For the calculation, $`N=32`$ grid cells per dimension are used (only 28 being active, see Section 2.3), with $`10^6`$ particles in the parent halo, and $`3\times 10^5`$ satellite particles. For the halo the resolution is thus 1.79 kpc/cell-length for $`r25`$ kpc, 6.71 kpc/cell-length for 25 kpc $`<r94`$ kpc and 25 kpc/cell-length for 94 kpc $`<r350`$ kpc. For the satellite the resolution is 0.36 kpc/cell-length for $`r_\mathrm{s}5`$ kpc, 1.43 kpc/cell-length for 5 kpc $`<r_\mathrm{s}20`$ kpc and 25 kpc/cell-length for 20 kpc $`<r_\mathrm{s}350`$ kpc, where $`r_\mathrm{s}`$ is the distance from the satellite’s density maximum. The integration step size is $`\mathrm{\Delta }t=t_{\mathrm{cr},33}/75`$.
The satellite is allowed to relax to virial equilibrium before being placed into the parent halo at an initial radius $`r(0)=|𝐫(t=0)|=150`$ kpc with an initial velocity $`v(0)=|𝐯(t=0)|=245`$ km/s perpendicular to the radius vector, $`𝐫`$.
The satellite’s galactocentric distance is $`r(t)=|𝐫_{\mathrm{sat}}(t)𝐫_{\mathrm{halo}}(t)|`$, and its velocity is $`v(t)=|𝐯_{\mathrm{sat}}(t)𝐯_{\mathrm{halo}}(t)|`$, where $`𝐫_{\mathrm{sat}}(t)`$ and $`𝐫_{\mathrm{halo}}(t)`$ are the position vectors of the density maxima of the satellite and halo, respectively, and $`𝐯_{\mathrm{sat}}(t)`$ and $`𝐯_{\mathrm{halo}}(t)`$ are the velocity vectors of the satellite’s and halo’s density maxima, respectively.
### 7.2 Chandrasekhar friction
The equation of motion of a satellite in a stationary and rigid isothermal parent halo is
$`{\displaystyle \frac{\mathrm{d}^2𝐫}{\mathrm{d}t^2}}`$ $`=`$ $`a_\mathrm{g}𝐫+\eta 𝐯,`$ (30)
$`a_\mathrm{g}`$ $`=`$ $`{\displaystyle \frac{V_\mathrm{c}^2}{\gamma ^2+r^2}},`$ (31)
$`\eta `$ $`=`$ $`4\pi G^2\mathrm{ln}\mathrm{\Lambda }{\displaystyle \frac{\rho (r)M_{\mathrm{sat}}}{v^3}}\left[\mathrm{erf}(x){\displaystyle \frac{2x\mathrm{e}^{x^2}}{\sqrt{\pi }}}\right],`$ (32)
where $`a_\mathrm{g}`$ is the acceleration in an isothermal halo, and $`\eta `$ is the deceleration due to dynamical friction, in which $`G=4.499\times 10^3\mathrm{pc}^3/(M_{}\mathrm{Myr}^2)`$ is the gravitational constant, $`\rho (r)=V_\mathrm{c}^2/\left[4\pi \mathrm{G}\left(\gamma ^2+r^2\right)\right]`$ is the halo density and erf$`(x)`$ is the error function with $`x=v/V_\mathrm{c}`$. A derivation of Chandrasekhar’s dynamical friction formula can be found in Binney & Tremaine (1987).
The numerical value of the Coulomb logarithm, ln$`\mathrm{\Lambda }=b_{\mathrm{max}}/b_{\mathrm{min}}`$, is somewhat ill-defined. It is calculated by integrating particle deflections over impact parameters ranging from a minimum ($`b_{\mathrm{min}}`$) to a maximum ($`b_{\mathrm{max}}`$) effective value. The minimum impact parameter is taken here to be $`b_{\mathrm{min}}=2.7`$ kpc. The maximum impact parameter, however, is something like the distance over which $`\rho (r)`$ falls off significantly, and is thus less-well defined. An upper limit is $`b_{\mathrm{max}}r(t=0)=150`$ kpc. It can also be estimated by assuming $`\rho (r+b_{\mathrm{max}}^+)=0.5\rho (r)`$. With $`r=150`$ kpc, $`b_{\mathrm{max}}^+=62`$ kpc. If, on the other hand, $`\rho (rb_{\mathrm{max}}^{})=2\rho (r)`$, then with $`r=150`$ kpc, $`b_{\mathrm{max}}^{}=44`$ kpc. Once the satellite orbit decays to $`r=30`$ kpc, then $`b_{\mathrm{max}}^+=13`$ kpc, while $`b_{\mathrm{max}}^{}=9`$ kpc. The Coulomb logarithm thus lies in the range ln$`\mathrm{\Lambda }`$ 1 to 4, and the above argument suggests that it may be a monotonically decreasing function of declining $`r(t)`$.
Apart from the somewhat arbitrary value of the Coulomb logarithm, the analytic derivation of Eq. 32 assumes a Maxwellian velocity distribution in the halo, neglects the self gravity of the wake induced by the satellite’s gravitational focusing, and the motion of the deflected halo particles is assumed to be governed only by the satellite – the gravitational field of the parent halo is neglected. Nevertheless, Eq. 32 has been shown to provide reliable results provided $`M_{\mathrm{sat}}/M_{\mathrm{halo}}\mathrm{¡}\mathrm{}0.2`$ and $`\gamma <r(t)<R_\mathrm{c}`$ (Binney & Tremaine 1987 and references therein). The Superbox model described in Section 7.1 conforms with these restrictions.
For the comparison with Superbox, $`V_\mathrm{c}=245`$ km/s, $`\gamma =3`$ kpc, and $`M_{\mathrm{sat}}=1\times 10^{10},7\times 10^{10}`$ and $`5\times 10^{11}M_{}`$. Eq. 30 is rewritten to four coupled first-order differential equations, and integrated numerically using the Runga-Kutta method, with a sufficiently small step size to ensure stability of the solution.
### 7.3 Results
Since the Coulomb logarithm is ill defined, it is necessary to first calibrate the numerical solution to Eq. 30 using one Superbox calculation. Fig. 15 shows the decay of the orbit of the satellite with $`M_{\mathrm{sat}}=7\times 10^{10}M_{}`$ according to Superbox and Eq. 30 under different assumptions for ln$`\mathrm{\Lambda }`$.
For the solution shown as the lower-dotted-$`r(t)`$ curve in Fig. 15, ln$`\mathrm{\Lambda }=4`$. The Superbox result requires a smaller Coulomb logarithm. Tests show that ln$`\mathrm{\Lambda }=1.6`$ leads to good agreement. To show a possible upper limit to the solution of the equation of motion, the case ln$`\mathrm{\Lambda }=(0.049r(t)/b_{\mathrm{min}})`$ is plotted as the upper-dotted-$`r(t)`$ curve. For this case, ln$`\mathrm{\Lambda }=1`$ initially, and it decreases with $`r(t)`$. The figure also shows the time-varying distance along the x-axis between the density maxima of the satellite and the halo. The damped oscillation is nicely evident, and the period of the orbit is initially 3.2 Gyr. The lower panel of Fig. 15 plots the solutions for $`v(t)`$. The satellite in the Superbox calculation retains an approximately constant velocity, as is expected from the theory of dynamical friction, the solutions from which are plotted with different lines corresponding to the cases discussed above.
The overall conclusion is that the Superbox result is consistent with theoretical expectations, provided ln$`\mathrm{\Lambda }=1.6`$. That this finding extends to other satellite masses is shown in Fig. 16. In all cases, the satellites have lost less than 20 per cent of their mass by the end of the calculation.
It is notable that one single value for the Coulomb logarithm applies to the whole radial range as well as to all three satellites, suggesting that non-local perturbations of the extended halo do not significantly alter the orbit of the satellites. This finding is in nice agreement with the results for ln$`\mathrm{\Lambda }`$ by Bontekoe & van Albada (1987) and Velazquez & White (1999), and support their finding that Chandrasekhar’s approximation yields a useful description of orbital evolution if ln$`\mathrm{\Lambda }1.5`$ is used.
Further calculations with more grid cells and a more detailed study will be needed to show if the oscillations in $`r(t)`$ and $`v(t)`$ seen in Figs. 15 and 16, as well as in similar figures in the literature (e.g. fig. 3 in Bontekoe & van Albada 1987) are physical, i.e. if they are manifestations of the dependency of $`\eta `$ (Eq. 32) on $`r`$ and $`v`$.
## 8 Conclusions
Superbox is an unconventional particle-mesh code that uses moving sub-grids to track and resolve high-density peaks in the particle distribution. This extension avoids the limitation in spatial resolution encountered with standard particle mesh codes that employ only a single grid. The code is efficient in that the computational overhead is kept as slim as possible. The code is also memory efficient by using only one set of grids to treat galaxies in succession.
In this paper we have explained the algorithms used for force evaluation, grid positioning and orbit integration. The nearest-grid-point scheme (NGP) is used to compute the density on a grid. The potential is obtained by FFT, from which accelerations at the position of each particle within a cell are calculated using second-order central-difference quotients with linear interpolation to the position of particles inside that cell. Such a scheme has been identified as a special higher-order NGP scheme and is competitive with the cloud-in-cell (CIC) scheme regarding the precision and continuity of the forces across grid boundaries. Together with the nested sub-grids, this leads to precise and reliable force calculations at the location of each particle. The sub-grids are positioned at each time-step on the density maximum or at the centre of mass of a galaxy, the former method being the more useful. For orbit integration the leap-frog scheme is used. Conservation of energy and angular momentum is excellent, and numerical relaxation is significantly longer than a Hubble time for a galaxy with a crossing time of about $`10^8`$ yr and particle numbers exceeding $`10^5`$. CPU-timing and memory consumption are such that calculations of galaxy encounters involving many $`\times 10^6`$ particles are possible on present-day desk-top computers with 20MB RAM or more (see Fig. 12). On the other hand, it is clear that using the resources of massively parallel supercomputers will significantly increase the capabilities of Superbox.
The most serious limitation of Superbox lies in the inability to resolve newly formed self-gravitating systems, because the grid structure is decided on at the beginning of a computation. Superbox has presently no scheme to allow the inclusion of additional grids during a calculation. Such objects can form in tidal arms and may be the precursors of some dwarf satellite galaxies. This limitation will be reduced as the number of grids that can be used increases owing to advances in computer technology.
Current and future developments of Superbox would include for example incorporating an individual time-step scheme per grid and galaxy. We intend to parallelise the code to run on CRAY T3E super-computers. Furthermore, additional grid-levels will be introduced to reduce the jump in resolution between the present middle and outer grids, and/or to further increase the resolution of the central regions of a galaxy. Finally, the current sticky-particle algorithm of Superbox can be extented (e.g. to include star formation).
Problems that are now being tackled with Superbox include computations of the tidal interactions of satellite and disk galaxies (J.M. Penarrubia, in preparation), and the possible formation of spheroidal satellite galaxies from stellar super clusters (Fellhauer et al.). Fellhauer & Kroupa (2000) report on a study of the dynamical evolution of clusters of dozens to hundreds of young massive and compact star clusters (i.e. stellar super clusters) in a tidal field, in an attempt to understand the evolution of such objects observed in HST images of the Antennae galaxies (see e.g. Lançon et al. 2000). Also, the parameter survey of dwarf galaxies orbiting in a massive dark halo is being completed (Kroupa, in preparation), to identify regions in parameter space that may lead to dSph-like systems without dark matter. Furthermore, collapse calculations are being performed to study the properties of the resulting object after violent relaxation and secular relaxation thereafter (Boily, in preparation).
Acknowledgements
P.K. and H.B. acknowledge financial support by SFB 328 and C.B, by SFB439 (at the University of Heidelberg) from the German Science Foundation (DFG). Support and advice on computational aspects of HLRS Stuttgart is gratefully acknowledged.
Superbox is available at
```
ftp:\\ftp.ari.uni-heidelberg.de/pub/mike/super.tar.gz.
```
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# Robustly Unstable Eigenmodes of the Magnetoshearing Instability in Accretion Disk
## 1. Introduction
Over the last several years, the presence of magnetic fields in a differentially rotating plasma has been proposed as a possible mechanism of accretion disk turbulence and its associated large anomalous angular momentum transport inside the disk. The presence of magnetic fields in a shear rotating gas cylinder makes the gas unstable against axisymmetric perturbations(Velikhov (1959); Chandrasekhar (1961)). The normal mode analysis(Kumar, Coleman & Kley (1994)) of this local magnetoshearing instability showed the existence of unstable axisymmetric eigenmodes. The presence of this robust instability was re-recognized (Balbus & Hawley (1991)) and confirmed by nonlinear ideal MHD simulations(Hawley & Balbus (1991), 1992; Hawley, Gammie & Balbus (1995)). The observational features of astrophysical accretion disks point to the need for large viscosity much beyond the one collisional mechanisms can yield. This robust instability has been invoked (Hawley & Balbus (1992)) as a most promising candidate mechanism for the viscosity puzzle(e.g., Tajima & Shibata (1997)).
Numerical investigation of nonaxisymmetric magnetoshearing modes has been carried out by adopting the shearing coordinates(e.g., Balbus & Hawley (1992)). Matsumoto & Tajima (1995) analyzed nonaxisymmetric nonlocal eigenmodes which are sandwiched by two Alfvén singularities around the corotational point and are not influenced by the disk edge boundaries and grow exponentially in time. These modes are distinct from the modes discussed by Ogilvie & Pringle (1996), which are nonaxisymmetric modes contained within cylindrical boundaries and which depend strongly on the boundary conditions.
In this paper, we concentrate our focus on the validity of the linear analysis of nonaxisymmetric eigenmodes questioned by Ogilvie & Pringle(1996). The analysis is performed in the frame rotating with the local angular velocity, which is adopted in nonaxisymmetric mode analysis(Ogilvie & Pringle (1996); Matsumoto & Tajima (1995)), since eigenmodes evolve exponentially in time. The resolution of this question is important in theory of accretion disk Unless this magnetoshearing instability is a robust mode unaffected by the boundary conditions, the long search for candidate mechanisms of anomalous viscosity of accretion disks needs to be reopened. The criticism of Ogilvie & Pringle (1996) is interesting because it reflects the difficulty and the extreme mathematical and physical subtlety involved in the nature of this mode around the Alfvén singularity. The problem is marred by the non-self-adjointness of the differential equation that describes eigenmodes, arising from the presence of shear flows. We know of no systematic mathematical theory on non-self-adjoint differential equations. Thus it takes us a development of mathematical and physical theory of such a system in order to understand the argument by Ogilvie & Pringle and to respond aptly to it. In the end such analysis has been developed and, assuringly, we find that the magnetoshearing instability is both robust and insensitive to the boundary condition, as we thought originally.
In §2, we derive the wave equation in a differentially rotating magnetized disk based on the analysis of Matsumoto & Tajima (1995), and basic properties of the wave equation are discussed. We show that discrete nonaxisymmetric eigenmodes exist, which are buffeted by the pair of Alfvén singularities where the Doppler-shifted wave frequency equals the Alfvén frequency. Our analysis finds that the eigenmodes oscillate indefinitely in the vicinity of the Alfvén singular points when the eigenvalue is real, whereas the eigenmodes are regular when the eigenvalue is complex. Numerical calculation of eigenmodes with the Alfvén frequency and wavenumber dependence is discussed in §3. We compare the results with the local dispersion relation, and show these eigenmodes are discrete. Astrophysical implications and conclusions are discussed in §4.
## 2. Analytical Properties of Non-Self-Adjoint Equation near Alfvén Singularity
We consider the MHD stability of magnetoshearing modes in the co-rotating frame of the fluid. The basic ideal MHD equations in the frame rotating with angular velocity $`\mathrm{\Omega }`$ are
$`\left({\displaystyle \frac{}{t}}+(𝒗)\right)𝒗`$ $`=`$ $`{\displaystyle \frac{1}{\rho }}P+{\displaystyle \frac{\times 𝑩\times 𝑩}{4\pi \rho }}+𝒈`$ (1)
$`+`$ $`2𝒗\times 𝛀+(𝛀\times 𝒓)\times 𝛀,`$
$$\frac{𝑩}{t}=\times (𝒗\times 𝑩),$$
(2)
where $`𝒈`$ is the gravitational acceleration and $`𝒓`$ is the position vector. We assume incompressibility for simplicity,
$$𝒗=0.$$
(3)
We also ignore self gravity, which is not essential for the magnetoshearing instability.
We use the local Cartesian coordinates ($`x,y,z`$) in the rotating frame where the $`x`$-axis is in the radial direction, the $`y`$-axis in the azimuthal direction, and the $`z`$-axis parallel to $`\mathrm{\Omega }`$. The uniform velocity shear $`v_y=(3\mathrm{\Omega }/2)x`$ is assumed for the Keplerian disk, where $`x=0`$ is the local co-rotating radial position. The wave equation is derived by linearizing the basic equations around the equilibrium state and assuming solution of the form $`\stackrel{~}{\varphi }(x,t)\text{exp}[\text{i}(k_yy+k_zz)]`$. In the unperturbed state, the density, pressure and magnetic field are assumed to be uniform. The assumption of $`v_x=v_z=B_x=0`$ in the unperturbed state yields the unperturbed momentum equation,
$$𝒈+2𝒗_0\times 𝛀+(𝛀\times 𝒓)\times 𝛀=0.$$
(4)
Next, the Laplace transform of the perturbation, $`\overline{\varphi }(x,\omega )`$, is employed,
$$\overline{\varphi }(x,\omega )=_0^{\mathrm{}}𝑑t\stackrel{~}{\varphi }e^{\text{i}\omega t}.$$
(5)
Substitution of the Laplace transformed momentum and induction equations into the continuity equation yields (see Matsumoto & Tajima (1995) for detail) the initial value equation
$`{\displaystyle \frac{d^2\overline{v}_x}{dx^2}}+{\displaystyle \frac{3\mathrm{\Omega }\omega _A^2k_y}{\omega _D(\omega _D^2\omega _A^2)}}{\displaystyle \frac{d\overline{v}_x}{dx}}+[(k_y^2+k_z^2)`$ (6)
$`{\displaystyle \frac{9\mathrm{\Omega }^2k_y^2\omega _A^2}{2\omega _D^2(\omega _D^2\omega _A^2)}}+\mathrm{\Omega }^2k_z^2{\displaystyle \frac{\omega _D^2+3\omega _A^2}{(\omega _D^2\omega _A^2)^2}}]\overline{v}_x=\mathrm{\Gamma }(x,\omega ),`$
where $`\omega _D`$ is the Doppler-shifted frequency,
$$\omega _D=\omega +\frac{3}{2}\mathrm{\Omega }k_yx,$$
(7)
and $`\omega _A`$ is the Alfvén frequency,
$$\omega _A^2=\frac{(𝒌𝑩)^2}{4\pi \rho }=k_{}^2v_A^2.$$
(8)
The initial condition enters through the source function $`\mathrm{\Gamma }(x,\omega )`$.
The wave equation is derived by expressing the homogeneous part of equation (6) in terms of the normalized radial coordinate
$$\xi =\frac{3\mathrm{\Omega }k_y}{2\omega _A}x$$
(9)
as
$`{\displaystyle \frac{d^2\overline{v}_x}{d\xi ^2}}+{\displaystyle \frac{2\omega _A^3}{\omega _D(\omega _D^2\omega _A^2)}}{\displaystyle \frac{d\overline{v}_x}{d\xi }}+[{\displaystyle \frac{4}{9}}(1+{\displaystyle \frac{1}{q}})\left({\displaystyle \frac{\omega _A}{\mathrm{\Omega }}}\right)^2`$ (10)
$`{\displaystyle \frac{2\omega _A^4}{\omega _D^2(\omega _D^2\omega _A^2)}}+{\displaystyle \frac{4}{9}}{\displaystyle \frac{\omega _A^2}{q}}{\displaystyle \frac{\omega _D^2+3\omega _A^2}{(\omega _D^2\omega _A^2)^2}}]\overline{v}_x`$
$`D(\omega ,\xi )\overline{v}_x=0,`$
where the ratio of the squares of the azimuthal and vertical wavenumber is defined as
$$q=\frac{k_y^2}{k_z^2}.$$
(11)
Unstable eigenmodes may exist when the solution satisfies the boundary condition
$$\underset{|\xi |\mathrm{}}{lim}\overline{v}_x=0$$
(12)
in the upper half of the complex $`\omega `$-plane reference. This boundary condition makes our eigenmodes distinct from the modes found by Ogilvie & Pringle (1996), which are confined in rigid cylindrical boundaries and strongly dependent on the boundary condition. In order to have an overall angular momentum transport across the entire disk, it is imperative to have unstable modes within the disk, not just on the boundaries of the disk. To investigate the interior of the accretion disk, eigenmodes should not depend on the edge boundary conditions. The eigenmodes which arise from finite boundaries may contribute to the angular momentum only near the edge of the disk. Our mode, which is not affected by the edge, grows whenever the eigenfunction is located between two Doppler-shifted Alfvén points. Since the positions of those points are determined in co-rotating frame, the origin of the frame can be anywhere in the disk. The boundary condition then assures that the momentum transport resulted from the superposition of the growing eigenmodes which can occur throughout the accretion disk.
The proper boundary condition interior of the disk can be easily examined by inspecting the aysmptotic form of the basic differential equation (6) of the system. This indicates that the leading radial dependence of eq. (6) leads to the exponential decay away from the co-rotation point and toward the Alfvén singularity. This mathematics is, of course, most reasonable and physical as well, because the instability energy is provided within the two Alfvén singular layers to the mode, which dissipates the energy at(or near) the singularities.
Since the wave equation (10) is not self-adjoint due to the existence of the flow shear, the square of the eigenvalue $`\omega ^2`$ is not guaranteed to be real, which any self-adjoint system always satisfies. The fact that the eigenvalue is not pure real or imaginary but in general complex prevents us from applying the Strum-Liouville theory to this system. Note that if $`k_y=0`$, equation (6) is reduced to self-adjoint,
$$\frac{d^2\overline{v}_x}{dx^2}+k_z^2\left[1+\mathrm{\Omega }^2\frac{\omega ^2+3\omega _A^2}{(\omega ^2\omega _A^2)^2}\right]\overline{v}_x=0,$$
(13)
and its eigenvalue constitutes the Alfvén continuum. The analysis of this mode has been done by Chandrasekhar(1961). We now assume $`k_y0`$.
To the best of our knowledge, no theory for non-self-adjoint operators exists, and we have to investigate the properties of non-self-adjoint systems in general. We find that the eigenmodes possess certain symmetry properties, which are originated in the character of the differential operator and radial symmetry. Since the differential operator $`D(\omega ,\xi )`$ is invariant under the operation $`(\omega ,\xi )(\omega ,\xi )`$ because of radial symmetry, the eigenfunction $`\overline{v}_x(\omega ,\xi )`$ is also the eigenfunction of $`D(\omega ,\xi )`$,
$$D(\omega ,\xi )\overline{v}_x(\omega ,\xi )=0.$$
(14)
Note that the operation $`\xi \xi `$ is the same as changing the direction of the rotation $`\mathrm{\Omega }\mathrm{\Omega }`$. Taking conjugate of the equation (14) and denoting $`\xi `$ to $`\xi `$ yields
$`D(\omega ^{},\xi )\overline{v}_x^{}(\omega ,\xi )`$ (15)
$`=D(\omega ^{},\xi )\overline{v}_x(\omega ^{},\xi )=0,`$
where $`\varphi ^{}`$ means the complex conjugate of $`\varphi `$. Thus, if $`\omega `$ is an unstable eigenvalue of the equation (10), $`\omega ^{}`$ is another unstable eigenvalue whose eigenfunction $`\overline{v}_x(\omega ^{},\xi )`$ satisfies the relation (15). Especially, when $`\omega `$ is pure imaginary, i.e., $`\omega ^{}=\omega `$, the real part of eigenfunction is symmetric and the imaginary part antisymmetric with respect to $`\xi =0`$,
$$\overline{v}_x^{}(\omega ,\xi )=\overline{v}_x(\omega ,\xi ).$$
(16)
These properties of the non-self-adjoint operator indicate that this system has complex eigenvalues in general, which differs from a self-adjoint system, and one unstable eigenvalue has another unstable and two stable companions. This symmetry of non-self-adjoint system and the comparison between self-adjoint and non-self-adjoint eigenvalues in complex-$`\omega `$ plane is shown in Figure 1.
Next, we look for the solution of the wave equation (10). The boundary conditions for these ideal MHD modes we are interested in are that the eigenmode resides and is differentiable around the corotational point and sandwiched by a pair of the Alfvén singularities, and decays toward $`\xi \pm \mathrm{}`$. Note that these physical boundary conditions preclude the usual Kelvin-Helmholtz(K-H) instability eigenmodes (Tajima et al. (1991); Noguchi et al. (1998)). The boundary conditions generic to the K-H mode is to match the dissipative structure at the corotational point. The properties of these eigenmodes may be investigated by analyzing the solution around the spatial coordinates of interest by using the Frobenius expansion (Morse & Feshbach, (1953)), in particular around the Alfvén singular points. It is of particular significance to examine analytical properties of the Frobenius expression around the Alfvén singularities in order to examine the question leveled by Ogilvie & Pringle (1996) against the eigenmodes obtained by Matsumoto & Tajima(1995). An appropriate, compatible numerical method for these eigenmodes is, therefore, the shooting method that starts from an exponentially decaying functional form at $`\xi =\pm \mathrm{}`$ and shoots toward the corotating point where two sides of the function should smoothly (differentiably) match.(This will be closely examined in the next section.)
We expand $`\overline{v_x}`$ in equation (10) using,
$$\overline{v}_x=\underset{n=0}{\overset{\mathrm{}}{}}a_n(\xi \xi _{A\pm })^{n+s}$$
(17)
in the vicinity of the Alfvén singularities $`\omega _D=\pm \omega _A`$ or $`\xi =\xi _{A\pm }`$, where $`\xi _{A\pm }`$ is defined by
$`\xi _{A\pm }`$ $`=`$ $`\pm 1{\displaystyle \frac{\omega }{\omega _A}}`$ (18)
$`=`$ $`\pm 1\omega _r^{}\text{i}\omega _i^{},`$
and $`\omega _r^{}(\omega _i^{})`$ is the real(imaginary) part of $`\omega /\omega _A`$. Applying the Frobenius method in the vicinity of the regular singular points $`\xi =\xi _{A\pm }`$, the indicial equation for the exponent $`s`$ for Eq. (10) is given by
$$s^2+\frac{\omega _i^{}(\omega _i^{}\pm 3\text{i})}{23\text{i}\omega _i^{}\omega _{i}^{}{}_{}{}^{2}}s+\frac{4}{9q}\left[\frac{42\text{i}\omega _i^{}\omega _{i}^{}{}_{}{}^{2}}{(2\text{i}\omega _i^{})^2}\right]=0.$$
(19)
Note that the eigenfunction is irregular and oscillates indefinitely at $`\xi =\xi _{A\pm }`$ whenever $`s`$ has imaginary component, provided Re($`s`$) is nonpositive integer,
$$\overline{v}_x=(\xi \xi _{A\pm })^{\text{Re}(s)}\mathrm{exp}\left[\text{Im}(s)\mathrm{log}|\xi \xi _{A\pm }|\right].$$
(20)
The singular points are on the real axis if and only if the eigenvalue $`\omega `$ is real. In general, the indices are given by
$`s={\displaystyle \frac{\omega _i^{}(\omega _i^{}\pm ^{}3\text{i})}{2(2^{}\text{i}\omega _i^{})(1^{}\text{i}\omega _i^{})}}\times [1`$ (21)
$`\pm \sqrt{1{\displaystyle \frac{16}{9q}}{\displaystyle \frac{(4^{}2\text{i}\omega _i^{}\omega _{i}^{}{}_{}{}^{2})(1^{}\text{i}\omega _i^{})^2}{\omega _{i}^{}{}_{}{}^{2}(\omega _i^{}\pm ^{}3\text{i})^2}}}],`$
where the sign $`\pm ^{}`$ indicates that we take upper sign when $`\xi =\xi _{A+}`$, and lower sign when $`\xi =\xi _A`$, respectively. Since $`s`$ is complex, the solution is not analytic at the Alfvén singular points.
Now, we investigate some special cases. First, when the eigenvalue is pure real, i.e., $`\omega _i^{}=0`$, the indices are purely imaginary,
$$s=\pm \frac{2}{3\sqrt{q}}\text{i},$$
(22)
the eigenfunction rapidly oscillates and is indefinite. Note that the boundary condition in this case becomes special in that we need to shoot outward from $`\xi =\xi _c`$. We can no longer shoot from the outer to the inner region. As the eigenmodes oscillate indefinitely around the Alfvén singularity, we find that the eigenvalue becomes continuous in this pure real case to form the Alfvén continuum. Second, when $`\omega _i^{}1`$, which is similar to the first case but now the singular points are not on the real axis, the solutions are
$$s=\pm \left[\frac{2}{3\sqrt{q}}\text{i}^{}\frac{3}{32}\sqrt{q}\omega _i^{}\right],$$
(23)
whose solution oscillates limited times in the vicinity of the singular points on the real axis, but not indefinitely, and the solution is regular on the real axis. Third, when the perturbation is nearly toroidal ($`k_yk_z,q1`$), as in the accretion disks far from their source, and the eigenvalue is complex, $`s`$ is given by
$`s={\displaystyle \frac{\omega _i^{}(\omega _i^{}\pm ^{}3\text{i})}{2(2^{}\text{i}\omega _i^{})(1^{}\text{i}\omega _i^{})}}\times [1`$ (24)
$`\pm (1{\displaystyle \frac{8}{9q}}{\displaystyle \frac{(4^{}2\text{i}\omega _i^{}\omega _{i}^{}{}_{}{}^{2})(1^{}\text{i}\omega _i^{})^2}{\omega _{i}^{}{}_{}{}^{2}(\omega _i^{}\pm ^{}3\text{i})^2}})].`$
When the perturbation is pure toroidal, the second term in the square brackets in equation (24) vanishes and $`s`$ is given by
$$s=0,\frac{\omega _i^{}(\omega _i^{}\pm ^{}3\text{i})}{(2^{}\text{i}\omega _i^{})(1^{}\text{i}\omega _i^{})}.$$
(25)
The eigenfunction corresponding to $`s=0`$ is regular, and the latter one diverges logarithmically at the singular point, when $`\omega `$ is real. Finally, when the perturbation is nearly poloidal($`k_yk_z,q1`$), as may occur in accretion disks close to their source, we find $`s`$ as
$$s=\pm \frac{2\text{i}}{3\sqrt{q}}\frac{\sqrt{4^{}2\text{i}\omega _i^{}\omega _{i}^{}{}_{}{}^{2}}}{2^{}\text{i}\omega _i^{}},$$
(26)
which reduces to the roots of the first case when the eigenvalue is pure real.
The exponent $`s`$ at the corotation point $`\omega _D=0`$ (or $`\xi _c=\omega /\omega _A`$) is given by
$$s=\frac{1}{2(\omega _{i}^{}{}_{}{}^{2}1)}\left[\omega _{i}^{}{}_{}{}^{2}3\pm \sqrt{\omega _{i}^{}{}_{}{}^{4}14\omega _{i}^{}{}_{}{}^{2}+17}\right].$$
(27)
When $`\omega _{i}^{}{}_{}{}^{2}7+4\sqrt{2}`$, the eigenvalue is regular at the corotation point, since both indices are real and positive. When $`0\omega _{i}^{}{}_{}{}^{2}<1`$, indices are still real, and one of them is positive. This is the solution that is consistent with the matching condition at $`\xi =0`$ of the shooting method discussed in §3. When $`1<\omega _{i}^{}{}_{}{}^{2}7+4\sqrt{2}`$, the corotation point is singluar since two indices are real and negative in the region $`1<\omega _{i}^{}{}_{}{}^{2}74\sqrt{2}`$, and complex in the region $`74\sqrt{2}<\omega _{i}^{}{}_{}{}^{2}<7+4\sqrt{2}`$. When $`\omega _{i}^{}{}_{}{}^{2}=1`$, the eigenfunction has irregular singularity at the corotation point. Again, even though the eigenfunction is irregular at the corotation point, the physical eigenmodes on the real $`\xi `$-axis is regular. Moreover, since all the coefficients of differential equation (10) are real at the corotation point even when the eigenvalue is complex, the eigenfunction should be real at $`\xi =\xi _c`$.
Although irregular in the vicinity of the Alfvén singularities or the corotation point in most cases, eigenfunctions are not irregular in the physical sense unless the singularities are on the real axis. Instead, the oscillatory behavior and amplitude of the eigenfunction around the Alfvén singularities directly reflects the physical eigenfunction behavior on the real $`\xi `$-axis, especially if $`\omega _i^{}`$ is small, i.e., magnetic field is strong.
When the eigenvalue $`\omega `$ and the index $`s`$ are both complex, the eigenfunction oscillates indefinitely in the vicinity of the singularity due to the imaginary component of the index $`s`$(see eq. ). In this case, the physical eigenfunction on the real $`\xi `$-axis also oscillates very rapidly in the vicinity of the point which is the projection of the complex singular point to the real $`\xi `$-axis, but the physical eigenfunction oscillates only finite times because the projected point is not a singular point.
When the real component of $`s`$ is negative, the eigenfunction diverges at the singularity(eq. ). The amplitude of the physical eigenfunction on the real $`\xi `$-axis is large at the projected singular point on the real $`\xi `$-axis. However, since the projected point is not a singular point, the eigenfunction does not diverge at this point. It is also clear that the eigenvalue is regular even if the singular point is a branch point, since we can choose the branch cut of the eigenvalue without crossing the real $`\xi `$-axis.
If the eigenvalue is real, the singularities are on the real $`\xi `$-axis and eigenfunction is irregular in the physical sense. We will discuss pure real eigenvalue cases in §3.
## 3. Robustly Unstable Magnetoshearing Eigenmodes
The eigenvalues of the wave equation (10) are calculated numerically by the shooting method with the boundary condition discussed in §2. Since the equation (10) may have three singularities(corotation point and Alfvén resonances), we choose complex initial value to avoid Alfvén singularities on real axis, and integrate(”shoot”) the equation (10) on the real $`\xi `$-axis from the left and right asymptotic boundaries(which are for removed from the any of particular singular behavior) to the corotation point, where their value and first derivative of the eigenmode are to be matched. If they are not matched, we change the eigenvalue appropriately until we can match. This iterative method is generally called the ”shooting method” for eigenvalue problems. Spatial steps of the integration are smaller in the vicinity of the point where the Alfvén point is projected on the real $`\xi `$-axis than other regions, in case the eigenvalue is almost real but still complex. By using the Newton method to decrease errors of a trial function, we obtain higher accuracy and faster convergence than the previous shooting codes. This allows us to search for subtle singular and regular eigenfunctions over a wide range of parameter values.
In order to satisfy the boundary condition (12), we impose $`\overline{v}_x^{}/\overline{v}_x=k_\pm `$ at the numerical boundaries $`\xi =\pm 10`$ (corresponding to the artificial infinity), where $`k_\pm `$ are the negative and positive solutions of the quadratic equation given by inserting the functional form $`\overline{v}_x=\mathrm{exp}(k_\pm x)`$ into the equation (10). On the numerical boundaries, the leading term of the equation (6) is the first term of the coefficient of $`\overline{v}_x`$, which is of the order of 1, and other terms are of the order of 0.01 or lower in our calculation. Our assumption for the boundary condition is justified as far as these estimations are valid.
Figure 2 shows examples of eigenfunctions $`\overline{v}_x`$ obtained by our shooting code when $`\omega _A=0.1\mathrm{\Omega }`$ and $`q=0.01`$. The solid and dashed curves represent the real and imaginary parts of the eigenfunction respectively. Figure 2a is for the fundamental pure imaginary eigenvalue and Fig. 2b is for the complex eigenvalue. Since the eigenvalue of Fig. 2a is pure imaginary, the real part of the eigenfunction is symmetric and the imaginary part antisymmetric with respect to $`\xi =0`$, which is consistent with the equation (16). Figure 2b is the eigenfunction with a complex eigenvalue, which makes a pair with the eigenvalue $`\omega ^{}`$ whose eigenfunction is derived from the relation (15). These eigenfunctions are confined between two Alfvén singularities located at $`\xi =\xi _{A\pm }\pm 1`$, and they are real at $`\xi =\xi _c`$.
Figure 3 shows the distribution of eigenvalues in the upper complex $`\omega `$-plane when $`\omega _A=0.01`$ and $`q=0.01`$. It shows only the eigenvalues in the region Re($`\omega )0`$, and all the complex eigenvalues have a paired eigenvalue $`\omega ^{}`$ in the region Re($`\omega )<0`$. It is obvious that this non-self-adjoint system has complex eigenvalues, which never appear in a self-adjoint system (see Fig. 1). There are only two pure imaginary eigenvalues, which will be shown to merge by changing $`\omega _A`$ and $`q`$. We find that complex eigenvalues, which Matsumoto & Tajima(1995) did not find, exist and have smaller imaginary part and grow slower in time than the fundamental eigenmode.
Figure 4 shows the dependence of unstable eigenvalues on $`\omega _A`$ when $`q=0.01`$. The solid(dashed) curves show the imaginary(real) part of the eigenvalues. When $`\omega _A`$ is small, there exist two purely growing modes, which merge at $`\omega _A0.66\mathrm{\Omega }`$ and form complex eigenvalues. These modes were found in Matsumoto & Tajima (1995) and the qualitative properties of this mode are about the same as found in Matsumoto & Tajima. However, the merging point is slightly greater than the earlier value and, more significantly, the growth rate does not decay significantly, even beyond $`\omega _A=1.584\mathrm{\Omega }`$, where it was calculated to vanish in Matsumoto et al.. The fundamental mode acquires its maximum growth rate just before it merges with another(secondary) pure imaginary mode. In addition, we find two new complex modes. These complex eigenmodes are also shown in Fig. 4, and the growth rate of all these modes saturate to $`\omega 0.15\mathrm{\Omega }\text{i}`$, the same saturation value for the fundamental mode with increasing $`\omega _A`$.
Figure 5 shows the dependence of eigenvalues on $`q`$ when $`\omega _A=0.01\mathrm{\Omega }`$ (Fig. 5a) and $`0.66\mathrm{\Omega }`$ (Fig. 5b). The fundamental(F) and secondary(S) pure imaginary modes and two complex modes(1, 2) are shown in Fig. 5, which correspond to the eigenmodes in Fig. 4, when $`q=0.01`$. Eigenvalues calculated from local analysis are also shown in Fig. 5, which we discuss later. Two pure imaginary modes are always distinct when $`\omega _A=0.01\mathrm{\Omega }`$ (two upper modes in Fig. 5a). However, when $`\omega _A=0.66\mathrm{\Omega }`$, these two modes merge and become complex at $`\mathrm{log}_{10}q1.8`$ and split to become pure imaginary again at $`\mathrm{log}_{10}q1.4`$. The eigenvalues of the other two complex modes become pure imaginary when $`q`$ exceeds a certain value ($`\mathrm{log}q=1.25`$ for $`\omega _A=0.01\mathrm{\Omega }`$, $`\mathrm{log}q=0.75`$ for $`\omega _A=0.66\mathrm{\Omega }`$), and the growth rate for those modes saturates with increasing $`q`$.
Next, we compare the nonlocal eigenfunction results with the local(Fourier) dispersion relation. By replacing $`d/dx`$ in equation (6) with a constant $`\text{i}k_x`$ around $`x=0`$ and assuming that the unperturbed magnetic field is toroidal($`B_x=B_z=0`$), the local solution in the regime
$$|\omega |\omega _A\omega _I\sqrt{\mathrm{\Omega }^2k_z^2/(k_x^2+k_y^2+k_z^2)}$$
is (Matsumoto & Tajima 1995)
$$\omega ^2=\frac{3}{2}\left[1\frac{3}{2}q\pm \frac{3}{2}\sqrt{(q2)(q\frac{2}{9})}\right]\omega _A^2.$$
(28)
This local dispersion relation (28) shows that pure real eigenmodes appear in the region $`q<\frac{2}{9}`$, pure imaginary eigenmodes in $`q>2`$, and complex in $`\frac{2}{9}<q<2`$. However, when $`q`$ is small, i.e., the perturbation is almost parallel to the magnetic field, nonlocal eigenmodes are unstable in both $`\omega _A=0.01\mathrm{\Omega }`$ and $`0.66\mathrm{\Omega }`$ (see Figs. 5a, b), and the growth rate of each mode does not have strong dependence on $`q`$. We conclude that replacing $`/x`$ by a single wavenumber $`k_x`$ is invalid for these modes since such eigenmodes oscillate very rapidly in the vicinity of the Alfvén points in a pronounced fashion(see Fig. 2b). In other words, the spatial variation of the wavenumber in the radial direction is essential for the modal analysis of the magnetoshear instability. Radial dependency of the wavelength also prevents us from applying the WKB method to this model. The WKB method requires the wavelength of the eigenmodes $`L_e`$ is smaller than the shear scale length $`L_s(L_e/L_s1)`$, which may be satisfied around the corotational and Alfvén singularities, but the wavelength is comparable to the shear scale length in other regions$`(L_e/L_s1)`$.
To show that these eigenmodes are discrete, let us first show the existence of the Alfvén continuum on the real $`\omega `$-axis. The wave equation (10) has a solution for any real $`\omega `$ for which $`\omega _D^2=\omega _A^2`$ for some $`x`$. It follows that the spectrum of this mode is continuous, and the Alfvén continuum extends to the entire real $`\omega `$ by choosing some $`𝑩,k_y`$ and $`k_z`$, which is different from the model chosen by Ogilvie & Pringle (1996) in which the Alfvén continuum is restricted by the boundary condition. The eigenmodes with the pure imaginary eigenvalue that we have shown are obviously not in this class.
When $`𝑩=0`$, the Kelvin-Helmholtz modes can be derived, which are stable in accretion disks. In this limit, equation (6) reduces to
$$\frac{d^2\overline{v}_x}{dx^2}+\left[k_{}^2+\frac{\mathrm{\Omega }^2k_z^2}{\omega _D^2}\right]\overline{v}_x=0,$$
(29)
and when $`k_z=0`$, it has a simple solution $`\overline{v}_x=\mathrm{exp}[k_y|x|]`$, which satisfies the boundary condition (12). The first derivative of this class of solutions is discontinuous at $`x=0`$, which vanishes with introducing dissipation. The Kelvin-Helmholtz instability also has continuous eigenvalues, but this class of solutions is eliminated in our calculation because of the matching condition of the shooting method, which requires the eigenfunction and its first derivative to be continuous. We search for eigenvalues by choosing initial values in the region $`0<\omega _r/\omega _A<1`$ and $`0<\omega _i/\omega _A<1`$. We find that of the initial values converge to one of the eigenvalues in Fig. 3 and we conclude that all of the eigenmodes are discrete.
In Fig. 6 we show another eigenmode whose eigenvalue gradually becomes real with increasing $`q`$, when $`\omega _A=0.66\mathrm{\Omega }`$. However, when the eigenvalue is real, we have already shown that the index of the eigenfunction $`s`$ is pure imaginary in the vicinity of the Alfvén singularities (Eq. (22)) and that the eigenfunction has the form
$$\overline{v}_x=\mathrm{exp}\left[\text{i}|s|\mathrm{log}(\xi \xi _{A\pm })\right],$$
(30)
from equation (20). Such eigenfunctions oscillate indefinitely toward the singularity and the function can take any value between $`1<\overline{v}_x<1`$. This indicates that the function in the inner region $`\xi _A<\xi <\xi _{A+}`$ and outer region $`\xi <\xi _A,\xi <\xi _{A+}`$ is discontinuous at the Alfvén singularities. Thus the boundary condition for the continua cannot be that of shooting from the outside $`|\xi |=\mathrm{}`$ toward the inside, but we should shoot from inside toward the singularities.
Figure 7 shows an example of eigenfunction in the inner region when $`\omega _A=0.01\mathrm{\Omega }`$, $`\omega =0.002\mathrm{\Omega }`$ and $`q=0.01`$(see eq. (22) and arguments for detail of this mode). However, the eigenmode in Fig. 6 is continuous even at the Alfvén singular points, since the integration by a finite spatial step brings in an effective dissipation, which is not the case for the pure real eigenvalue. Instead, the eigenmode becomes continuous because of the numerical dissipation. Although this numerical eigenfunction is different from the theoretical eigenfunction beyond the passage of the singularity, the fact of continua remains the same for two different reasons. It should be pointed out that in real physical situation there always exists a dissipation even for a nearly ideal MHD system. The dissipation prevents the eigenmode from blowing up on the Alfvén singular points, keeping the energy of the eigenmode finite. The numerical dissipation affects this eigenmode in the same manner mathematically as the real dissipation does, by passing oscillations through the singularity barrier and avaraging oscillation in a finite spatial step. Thus the numerically obtained eigenmode, though different from theoretically expected continua, may be regarded as realistic.
Finally, we describe the physical behavior of the eigenmodes in accretion disks. The expression of $`\overline{v}_y`$ in terms of $`\overline{v}_x`$ is derived from the continuity equation(3),
$`\overline{v}_y={\displaystyle \frac{\text{i}}{1+q}}[{\displaystyle \frac{q}{k_z}}{\displaystyle \frac{}{\xi }}`$ (31)
$`{\displaystyle \frac{\omega _D\mathrm{\Omega }}{2(\omega _D^2\omega _A^2)}}(3{\displaystyle \frac{\omega _A^2}{\omega _D^2}}+1)]\overline{v}_x.`$
Figure 8a shows $`\overline{v}_y`$ calculated from the fundamental pure imaginary mode (Fig. 2a) and the velocity field created in the $`xy`$-plane by the fundamental eigenmode is shown in Fig. 8b. The eigenfunction of $`\overline{v}_y`$ is also trapped between two Alfvén singularities. Since both $`\overline{v}_x`$ and $`\overline{v}_y`$ are almost out of phase with each other, the velocity field created by the fundamental eigenmode consists of vortices in $`xy`$-plane which are the seeds of nonlinear instability(Matsumoto & Tajima (1995)). Note that since we use the frame rotating with angular velocity $`\mathrm{\Omega }`$, there is no unique origin $`x`$ in the $`xy`$-plane. Thus such unstable eigenmodes excited at various $`x`$-positions will overlap with each other to expand the unstable region in the $`x`$-direction.
## 4. Summary
We have shown the existence of the discrete unstable nonaxisymmetric magnetoshearing instability eigenmodes. Since we assumed the exponentially decaying boundary condition (Eq. ) for the radial component of velocity, our modes are independent of the effect of the boundary condition. In our accretion disk theory, this robust instability occurs without any unrealistic disk edge boundary condition in infinite linear shear flow. The scalelength of a single eigenmode in the $`x`$ direction $`\mathrm{\Delta }x`$ is determined by the local strength of magnetic field, the direction and amplitude of the wavenumber and the magnitude of the angular velocity (see Eqs. and ), since the eigenfunction is buffeted by the Alfvén singular points $`\omega _D=\pm \omega _A`$. If the magnetic field is pure toroidal, $`\mathrm{\Delta }x=2\omega _A/3\mathrm{\Omega }k_y=2v_A/3\mathrm{\Omega }2/3(v_A/C_s)H`$, where $`C_s`$ is the sound speed and $`H=C_s/\mathrm{\Omega }`$ is the thickness of the disk. When $`v_AC_s`$, the mode is localized in the radial direction with the scalelength smaller than the thickness of the disk. If the magnetic field has azimuthal component, $`\mathrm{\Delta }x`$ is propotional to $`k_{}/k_y`$(Matsumoto & Tajima (1995)). In both cases, our infinite boundary condition is sufficient if the scalelength of the eigenmode is smaller than $`H`$. The curvature of the magnetic field is also small if $`\mathrm{\Delta }xH`$. However, for nearly axisymmetric perturbations($`k_yk_{}`$), the eigenmode have a large radial scalelength, and the infinite boundary condition may not valid. Density gradient and geometrical effects become important in this case.
Our result of the growth rate of unstable modes agrees with that of Matsumoto & Tajima (1995) in the region $`\omega _A\mathrm{\Omega }`$. We have found complex eigenvalues with smaller growth rates than the fundamental pure imaginary eigenmode. When $`\omega _A`$ is larger than $`\mathrm{\Omega }`$, two pure imaginary eigenmodes merge, the results of which is the same as found in Matsumoto & Tajima. However, our result shows that the growth rate saturates with increasing $`\omega _A`$. This indicates that the accretion disk is unstable even if the Alfvén frequency is comparable to the angular velocity, a case of strong magnetic fields.
The comparison of the nonlocal and local dispersion relations demonstrates where and how the local Fourier mode approximation fails to be accurate for this nonaxisymmetric instability. The wavenumber dependence of the eigenvalue shows that the nonlocal modes are unstable even in the region $`q=k_y^2/k_z^21`$, where the local dispersion relation has only a stable solution. Furthermore, we have found that two pure imaginary eigenvalues merge to be complex and split into two pure imaginary again with increasing $`q`$ in a region where the solutions of the local dispersion relation are purely real. Overall, as Fig. 5 indicates, the discrepancy of the local theory from the correct nonlocal theory amounts to not just a quantitative level but a qualitative deviation.
We have developed the mathematical and physical theory of a system of nonself-adjoint differential equations for the first time. In astrophysics, nonself-adjointness always appears whenever there is a shear flow, which is mathematically unsolved so far. We find the relationship of complex eigenvalues, which a self-adjoint system does not have, by a general approach to the nonself-adjoint system. A pair of eigenvalues $`\omega `$ and $`\omega ^{}`$ relate to each other in our model, since our model is symmetric with respect to the $`x`$-axis. Even if there is no such symmetry, we conclude that four eigenmodes make up a group in the nonself-adjoint system in general.
Although our model of the nonaxisymmetric mode has been linear in Cartesian coordinates and ignored the effect of diffusion, analysis in Section 3 suggests how the eigenmode grows to enter a nonlinear stage, and how it explains momentum transport in accretion disks. Since there is no particular sense in the radial direction in Cartesian coordinates, there exists no specific calibration of momentum transport in the linear stage. The eigenmodes can be excited, however, between any pair of Alfvén singularities in the radial direction (the $`x`$-axis) and create vortices in the disk plane (the $`xy`$-plane), as shown in Fig. 8b. The eigenmode with the fundamental eigenvalue dominates in time for a given radial co-rotation point. For another (arbitrary) co-rotation point, the same applies. These eigenmodes can overlap with each other to form greater vortices. In this stage, the nonlinear effect gives rise to anomalous magnetic viscosity that underlies the momentum transport needed to explain astrophysical disks. Matsumoto & Tajima demosntrated non-linear evolution of the eigenmode by three-dimentional MHD simulation with the shearing-box model. The overlap of eigenmodes excited at various $`x`$s’ was shown. They also calculated the magnetic viscosity parameter
$$\alpha _B=\frac{\delta B_x\delta B_y}{4\pi \rho C_s^2}<\frac{\delta B^2}{4\pi \rho C_s^2}\frac{\omega _A^2}{k_{}^2C_s^2}<\frac{1}{(k_{}H)^2}$$
(32)
where the notation $`\delta B_x\delta B_y`$ etc., denotes the spatial avarage. They found that when the poloidal field is dominant, the magnetic viscosity is $`\alpha _BO(0.1)`$, which corresponds to $`\alpha `$ in dwarf novae during the bursting phase.
We conclude that the results of Matsumoto & Tajima are correct and that the robust mode in magnetized accretion disks is of the magnetoshearing origin. This mode should be dominant in nonlinear theory and our linear analysis supports the results from the three-dimensional simulation in Matsumoto & Tajima, which explained anomalous momentum transport in accretion disks.
The work was supported by the US Department of Energy and NSF ATM 98-15809. TT is currently on leave at LLNL.
## Appendix A Exisitence of the Localized Growing Mode in the Infinite Shearing Flow
Ogilvie & Pringle’s argument of the non-existence of localized growing modes (Ogilvie & Pringle (1996), Appendix C), is incorrect for the following reasons.
First, we rewrite the integrated wave equation (Ogilvie & Pringle (1996), \[C8\]) without introducing a parameter $`\lambda `$ for simplicity,
$$_{\mathrm{}}^+\mathrm{}\left[\left|\frac{dy}{dx}\right|^2+(k^2f)|y|^2\right]𝑑x=0,$$
(A1)
where $`k^2=k_y^2+k_z^2`$ and the eigenfunction $`y(x)`$ is assumed to decay exponentially as $`|x|\mathrm{}`$, as we assume for the boundary condition. They also assume a uniform magnetic field, which makes equation (A1) qualitatively equivalent to our wave equation (10). The function $`f(x)`$ is given by
$`f(x)={\displaystyle \frac{(k_yu^{})^2\omega _A^2}{[(i\omega _ik_yu^{}x)^2\omega _A^2]^2}}+k_z^2[{\displaystyle \frac{2\mathrm{\Omega }u^{}}{[(i\omega _ik_yu^{}x)^2\omega _A^2]}}`$ (A2)
$`+{\displaystyle \frac{4\mathrm{\Omega }^2(i\omega _ik_yu^{}x)^2}{[(i\omega _ik_yu^{}x)^2\omega _A^2]}}],`$
where $`u(x)`$ is the shearing velocity, and $`f(x)`$ has two Alfvén singular points in the lower half-plane. We show here that equation (A1) can be satisfied because of those two Alfvén singular points, which makes the integral (A1) equal to zero if the complex frequency is properly chosen. The positions of the Alfvén singular points are given by
$$x=\frac{\text{i}\omega _i\pm \omega _A}{k_yu^{}}$$
(A3)
where $`\omega _i`$ is the imaginary part of the eigenfrequency, and $`k_y`$ is the wavenumber parallel to the shear flow, respectively. Note that the total integrand of the equation (A1) becomes positive and decays to zero as $`|x|\mathrm{}`$ since $`f(x)`$ decays to zero as $`|x|\mathrm{}`$.
Now, if $`k^2>|f(x)|`$ for any $`x`$, the integrand is always positive and equation (A1) can not be satisfied with any eigenfrequency. However, if the singularity is near enough to or on the real axis, $`|f(x)|`$ becomes large enough to satisfy $`|f(x)|>k^2`$ in the vicinity of the singularities and the total integral can take any value. Note that the real part of the eigenfrequency does not affect the position of the singularity. Furthermore, if the eigenvalue is complex, as we calculated in §3, the integral can equal to zero without crossing any branch cut of the singularities. Therefore, localized growing modes should exist by choosing eigenfrequencies which makes the integral of the equation (A1) equal to zero. In other words, the contour of the integral (A1) is not free to pass around the circle in the upper half plane, since the complex conjugate of the eigenfunction $`y^{}`$ has two singularities in the upper half plane with branch cuts, which any half circle in the upper half plane with sufficient large radius should cross.
Thus the statement ”The integral relation (C8) therefore cannnot be satisfied” in line 25 of p.164 of Ogilvie & Pringle was incorrect, and our calculation is valid.
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# 1 Introduction
## 1 Introduction
The path integral formulation of quantum mechanics is quite subtle when applied to particles moving in a curved space . It can be used to evaluate anomalies in quantum field theories, but only when the corresponding quantum mechanical models are defined on a finite interval of the worldline. When viewed as one dimensional QFTs on the worldline (but with higher-dimensional target spaces), one is dealing with nonlinear sigma models with double-derivative interactions. Such theories are super-renormalizable: though certain one- and two-loop Feynman diagrams are superficially divergent and regularization is necessary. However, there is no need to renormalize infinities away because the infinities of different graphs cancel each other and quantum mechanics is finite. Different regularization prescriptions give in general different finite answers for the same Feynman diagram. This situation is rather familiar in QFT: it simply means that there are free parameters entering in the theory (which are equivalent to the ordering ambiguities of canonical quantization) that can only be fixed by requiring further constraints (finite renormalization conditions). The latter simply parametrize different physical phenomena which can be described by the quantum mechanical model under consideration.
It is sometimes claimed that one does not need any “artificial” counterterm at all because the theory has no divergences. As the results of this article show, in all regularization schemes studied so far one always needs finite local counterterms. In fact, finite local counterterms are in general to be expected because they just amount to finite additive renormalizations needed to implement the renormalization conditions.
In the recent past, two different regularization schemes for nonlinear sigma models on a finite time interval have been discussed carefully: mode regularization and time discretization . A detailed comparison carried out to three loops shows that both schemes produce the same physics . However, they both break manifest general coordinate invariance at intermediate stages and require noncovariant counterterms to restore that symmetry in the final result. It is important to stress that these counterterms are unambiguously determined in each scheme. Nevertheless, lack of manifest covariance is annoying and constitutes a technical limitation: at higher loops one must expand the non-covariant counterterms to get the corresponding vertices but one cannot employ covariant techniques to simplify that computation.
Recently, dimensional regularization has been employed to define a new regulated version of the path integral for an infinite time interval . By evaluating the partition function of a particular massive nonlinear sigma model with a one-dimensional flat target space, it was found that no noncovariant counterterms were needed to obtain the correct result. Since target space in was only one-dimensional, covariant counterterms could not be detected since these are proportional to the scalar curvature $`R`$. It is the purpose of this letter to extend the proposal of ref. to a higher dimensional target space and to demonstrate that a covariant counterterm is needed. This counterterm turns out to be $`V_{DR}=\frac{\mathrm{}^2}{8}R`$.
Let us present first a discussion on the limits of dimensional regularization applied to quantum mechanics as used in . The main problem is that it seems to require an infinite propagation time. In fact, one obtains a continuum momentum space (the energy in one dimension) only upon Fourier transforming the infinite time dimension. Integrals in momentum space are regulated dimensionally afterwards . Instead, it would be desirable to regulate and compute the path integral for a finite propagation time. The latter could be interpreted as a proper time, thus making it useful for relativistic applications in the world line approach to QFT . A related problem is that the infinite propagation time introduces infrared divergences in massless models, and requires a harmonic term as infrared regulator. In ref. only a massive model was considered. The harmonic term ruins general coordinate invariance: a potential of the form $`V\omega ^2g_{ij}(x)x^ix^j`$ is not a scalar since the coordinates $`x^i`$ do not transform as the components of a vector. Invariance in the final result could be recovered in the limit $`\omega 0`$ if the propagation time would be kept finite, but that limit is not possible in the dimensional regularization described above which requires an infinite propagation time. Given that general coordinate invariance is necessarily softly broken, one may use as well a potential $`V\omega ^2g_{ij}(0)x^ix^j`$ as infrared regulator. The latter is quadratic even far away from the origin of the chosen coordinate system and will not modify the interaction vertices. This soft breaking of general coordinate invariance is not expected to modify the counterterm $`V_{CT}`$ since such a counterterm is sensitive only to the ambiguities due to ultraviolet divergences.
We now proceed to test the proposal of ref. in a class of sufficiently general models and relate it to the other regularization methods mentioned above. The calculation in is enough to indicate that possible counterterms will be covariant, but since it involves a single coordinate it misses terms proportional to the curvature. Our strategy will be to compute terms in the effective action using both mode regularization (MR) and dimensional regularization (DR). Equating the results fixes the counterterm needed in dimensional regularization to be $`V_{DR}=\frac{\mathrm{}^2}{8}R`$.
First, let us briefly review some known facts. Quantization of a free particle on a curved space produces in the quantum Hamiltonian $`\widehat{H}`$ an undetermined term proportional to the scalar curvature, $`\widehat{H}=\frac{\mathrm{}^2}{2}\mathrm{\Delta }+\alpha \mathrm{}^2R`$. This is easily seen using canonical (operatorial) methods: ordering ambiguities are encountered in the construction of the quantum Hamiltonian from the classical one and give rise to terms with at most two derivatives on the metric. Then, requiring general coordinate invariance leaves only a term proportional to the scalar curvature. Using path integrals this arbitrary coupling will appear as a correction to the effective action proportional to the scalar curvature
$`0|e^{\beta \widehat{H}/\mathrm{}}|0={\displaystyle }𝒟x\mathrm{e}^{S[x]/\mathrm{}}=\mathrm{e}^{\mathrm{\Gamma }/\mathrm{}}=\mathrm{exp}\{{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _0^\beta }dt[\mathrm{}+(\alpha +{\displaystyle \frac{1}{12}})\mathrm{}^2R+\mathrm{}]]\}`$ (1)
where the first equality reminds us of the equivalence of canonical and path integral quantization ($`|0`$ and $`0|`$ are eigenstates of the position operator $`\widehat{x}`$ with eigenvalue zero) and in the second equality we have the definition of the effective action $`\mathrm{\Gamma }`$. The term $`\frac{\mathrm{}^2}{12}R`$ is partially due to the counterterm and partially due to two-loop diagrams, see eq. (36). Henceforth we set $`\mathrm{}=1`$.
We are going to compute the corrections to the effective action $`\mathrm{\Gamma }`$ as function of the various couplings using both mode and dimensional regularization. In the former we can be general and allow for a finite propagation time $`\beta `$. Then we take the limit $`\beta \mathrm{}`$, which is safe in the presence of an infrared regulator, and compare the result with dimensional regularization. It is known that the former requires the counterterm
$`V_{MR}={\displaystyle \frac{1}{8}}R{\displaystyle \frac{1}{24}}g^{ij}g^{kl}g_{mn}\mathrm{\Gamma }_{ik}^m\mathrm{\Gamma }_{jl}^n`$ (2)
to produce a general coordinate invariant result with $`\alpha =0`$ . We will see that dimensional regularization will match the result when using a counterterm
$`V_{DR}={\displaystyle \frac{1}{8}}R`$ (3)
which is manifestly covariant. For comparison we mention that the counterterm for time-slicing, needed to obtain the same result as mode regularization, is different, see eq. (40).
## 2 The 3-loop calculation with Riemann normal coordinates
The model we analyze is given by
$`S[x^i]={\displaystyle _{t_i}^{t_f}}𝑑t\left[{\displaystyle \frac{1}{2}}g_{ij}(x)\dot{x}^i\dot{x}^j+{\displaystyle \frac{1}{2}}\omega ^2g_{ij}(0)x^ix^j+V_{CT}\right]`$ (4)
where $`\omega `$ is a frequency needed as infrared regulator and $`V_{CT}`$ is the counterterm for the regularization scheme chosen. Using Riemann normal coordinates, we will need to compute up to three loops since the noncovariant part of the counterterm (2), when expanded around the origin of the coordinates, only gives contributions from 3 loops onwards. We want to make sure that noncovariant counterterms are not required when using dimensional regularization, as noticed in ref. . In the next section we shall repeat the calculation below for general coordinates but with only two-loop graphs. Since in general coordinates the first derivatives of the metric do not vanish, we get nonvanishing contributions from the noncovariant parts of (2) already at the two-loop level. This gives an additional nontrivial check on the covariance of the counterterm of dimensional regularization on the infinite time interval.
The counterterm is effectively of order $`\mathrm{}^2`$ since it first appears at two loops, but for notational convenience we are using units where $`\mathrm{}=1`$. As already mentioned, the harmonic potential breaks general coordinate invariance since it selects those coordinates in which the potential is quadratic. We have chosen them to be Riemann normal coordinates as a definition of our model, so that the metric has the expansion
$`g_{ij}(x)`$ $`=`$ $`\delta _{ij}+{\displaystyle \frac{1}{3}}R_{kijl}(0)x^kx^l+{\displaystyle \frac{1}{6}}_mR_{kijl}(0)x^kx^lx^m`$ (5)
$`+({\displaystyle \frac{1}{20}}_m_nR_{kijl}(0)+{\displaystyle \frac{2}{45}}R_{kipl}R_{mj}{}_{}{}^{p}{}_{n}{}^{}(0))x^kx^lx^mx^n+O(x^5).`$
We find it convenient to use a rescaled time parameter $`\tau `$ with $`t=\beta \tau +t_f`$ and $`\beta =t_ft_i`$, so that $`1\tau 0`$. An infinite propagation time will be recovered in the limit $`\beta \mathrm{}`$, while for finite $`\beta `$ this setting allows us to compare easily with the results for $`\omega =0`$ which were reported in using similar notations <sup>4</sup><sup>4</sup>4 Our conventions follow from $`[_i,_j]V^k=R_{ij}{}_{}{}^{k}{}_{l}{}^{}V_{}^{l}`$, $`R_{ij}=R_{ik}{}_{}{}^{k}_j`$. Thus, the scalar curvature $`R=R_i^i`$ of a sphere is negative..
With this rescaling, and introducing the ghost $`a^i,b^i,c^i`$ for a correct treatment of the measure , we aim to compute the following path integral with two different regularization schemes, mode regularization (MR) and dimensional regularization (DR),
$`{\displaystyle 𝒟x𝒟a𝒟b𝒟c\mathrm{e}^{\frac{1}{\beta }S}}`$ (6)
with
$`SS[x,a,b,c]={\displaystyle _1^0}𝑑\tau \left({\displaystyle \frac{1}{2}}g_{ij}(x)(\dot{x}^i\dot{x}^j+a^ia^j+b^ic^j)+{\displaystyle \frac{1}{2}}(\beta \omega )^2g_{ij}(0)x^ix^j+\beta ^2V_{CT}(x)\right)`$ (7)
and with the boundary conditions that all fields vanish at $`t=t_i,t_f`$, (i.e. at $`\tau =1,0`$).
For the perturbative evaluation (in the coupling constants contained in the metric $`g_{ij}(x)`$) it is convenient to split the action into a quadratic part $`S_2`$ and an interacting part $`S_{int}=S_3+S_4+S_5+S_6+\mathrm{}`$
$`S_2`$ $`=`$ $`{\displaystyle _1^0}𝑑\tau \left[{\displaystyle \frac{1}{2}}\delta _{ij}(\dot{x}^i\dot{x}^j+a^ia^j+b^ic^j)+{\displaystyle \frac{1}{2}}\delta _{ij}(\beta \omega )^2x^ix^j\right]`$ (8)
$`S_3`$ $`=`$ $`0`$ (9)
$`S_4`$ $`=`$ $`{\displaystyle _1^0}𝑑\tau \left[{\displaystyle \frac{1}{6}}R_{kijl}x^kx^l(\dot{x}^i\dot{x}^j+a^ia^j+b^ic^j)+\beta ^2V_{CT}\right]`$ (10)
$`S_5`$ $`=`$ $`{\displaystyle _1^0}𝑑\tau \left[{\displaystyle \frac{1}{12}}_mR_{kijl}x^kx^lx^m(\dot{x}^i\dot{x}^j+a^ia^j+b^ic^j)+\beta ^2x^i_iV_{CT}\right]`$ (11)
$`S_6`$ $`=`$ $`{\displaystyle _1^0}d\tau [({\displaystyle \frac{1}{40}}_m_nR_{kijl}+{\displaystyle \frac{1}{45}}R_{kipl}R_{mj}{}_{}{}^{p}{}_{n}{}^{})x^kx^lx^mx^n(\dot{x}^i\dot{x}^j+a^ia^j+b^ic^j)`$ (12)
$`+{\displaystyle \frac{\beta ^2}{2}}x^ix^j_i_jV_{CT}].`$
Note that all structures like $`R_{ijkl}`$, $`V_{CT}`$ and derivatives thereof are evaluated at the origin of the Riemann coordinate system, but for notational simplicity we do not indicate so explicitly from now on.
From $`S_2`$ one recognizes the propagators
$`x^i(\tau )x^j(\sigma )`$ $`=`$ $`\beta \delta ^{ij}\mathrm{\Delta }(\tau ,\sigma )`$
$`a^i(\tau )a^j(\sigma )`$ $`=`$ $`\beta \delta ^{ij}\mathrm{\Delta }_{gh}(\tau ,\sigma )`$ (13)
$`b^i(\tau )c^j(\sigma )`$ $`=`$ $`2\beta \delta ^{ij}\mathrm{\Delta }_{gh}(\tau ,\sigma )`$
where the functions $`\mathrm{\Delta }(\tau ,\sigma )`$, $`\mathrm{\Delta }_{gh}(\tau ,\sigma )`$ are to be defined shortly in each regularization scheme. Then, the transition element, eq. (1), at three loops is given by
$`𝒵=\mathrm{e}^\mathrm{\Gamma }=A\mathrm{exp}\left\{{\displaystyle \frac{1}{\beta }}(S_4+S_5+S_6)+{\displaystyle \frac{1}{2\beta ^2}}S_4^2_{con}+\mathrm{}\right\}`$ (14)
where the subscript $`\mathrm{𝑐𝑜𝑛}`$ refers to connected diagrams only. The constant $`A`$ is the normalization of the exact path integral for $`S_2`$ which describes a harmonic oscillator in $`D`$ dimensions
$`A=\left({\displaystyle \frac{\omega }{2\pi \mathrm{sinh}(\beta \omega )}}\right)^{\frac{D}{2}}.`$ (15)
For $`\omega =0`$ this term becomes the familiar Feynman measure for a free particle $`(2\pi \beta )^{D/2}`$. The perturbative contributions are obtained by computing the various Wick contractions. We record the results in terms of $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }_{gh}`$ through three loops; the symbol * denotes counterterms. The nonzero contributions are
$`{\displaystyle \frac{1}{\beta }}S_4=\text{}=I_2{\displaystyle \frac{\beta }{6}}R\beta V_{CT}`$ (16)
$`{\displaystyle \frac{1}{\beta }}S_6=\text{}=I_8\beta ^2\left({\displaystyle \frac{1}{20}}^2R+{\displaystyle \frac{1}{45}}R_{ij}^2+{\displaystyle \frac{1}{30}}R_{ijmn}^2\right)+I_9{\displaystyle \frac{\beta ^2}{2}}^i_iV_{CT}`$ (17)
$`{\displaystyle \frac{1}{2\beta ^2}}S_4^2_{con}=\text{}=I_{14}{\displaystyle \frac{\beta ^2}{36}}R_{ij}^2+I_{15}{\displaystyle \frac{\beta ^2}{24}}R_{ijmn}^2`$ (18)
being $`S_5`$ proportional to at least one classical field that is zero since $`x_i=x_f=0`$. The integrals $`I_n`$ are given by
$`I_2`$ $`=`$ $`{\displaystyle _1^0}𝑑\tau \left(\mathrm{\Delta }({}_{}{}^{}\mathrm{\Delta }_{}^{}+\mathrm{\Delta }_{gh}){}_{}{}^{}\mathrm{\Delta }_{}^{2}\right)|_\tau `$ (19)
$`I_8`$ $`=`$ $`{\displaystyle _1^0}𝑑\tau \left(\mathrm{\Delta }^2({}_{}{}^{}\mathrm{\Delta }_{}^{}+\mathrm{\Delta }_{gh}){}_{}{}^{}\mathrm{\Delta }_{}^{2}\mathrm{\Delta }\right)|_\tau `$ (20)
$`I_9`$ $`=`$ $`{\displaystyle _1^0}𝑑\tau \mathrm{\Delta }|_\tau `$ (21)
$`I_{14}`$ $`=`$ $`{\displaystyle _1^0}d\tau {\displaystyle _1^0}d\sigma (\mathrm{\Delta }|_\tau ({}_{}{}^{}\mathrm{\Delta }_{}^{}{}_{}{}^{2}\mathrm{\Delta }_{gh}^2)\mathrm{\Delta }|_\sigma 4\mathrm{\Delta }|_\tau {}_{}{}^{}\mathrm{\Delta }_{}^{}{}_{}{}^{}\mathrm{\Delta }\mathrm{\Delta }^{}|_\sigma `$ (22)
$`+2\mathrm{\Delta }|_\tau {}_{}{}^{}\mathrm{\Delta }_{}^{2}({}_{}{}^{}\mathrm{\Delta }_{}^{}+\mathrm{\Delta }_{gh})|_\sigma +2\mathrm{\Delta }^{}|_\tau \mathrm{\Delta }{}_{}{}^{}\mathrm{\Delta }_{}^{}\mathrm{\Delta }^{}|_\sigma +2\mathrm{\Delta }^{}|_\tau {}_{}{}^{}\mathrm{\Delta }\mathrm{\Delta }^{}\mathrm{\Delta }^{}|_\sigma `$
$`4\mathrm{\Delta }^{}|_\tau \mathrm{\Delta }{}_{}{}^{}\mathrm{\Delta }({}_{}{}^{}\mathrm{\Delta }_{}^{}+\mathrm{\Delta }_{gh})|_\sigma +({}_{}{}^{}\mathrm{\Delta }_{}^{}+\mathrm{\Delta }_{gh})|_\tau \mathrm{\Delta }^2({}_{}{}^{}\mathrm{\Delta }_{}^{}+\mathrm{\Delta }_{gh})|_\sigma )`$
$`I_{15}`$ $`=`$ $`{\displaystyle _1^0}d\tau {\displaystyle _1^0}d\sigma (\mathrm{\Delta }^2({}_{}{}^{}\mathrm{\Delta }_{}^{}{}_{}{}^{2}\mathrm{\Delta }_{gh}^2)+{}_{}{}^{}\mathrm{\Delta }_{}^{2}\mathrm{\Delta }^{}{}_{}{}^{2}2\mathrm{\Delta }{}_{}{}^{}\mathrm{\Delta }\mathrm{\Delta }^{}{}_{}{}^{}\mathrm{\Delta }_{}^{}).`$ (23)
We have kept the same names and notations for the integrals $`I_n`$ as in to facilitate comparison for the limit $`\omega 0`$ possible in mode regularization when $`\beta `$ is kept finite. We recall that $`\mathrm{\Delta }|_\tau \mathrm{\Delta }(\tau ,\tau )`$ and $`{}_{}{}^{}\mathrm{\Delta }\frac{}{\tau }\mathrm{\Delta }(\tau ,\sigma )`$ while $`\mathrm{\Delta }^{}\frac{}{\sigma }\mathrm{\Delta }(\tau ,\sigma )`$.
Let us first consider mode regularization. Here one expands all fields in a Fourier sine series and keeps all modes up to a large mode number $`M`$. The limit $`M\mathrm{}`$ is taken after having computed all integrals. In practice, one manipulates the integrals by partial integration to put them into a form which can be computed directly and without ambiguities in the continuum. One partially integrates such that all double derivatives of $`\mathrm{\Delta }`$, namely $`{}_{}{}^{}\mathrm{\Delta }_{}^{}`$ and $`\mathrm{\Delta }_{gh}{}_{}{}^{}\mathrm{\Delta }_{0}^{}`$ are removed. If this is not possible, one casts the expressions in a form such that the integrands vanish at the end-points. In the latter case, singularities like $`\delta (\tau )`$ and $`\delta (\tau +1)`$ are neutralized. With this prescription one recognizes that the function $`\mathrm{\Delta }(\tau ,\sigma )`$ appearing in the propagator is given by
$`\mathrm{\Delta }(\tau ,\sigma )={\displaystyle \underset{m=1}{\overset{M}{}}}\left[{\displaystyle \frac{2}{(\pi m)^2+(\beta \omega )^2}}\mathrm{sin}(\pi m\tau )\mathrm{sin}(\pi m\sigma )\right]`$ (24)
while as anticipated one can represent $`\mathrm{\Delta }_{gh}(\tau ,\sigma )={}_{}{}^{}\mathrm{\Delta }_{0}^{}(\tau ,\sigma )`$ with
$`\mathrm{\Delta }_0(\tau ,\sigma )={\displaystyle \underset{m=1}{\overset{M}{}}}\left[{\displaystyle \frac{2}{(\pi m)^2}}\mathrm{sin}(\pi m\tau )\mathrm{sin}(\pi m\sigma )\right].`$ (25)
Their continuum limit ($`M\mathrm{}`$) is given by
$`\mathrm{\Delta }(\tau ,\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{\beta \omega \mathrm{sinh}(\beta \omega )}}[\theta (\tau \sigma )\mathrm{sinh}(\beta \omega \tau )\mathrm{sinh}(\beta \omega (\sigma +1))+`$ (26)
$`+\theta (\sigma \tau )\mathrm{sinh}(\beta \omega \sigma )\mathrm{sinh}\left(\beta \omega (\tau +1)\right)]`$
$`\mathrm{\Delta }_{gh}(\tau ,\sigma )`$ $`=`$ $`\delta (\tau \sigma ).`$ (27)
It is easy to check that $`\left[_\tau ^2(\beta \omega )^2\right]\mathrm{\Delta }(\tau ,\sigma )=\delta (\tau \sigma )`$ and $`\mathrm{\Delta }(0,\sigma )=\mathrm{\Delta }(1,\sigma )=\mathrm{\Delta }(\tau ,0)=\mathrm{\Delta }(\tau ,1)=0`$.
Now, we can compute the various $`I_n`$ and obtain the results summarized in Table 1,
where we have found it convenient to define the function
$`I(a)={\displaystyle \frac{1a\mathrm{coth}(a)}{a^2}}.`$ (28)
As an example how these results are obtained, consider the “clover leaf” graph in (17) corresponding to $`I_8`$. Using that in mode regularization
$`({}_{}{}^{}\mathrm{\Delta }_{}^{}+\mathrm{\Delta }_{gh})|_\tau =_\tau (({}_{}{}^{}\mathrm{\Delta })|_\tau )(\beta \omega )^2\mathrm{\Delta }|_\tau ,`$ (29)
the first term in $`I_8`$ yields
$`{\displaystyle _1^0}𝑑\tau (4{}_{}{}^{}\mathrm{\Delta }_{}^{2}\mathrm{\Delta }(\beta \omega )^2\mathrm{\Delta }^3)|_\tau .`$ (30)
Hence
$`I_8={\displaystyle _1^0}𝑑\tau (5{}_{}{}^{}\mathrm{\Delta }_{}^{2}\mathrm{\Delta }(\beta \omega )^2\mathrm{\Delta }^3)|_\tau .`$ (31)
Then from eq. (26) we obtain
$`\mathrm{\Delta }|_\tau `$ $`=`$ $`{\displaystyle \frac{\mathrm{sinh}(\beta \omega \tau )\mathrm{sinh}(\beta \omega (\tau +1))}{\beta \omega \mathrm{sinh}(\beta \omega )}}`$ (32)
$`({}_{}{}^{}\mathrm{\Delta })|_\tau `$ $`=`$ $`{\displaystyle \frac{\mathrm{sinh}(\beta \omega (2\tau +1))}{2\mathrm{s}\mathrm{i}\mathrm{n}\mathrm{h}(\beta \omega )}}`$ (33)
and substitution into (31) yields the result for $`I_8`$ as given in Table 1.
In this regularization scheme the counterterm to be used is $`V_{MR}`$ as given in eq. (2). When evaluated at the origin of the Riemann normal coordinates it produces
$`V_{MR}={\displaystyle \frac{1}{8}}R`$ (34)
$`^i_iV_{MR}={\displaystyle \frac{1}{8}}^2R{\displaystyle \frac{1}{36}}R_{ijkl}R^{ijkl}.`$ (35)
As an aside, we can check the correctness of the $`\omega 0`$ limit. Since $`I(\beta \omega )\frac{1}{3}`$ for $`\omega 0`$, one can verify that the results in are reproduced
$`𝒵=A\mathrm{exp}\left\{\left[\beta {\displaystyle \frac{1}{12}}R+\beta ^2\left({\displaystyle \frac{1}{120}}^2R+{\displaystyle \frac{1}{720}}R_{ij}^2{\displaystyle \frac{1}{720}}R_{ijkl}^2\right)+\mathrm{}\right]\right\}.`$ (36)
This result is expected to be covariant and the use of Riemann normal coordinates shows immediately which is the covariant form of the effective action.
On the other hand, for $`\omega 0`$ and $`\beta \mathrm{}`$ one gets $`\beta I(\beta \omega )\frac{1}{\omega }`$, and thus
$`𝒵=A\mathrm{exp}\left\{\beta \left[{\displaystyle \frac{1}{12}}R+{\displaystyle \frac{1}{\omega }}\left({\displaystyle \frac{1}{40}}^2R+{\displaystyle \frac{1}{240}}R_{ij}^2{\displaystyle \frac{1}{240}}R_{ijkl}^2\right)+\mathrm{}\right]\right\}.`$ (37)
Now, this result in not expected to be covariant because of the presence of the mass term $`\omega `$. The apparent covariance of (37) is just a coordinate artifact of the Riemann normal coordinates (this point will be self-evident in the calculations of the next section). The result (37) is what one should obtain in dimensional regularization as well.
Thus, let us turn to dimensional regularization. The propagators are represented as in (13) with
$`\mathrm{\Delta }(\tau ,\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{dk}{2\pi }\frac{\mathrm{e}^{ik\beta (\tau \sigma )}}{k^2+\omega ^2}}`$ (38)
$`\mathrm{\Delta }_{gh}(\tau ,\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \frac{dk}{2\pi }\mathrm{e}^{ik\beta (\tau \sigma )}}.`$ (39)
Note that, strictly speaking, one should use an infinite $`\beta `$, which anyway cancels in (13), and a finite $`t\beta \tau `$ and $`s\beta \sigma `$. Now one can use dimensional regularization to compute the various integrals (with momenta contracted as suggested by the kinetic term continued to $`D`$ dimensions) and then take the limit $`D1`$. Using the formulas given in (and also in where dimensional regularization is used in configuration space), one recognizes that the ghosts are effectively regulated to give a vanishing contribution (this is due to the fact that $`\delta ^{(n)}(0)`$ is zero in dimensional regularization), while the remaining integrals give the results summarized in Table 2.
It is immediate to verify that the result (37) is reproduced once one uses the counterterm $`V_{DR}=\frac{1}{8}R`$ (of course, in the limit of infinite $`\beta `$ this result is unaffected by the infrared divergence related to the infinite time integral and remains finite). Thus, we conclude that $`V_{DR}=\frac{1}{8}R`$ is the counterterm needed in dimensional regularization to have $`\alpha =0`$ in eq. (1).
Of course, we could have compared as well dimensional regularization with time slicing regularization and obtain the same result. In that case, one should remember that time slicing (TS) requires different rules to compute the integrals in eqs. (19-23) but also a different counterterm
$`V_{TS}={\displaystyle \frac{1}{8}}R+{\displaystyle \frac{1}{8}}g^{ij}\mathrm{\Gamma }_{ik}^l\mathrm{\Gamma }_{jl}^k.`$ (40)
As an extra check, in what follows we also verify the necessity of the counterterm $`V_{DR}`$ at two loops but using arbitrarily chosen coordinates.
## 3 The two-loop calculation with general coordinates
In this section we repeat the calculation for the amplitude (1) using general coordinates going as far as two loops. Again we perform the calculation using mode regularization along with the counterterm (2) and dimensional regularization with the counterterm (3) applied to the model (4) where $`x^i`$ are now general coordinates. Writing (2) explicitly in terms of the metric tensor
$`V_{MR}`$ $`=`$ $`{\displaystyle \frac{1}{8}}R{\displaystyle \frac{1}{24}}g^{ij}g^{kl}g_{mn}\mathrm{\Gamma }_{ik}^m\mathrm{\Gamma }_{jl}^n={\displaystyle \frac{1}{8}}R{\displaystyle \frac{1}{32}}\left(_ig_{jk}\right)^2+{\displaystyle \frac{1}{48}}\left(_ig_{jk}\right)\left(_jg_{ik}\right)`$ (41)
makes it clear that one will get nonzero contribution from the noncovariant parts of the counterterms already at the two-loop level. Indeed the derivatives of the metric do not vanish at the origin of an arbitrary system of coordinates contrarily to what happens in Riemann normal coordinates where they do vanish. The expansion of the metric $`g_{ij}(x)`$ around the origin gives the same quadratic action of the previous section and thus the propagators are the same as well. The interacting part is $`S_{int}=S_3+S_4+\mathrm{}`$, being
$`S_3`$ $`=`$ $`{\displaystyle _1^0}𝑑\tau {\displaystyle \frac{1}{2}}_kg_{ij}x^k\left(\dot{x}^i\dot{x}^j+a^ia^j+b^ic^j\right)`$ (42)
$`S_4`$ $`=`$ $`{\displaystyle _1^0}𝑑\tau \left[{\displaystyle \frac{1}{4}}_k_lg_{ij}x^kx^l\left(\dot{x}^i\dot{x}^j+a^ia^j+b^ic^j\right)+\beta ^2V_{CT}\right]`$ (43)
where metric and derivative thereof and $`V_{CT}`$ are evaluated at the origin of the system of coordinates. The transition element at two-loop is given by
$`𝒵=A\mathrm{exp}\left\{{\displaystyle \frac{1}{\beta }}\left(S_3+S_4\right)+{\displaystyle \frac{1}{2\beta ^2}}S_3^2_{con}+\mathrm{}\right\}`$ (44)
where $`S_3`$ vanishes because it contains an odd number of quantum fields while
$`{\displaystyle \frac{1}{\beta }}S_4={\displaystyle \frac{\beta }{4}}\left[A_1^2g+2A_2^jg_j\right]\beta V_{CT}`$ (45)
$`{\displaystyle \frac{1}{2\beta ^2}}S_4_{con}={\displaystyle \frac{\beta }{8}}\left[B_1(_ig)^2+4B_2(_jg)g^j+2B_3(_ig_{jk})^2+4B_4(_ig_{jk})_jg_{ik}+4B_5g_j^2\right].`$ (46)
We have used the shorthand notation: $`^2gg^{ij}g^{kl}_k_lg_{ij}`$, $`_kgg^{ij}_kg_{ij}`$, $`g_kg^{ij}_ig_{jk}`$, $`^jg_jg^{ik}g^{jl}_k_lg_{ij}`$. The results obtained from this calculation are summarized in Table 3, where the column “Result” refers to the computations done using dimensional regularization (DR) and mode regularization (MR) of the integrals shown in the column aside. In the same line we also report a pictorial representation and the “tensor structure” associated to each diagram. Recalling that the scalar curvature is given by
$`R=^2g^jg_j{\displaystyle \frac{3}{4}}\left(_kg_{ij}\right)^2+{\displaystyle \frac{1}{2}}\left(_ig_{jk}\right)_jg_{ik}+{\displaystyle \frac{1}{4}}\left(_jg\right)^2\left(_jg\right)g^j+g_j^2`$ (47)
and using the results from Table 3, the amplitude $`𝒵`$ reads
$`𝒵=A\mathrm{exp}\left\{{\displaystyle \frac{\beta }{16}}^2g+{\displaystyle \frac{\beta }{8}}^jg_j+{\displaystyle \frac{\beta }{48}}(_ig_{jk})^2{\displaystyle \frac{\beta }{12}}(_ig_{jk})_jg_{ik}+{\displaystyle \frac{\beta }{8}}(_jg)g^j{\displaystyle \frac{\beta }{8}}g_j^2\right\}`$ (48)
for both regularization schemes. Therefore, also in this case, dimensional regularization yields the same transition amplitude as mode regularization only requiring the covariant counterterm $`\frac{1}{8}R`$.
Note that in Riemann normal coordinates $`^2g=\frac{2}{3}R`$ and $`^jg_j=\frac{1}{3}R`$ at the origin; substituting these identities, (48) reduces to the two-loop part of (36). Obviously the result is not covariant as the covariance of the model in eq. (4) is explicitly broken by the mass term and cannot be recovered even in the limit $`\omega 0`$ since (48) is $`\omega `$ independent. Therefore dimensional regularization of the path integral on an infinite time interval does not preserve target space general covariance, contrarily to what stated in .
## 4 Conclusions
In this letter we considered quantum mechanical path integrals in curved space with an infinite propagation time. We computed transition amplitudes both using dimensional regularization (DR) and other (in this context more established) regularization schemes. We showed that DR does not need noncovariant counterterms in order to reproduce the correct answer, as already noticed in a simpler model in , but it does need a covariant two-loop counterterm, namely $`V_{DR}=\frac{\mathrm{}^2}{8}R`$. We took an infinite propagation time in order to have a continuous momentum spectrum and to be able to use DR in the usual way. This forced us to add an infrared regulator: a mass term. The unpleasant feature of this term is that it breaks manifest general covariance. Furthermore, for applications to quantum field theories such as computations of anomalies, one needs path integrals on a finite time interval. We are at present working on an approach to use dimensional regularization at finite $`\beta `$. The crucial question is whether again only covariant counterterms are needed.
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# Mesonic fluctuations in a nonlocal Nambu–Jona-Lasinio model
## I Introduction
Until recently most calculations performed with the Nambu–Jona-Lasinio (NJL) model have been restricted to tree-level in the mesons or, in other words, to leading order (LO) in the $`1/N_c`$ expansion. Although the possible importance of quantum fluctuations of the mesons has been generally recognized (see, for example, Refs.), a satisfactory treatment of such effects has required the development of schemes that are faithful to the chiral invariance of the model. This is crucial, given that the most important motivation for studying the model is its description of the dynamical symmetry breakdown. Schemes have been introduced on the basis of an effective action method, an appropriate selection of Feynman diagrams and a bosonized approximation.
Calculations at next-to-leading order (NLO) in $`1/N_c`$ are certainly more involved than at LO, both analytically and numerically. However, since $`1/N_c`$ is such a modest expansion parameter, it has rightly been seen as important to estimate the size of some NLO effects. Even if the calculations were of interest for no other reason, a check on this perturbative approach would be very valuable. Encouragingly, there seems to be a general consensus emerging from recent studies that NLO corrections are indeed relatively small for standard choices of the model parameters, at least for properties of the ground state and of the pion. A contrary view, however, has been expressed by Kleinert and Van den Bossche.
The most compelling reason for a NLO analysis, however, is that the LO treatment of an NJL-type model neglects physical processes that are expected to be qualitatively important. For instance, several of the particles described by such models ($`\sigma `$, $`\rho `$, $`a_1`$) are broad resonances, yet the model meson propagators at LO are purely real. At NLO the particle widths are incorporated in a natural way, by including diagrams with two-meson intermediate states in the Bethe–Salpeter equation (BSE). Although these widths can be estimated from the LO decay amplitudes, the effects on the position of a resonance require a NLO treatment. Such effects may well prove important in model descriptions of, for example, the $`\rho `$ meson.
A difficulty with the model at NLO concerns the specification of model parameters associated with regularization of the loop integrals. At LO the model contains three parameters: the current quark mass, a coupling constant and a cutoff parameter that regularizes quark loops. Conventionally two of these parameters are fitted to the pion mass and decay constant, while the third is at least constrained by the value of the quark condensate. However, the regularization of the model must be further specified at NLO since meson and quark loop integrals can be regulated quite independently without affecting chiral symmetry relations. The parameter space can be limited to some extent because the NLO corrections produce instabilities for large values of the meson cutoff. Nonetheless the freedom to introduce another cutoff parameter produces significant uncertainty. Oertel et al. have determined the meson cutoff by studying the pion electromagnetic form factor and find values well away from the region of instabilities. This is a promising result, but to put it on a firmer basis one would have to relax assumptions that were made in Ref. for numerical reasonsThe $`\rho `$ meson was included in the model but vector and axial-vector degrees of freedom were not included as intermediate states in the NLO diagrams. In addition, $`\pi a_1`$ mixing was neglected, even though this is a LO effect.. An attractive alternative is to perform the regularization at the level of the model action rather than in the loop integrals. In this approach the regularization is specified from the outset and does not have to be imposed at each order of the expansion, a point which has been stressed by Ripka.
In this paper, we make a NLO analysis of a model that was originally proposed in Ref.. It has subsequently been studied at LO in the meson, baryon (soliton) and quark matter sectors. The two-body interaction vertex is nonlocal and separable, a Gaussian form factor being associated with each quark field. Such a model successfully eliminates several of the traditional problems of the NJL model whilst nevertheless retaining much of the simplicity that is its chief merit. The separable nature of the interaction is motivated in part by instanton liquid studies. The interaction form factors ensure the convergence of all loop integrals and so the NLO corrections are unambiguous. In addition, a practical advantage is conferred by the shape of the form factor, which allows complicated NLO diagrams to be evaluated efficiently with Gaussian numerical techniques. Such properties make this model a particularly convenient one with which to examine NLO effects. NLO quantum fluctuations of the quark condensate in a variant of this model were studied in Ref.. Here we investigate the NLO treatment of mesons.
In our approach we keep the full momentum dependence of all quark loop diagrams, and avoid expanding them in powers of meson momenta. Such an expansion may be useful in the determination of pion properties, but it is not valid for heavy mesons, such as the scalar, isoscalar sigma.
## II The nonlocal model
In this section we define the model used and recall some important aspects of its behaviour at LO. This is entirely standard and is included largely in order to establish our notation. The action of a quark model with a four-point interaction may be written
$$S=d^4x\overline{\psi }(x)(i\text{/}m_c)\psi (x)+\underset{i}{}\underset{n}{}d^4x_nH_i(x_1,x_2,x_3,x_4)$$
$$\times \overline{\psi }(x_1)\mathrm{\Gamma }_i^\alpha \psi (x_3)\overline{\psi }(x_2)\mathrm{\Gamma }_{i\alpha }\psi (x_4),$$
(1)
where $`\mathrm{\Gamma }_i^\alpha `$ denotes Dirac, colour and isospin matrices. Chiral symmetry constrains certain of the possible Dirac and isospin structures to appear in particular combinations. The admissible interaction terms are
$`H_1(11+i\gamma _5\tau ^ai\gamma _5\tau ^a),`$ $`H_2(\gamma ^\mu \tau ^a\gamma _\mu \tau ^a+\gamma ^\mu \gamma _5\tau ^a\gamma _\mu \gamma _5\tau ^a),`$ (2)
$`H_3(\gamma ^\mu \gamma _\mu ),`$ $`H_4(\gamma ^\mu \gamma _5\gamma _\mu \gamma _5),`$ (3)
$`H_5(\tau ^a\tau ^a+i\gamma _5i\gamma _5),`$ $`H_6(\sigma _{\mu \nu }\sigma ^{\mu \nu }\sigma _{\mu \nu }\tau ^a\sigma ^{\mu \nu }\tau ^a).`$ (4)
We shall be primarily concerned with the simplest version of the model, including the pions, and their chiral partner, through the $`H_1`$ interaction. In the usual NJL model the $`H_i`$ are simply constants. However, the use of a nonlocal interaction is appealing, not least because it enables one to avoid some of the ambiguities associated with regularization. Moreover, a nonlocal interaction of separable form is suggested by studies of the instanton-liquid picture of the QCD vacuum. In this case, the momentum space form of the $`H_i`$ is
$$H_i(p_1,p_2,p_3,p_4)=\frac{1}{2}(2\pi )^4G_if(p_1)f(p_2)f(p_3)f(p_4)\delta (p_1+p_2p_3p_4),$$
(5)
with the $`G_i`$ being constants which, for the purposes of counting, we take to be $`𝒪(1/N_c)`$. Although the instanton-liquid model predicts a particular form for the form-factor function $`f(p)`$, other authors have considered a function which is a Gaussian in Euclidean space,
$$f(p_E)=\mathrm{exp}(p_E^2/\mathrm{\Lambda }^2).$$
(6)
The Gaussian shape produces a good phenomenology of mesons and nucleons in the LO approximation and is used for the numerical results in this paper.
At leading order in the $`1/N_c`$ expansion the quark propagator is dressed by a single quark-loop self-energy diagram. This results in a momentum-dependent quark “mass”, $`m(p)`$:
$$S^1(p)=\text{/}pm(p);$$
(7)
$$m(p)=m_c+\left(m(0)m_c\right)f^2(p),$$
(8)
where the constant $`m(0)`$ is obtained from the Schwinger–Dyson equation (SDE). If $`m(0)`$ is sufficiently large, the poles in the quark propagator are shifted into the complex $`p^2`$ plane. A feature of models with nonlocal interactions is the appearance of additional, unphysical poles in the quark propagator. For the parameter sets of most interest these are located well away from the real axis, but instabilities of the model ground state can appear for small values of the cutoff.
Mesonic bound states can be found from the poles in the $`\overline{q}q`$ scattering matrix, $`T`$. For a separable interaction, it is convenient to factor out interaction form factors of the external quark momenta and to define
$$T(p_1,p_2,p_3,p_4)=\delta (p_1+p_2p_3p_4)\underset{n}{}f(p_n)\underset{i,j}{}\mathrm{\Gamma }_i\widehat{T}_{ij}(q)\overline{\mathrm{\Gamma }}_j,$$
(9)
where $`q=p_1p_3=p_4p_2`$ denotes the total momentum of the $`\overline{q}q`$ pair. The notation $`\overline{\mathrm{\Gamma }}_j`$ is needed in the longitudinal components of the vector and axial channels where, for example, $`\mathrm{\Gamma }_V=\overline{\mathrm{\Gamma }}_V=i\text{/}q`$. In all other cases $`\mathrm{\Gamma }_j`$ and $`\overline{\mathrm{\Gamma }}_j`$ are the same (see Eq. (23) of Ref.). The Bethe–Salpeter equation (BSE) at LO is given by
$$\widehat{T}(q)=G+GJ(q)\widehat{T}(q),$$
(10)
where $`G`$ is simply a matrix of the coupling constants from the action and $`J(q)`$ is the matrix of polarization loop integrals
$$J_{ij}(q)=i\text{Tr}\frac{d^4p}{(2\pi )^4}f^2(p_+)f^2(p_{})\mathrm{\Gamma }_jS(p_{})\overline{\mathrm{\Gamma }}_iS(p_+),$$
(11)
where we have introduced $`p_\pm =p\pm \frac{1}{2}q`$. The symbol ‘Tr’ is used to denote a trace over flavour, colour and Dirac indices.
Close to the pole corresponding to a particular particle, the amplitude in the relevant channel may be written
$$\frac{\overline{V}(q)V(q)}{m^2q^2},$$
(12)
where $`V(q)`$ denotes the particle vertex function. In the simplest form of the model, with only the $`G_1`$ coupling included, these functions are:
$$V_\pi (q)=g_{\pi qq}i\gamma _5\tau ^a,V_\sigma (q)=g_{\sigma qq},$$
(13)
where:
$$\frac{1}{g_{iqq}^2}=\frac{dJ_{ii}}{dq^2}|_{q^2=m^2}.$$
(14)
The vertex functions appear in calculations of mesonic properties, such as the pion decay constant, $`f_\pi `$.
In order to obtain symmetry currents with the same divergences as the corresponding local currents in QCD, and hence to respect the corresponding Ward identities, one has to include additional nonlocal terms in the currents. These arise as a consequence of the nonlocality of the action. Details of the nonlocal terms required in both vector and axial currents can be found in Refs.. Although the method developed in those papers is not unique and there is an ambiguity in defining the transverse parts of the currents, this does not affect the longitudinal components.
As an example, consider the pion decay constant. The diagrams in Fig. 1 represent the contributions from the local and nonlocal pieces of the axial current. The longitudinal part of the current, $`q^\mu A_\mu ^a`$, contains the following term
$`{\displaystyle \frac{i}{2(2\pi )^{12}}}G_1{\displaystyle }`$ $`{\displaystyle \underset{n}{}}d^4p_n\overline{\psi }(p_1)i\gamma _5\tau ^a\psi (p_3)\overline{\psi }(p_2)\psi (p_4)\delta (p_1+p_2+qp_3p_4)`$ (17)
$`\times [f(p_1)f(p_2)f(p_3)f(p_4q)+f(p_1)f(p_2)f(p_3q)f(p_4)`$
$`f(p_1)f(p_2+q)f(p_3)f(p_4)f(p_1+q)f(p_2)f(p_3)f(p_4)].`$
This term contributes to the pion decay constant since the operator $`\overline{\psi }(p_2)\psi (p_4)`$ has a nonzero vacuum expectation value (represented by the bubble loop in the second diagram of Fig. 1). The rest of the operator, $`\overline{\psi }(p_1)i\gamma _5\tau ^a\psi (p_3)`$, contributes to the pion-to-vacuum matrix element of the current. Such terms are needed in order to satisfy the Gell-Mann–Oakes–Renner (GMOR) relation, as discussed in Ref. , and they can make significant numerical contributions to various observables.
## III NLO contributions
### A Symmetry currents
At leading order in the $`1/N_c`$ expansion the fields at $`x_1`$ and $`x_3`$ in the interaction term of the action (1) are always connected to the same quark loop, and similarly for the fields at $`x_2`$ and $`x_4`$. Beyond this order, there are also exchange or “Fock” terms where the fields at $`x_1`$ and $`x_4`$ are connected to the same loop. These can be constructed by first performing a Fierz transformation on the action and then using the transformed action in just the same way that one uses the original action at LO. The Fierz transformed action of this model contains the following colour-singlet terms:
$`(4N_c)^1(G_12G_3+2G_4G_5+12G_6)(11+i\gamma _5\tau ^ai\gamma _5\tau ^a),`$ (18)
$`(4N_c)^1(2G_2+G_3+G_4)(\gamma _\mu \tau ^a\gamma ^\mu \tau ^a+\gamma _\mu \gamma _5\tau ^a\gamma ^\mu \gamma _5\tau ^a),`$ (19)
$`(4N_c)^1(2G_1+6G_2+G_3+G_42G_5)(\gamma _\mu \gamma ^\mu ),`$ (20)
$`(4N_c)^1(2G_1+6G_2+G_3+G_4+2G_5)(\gamma _\mu \gamma _5\gamma ^\mu \gamma _5),`$ (21)
$`(4N_c)^1(G_12G_3+2G_4+G_512G_6)(\tau ^a\tau ^a+i\gamma _5i\gamma _5),`$ (22)
$`(8N_c)^1(G_1G_54G_6)(\sigma _{\mu \nu }\sigma ^{\mu \nu }\sigma _{\mu \nu }\tau ^a\sigma ^{\mu \nu }\tau ^a).`$ (23)
(There are also colour-octet terms which we do not use here.)
The nonlocal terms in the symmetry currents of the model also involve four quark fields and so will also be subject to exchange effects at NLO. These Fock pieces in the currents will introduce further ambiguity through the definition of their transverse parts. One way to construct the NLO current terms is to take the Fierz transformed action and to apply the same method as was used in Refs. to form the nonlocal currents from the original model action. This method leads to nonlocal NLO current terms of the same structures as those presented in Ref., with the appropriate coupling constants in the original terms replaced by the corresponding combinations of couplings from the Fierz transformed action. For example, as well as a term $`G_1(i\gamma _5\tau ^a1)`$ in the nonlocal axial current constructed from the original action (see Eq. (17)), there is a Fock term of the same structure where $`G_1`$ is replaced by $`(4N_c)^1(G_12G_3+2G_4G_5+12G_6)`$.
An alternative approach to constructing the Fock terms of the model’s currents is simply to Fierz transform the LO currents of Ref.. The differences between this procedure and the one outlined above (where the Noether-like method for deriving a symmetry current is applied after the Fierz transformation) lie in the combinations of form factors which are present in the currents. However, these affect only the transverse parts of the currents; both definitions satisfy the appropriate Ward identities. The equivalence of the longitudinal components has been checked explicitly. In practice, the currents obtained from the Fierz-transformed action are somewhat easier to work with, since they retain the same structures as the nonlocal terms in the LO currents, with appropriate substitutions for the overall coefficients.
### B Quark propagator
We are now in a position to consider the quark and meson propagators at NLO. Like Oertel et al. and Ripka we use a strict NLO scheme in which all propagators and matrix elements are expanded to NLO in $`1/N_c`$. This scheme has also been presented by Dmitras̆inović et al., who themselves preferred to adopt an alternative scheme where certain NLO terms in the quark self-energy are treated self-consistently in the quark Schwinger–Dyson equation. Their scheme resums a subset of terms to all orders in $`1/N_c`$. Both schemes satisfy chiral low-energy theorems. The Feynman diagrams required by the strict NLO scheme in the nonlocal NJL model are briefly presented below.
At NLO the quark Schwinger–Dyson equation is coupled to the meson Bethe–Salpeter equation and takes the form
$`S_F^1(p)`$ $`=`$ $`\text{/}pm_ciG_1f^2(p)\mathrm{Tr}{\displaystyle \frac{d^4k}{(2\pi )^4}S_F(k)f^2(k)}`$ (25)
$`+if^2(p){\displaystyle \underset{i,j}{}}{\displaystyle \frac{d^4k}{(2\pi )^4}\widehat{T}_{ij}(k)\mathrm{\Gamma }_iS_F(pk)\overline{\mathrm{\Gamma }}_jf^2(pk)},`$
where $`\widehat{T}_{ij}(k)`$ denotes the element of the $`\overline{q}q`$ scattering matrix which describes scattering from the state with matrix structure $`\mathrm{\Gamma }_j`$ to the state with structure $`\mathrm{\Gamma }_i`$. This matrix is diagonal apart from the mixing between the pseudoscalar and axial channels which is present in the extended model with vector and axial interactions. It is to be understood as the LO scattering matrix, which is of order $`N_c^1`$ (see Eq. (10)).
Although the SDE as written above contains all of the required terms at LO and NLO, it also includes some unwanted higher order terms. The NLO terms can be obtained by replacing the quark propagator in the final integral by its LO form, and by expanding the propagator in the quark “bubble” integral to NLO:
$$S_F(p)S(p)+S(p)\mathrm{\Sigma }^{(1)}(p)S(p),$$
(26)
where $`\mathrm{\Sigma }^{(1)}`$ denotes the NLO contribution to the quark self-energy. This leads to an equation for $`\mathrm{\Sigma }^{(1)}`$ of the form
$`\mathrm{\Sigma }^{(1)}(p)`$ $`=`$ $`iG_1f^2(p)\mathrm{Tr}{\displaystyle \frac{d^4k}{(2\pi )^4}S(k)\mathrm{\Sigma }^{(1)}(k)S(k)f^2(k)}`$ (28)
$`if^2(p){\displaystyle \underset{i,j}{}}{\displaystyle \frac{d^4k}{(2\pi )^4}\widehat{T}_{ij}(k)\mathrm{\Gamma }_iS(pk)\overline{\mathrm{\Gamma }}_jf^2(pk)}.`$
The first term comes from an insertion of the NLO self-energy into the quark bubble diagram. The second, involving the scattering matrix, includes both Fock terms and dressing of the quark by virtual mesons.
The separable nature of the interaction means that this equation can be solved straightforwardly to get
$$\mathrm{\Sigma }^{(1)}(p)=f^2(p)Cif^2(p)\underset{i,j}{}\frac{d^4k}{(2\pi )^4}\widehat{T}_{ij}(k)\mathrm{\Gamma }_iS(pk)\overline{\mathrm{\Gamma }}_jf^2(pk),$$
(29)
where
$`C`$ $`=`$ $`{\displaystyle \frac{G_1}{1G_1J_{\sigma \sigma }(0)}}{\displaystyle \underset{i,j}{}}\mathrm{Tr}{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{d^4\mathrm{}}{(2\pi )^4}\widehat{T}_{ij}(k\mathrm{})}`$ (31)
$`\times S(k)\mathrm{\Gamma }_iS(\mathrm{})\overline{\mathrm{\Gamma }}_jS(k)f^4(k)f^2(\mathrm{}).`$
The two terms correspond to the tadpole and meson cloud diagrams shown in Fig. 2. The tadpole diagram is responsible for the contribution $`f^2(p)C`$. It includes the exchange of a zero-momentum sigma meson between the quark and a virtual meson. Since $`C`$ is a momentum-independent constant, this contribution has the same form as the LO quark self-energy. The other diagram describes the emission and subsequent reabsorption of a virtual meson. It has a nontrivial dependence on the momentum of the quark and generates a wave function renormalization as well as a scalar term. Note that although we shall often refer to the upper double line in such diagrams as a meson propagator, it is really a $`\overline{q}q`$ scattering amplitude and so the diagrams contain Fock terms as well as virtual meson contributions.
### C Meson propagators
At LO, the mesonic bound states are constructed from the ladder Bethe–Salpeter equation (Eq. (10)). The NLO extension of the BSE can be expressed in terms of corrections to the basic two-quark loopWe refer to a loop integral containing $`n`$ LO quark propagators as an $`n`$-quark loop. of Eq. (11). If the quantities $`J_{ij}`$ are redefined to include these NLO contributions, the scattering matrix can still be written in the form of Eqs. (9) and (10).
An obvious set of NLO terms in a polarization loop integral $`J_{ij}(q)`$ consists of insertions of the NLO quark self-energies of Fig. 2 on either the quark or antiquark line. These contribute
$`i\mathrm{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\mathrm{\Gamma }_jS(p_{})\mathrm{\Sigma }^{(1)}(p_{})S(p_{})\overline{\mathrm{\Gamma }}_iS(p_+)f^2(p_{})f^2(p_+)}`$ (32)
$`+`$ $`i\mathrm{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\mathrm{\Gamma }_jS(p_{})\overline{\mathrm{\Gamma }}_iS(p_+)\mathrm{\Sigma }^{(1)}(p_+)S(p_+)f^2(p_{})f^2(p_+)}.`$ (33)
A second kind of NLO contribution arises from the exchange of a $`t`$-channel virtual meson between the quark and the antiquark. The corresponding diagram is shown in Fig. 3 and makes the following contribution to $`J_{ij}(q)`$:
$`{\displaystyle \underset{l,m}{}}\mathrm{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\widehat{T}_{lm}(pk)\mathrm{\Gamma }_jS(k_{})\mathrm{\Gamma }_lS(p_{})\overline{\mathrm{\Gamma }}_iS(p_+)\overline{\mathrm{\Gamma }}_mS(k_+)}`$ (34)
$`\times f^2(k_{})f^2(k_+)f^2(p_{})f^2(p_+).`$ (35)
Finally, there are NLO contributions that involve intermediate two-meson states, represented by the second diagram in Fig. 3. These allow for the instability of a meson against two-body decays and so can introduce an imaginary part in the corresponding propagator above the threshold energy for the two-particle final states. They are constructed by connecting two LO three-meson vertices with two meson propagators. The resulting imaginary part of a meson mass can be expressed in terms of the square of the LO decay amplitude. These contributions also generate a shift in the real part of the mass which involves a new type of loop integral.
In writing explicit expressions for these contributions, it is convenient to introduce functions $`L_{ij,k}`$ and $`\overline{L}_{i,jk}`$ to describe the LO three-point vertices. The function $`L_{ij,k}`$ describes the process $`kij`$ and is defined as
$`L_{ij,k}(q_1,q_2,q)`$ $`=`$ $`i\mathrm{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\mathrm{\Gamma }_kS(p_{})\overline{\mathrm{\Gamma }}_jS(p\frac{1}{2}(q_1q_2))\overline{\mathrm{\Gamma }}_iS(p_+)}`$ (39)
$`\times f^2(p_{})f^2(p_+)f^2(p\frac{1}{2}(q_1q_2))`$
$`+i\mathrm{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\mathrm{\Gamma }_kS(p_{})\overline{\mathrm{\Gamma }}_iS(p+\frac{1}{2}(q_1q_2))\overline{\mathrm{\Gamma }}_jS(p_+)}`$
$`\times f^2(p_{})f^2(p_+)f^2(p+\frac{1}{2}(q_1q_2)),`$
while $`\overline{L}_{i,jk}`$ describes the process $`jki`$,
$`\overline{L}_{i,jk}(q,q_1,q_2)`$ $`=`$ $`i\mathrm{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\overline{\mathrm{\Gamma }}_iS(p_+)\mathrm{\Gamma }_jS(p\frac{1}{2}(q_1q_2))\mathrm{\Gamma }_kS(p_{})}`$ (43)
$`\times f^2(p_{})f^2(p_+)f^2(p\frac{1}{2}(q_1q_2))`$
$`+i\mathrm{Tr}{\displaystyle \frac{d^4p}{(2\pi )^4}\overline{\mathrm{\Gamma }}_iS(p_+)\mathrm{\Gamma }_kS(p+\frac{1}{2}(q_1q_2))\mathrm{\Gamma }_jS(p_{})}`$
$`\times f^2(p_{})f^2(p_+)f^2(p+\frac{1}{2}(q_1q_2)).`$
The contribution to $`J_{ij}(q)`$ from the two-meson diagrams can then be written in the form
$$\frac{i}{2}\underset{k,l,m,n}{}\frac{d^4p}{(2\pi )^4}\overline{L}_{i,km}(q,p_+,p_{})\widehat{T}_{kl}(p_+)\widehat{T}_{mn}(p_{})L_{ln,j}(p_+,p_{},q).$$
(44)
Once the full NLO matrix $`J^{(1)}`$ has been constructed, we can use it to find the corrections to meson masses and meson-quark couplings. As an example, consider the pion in the version of the model without vector and axial couplings. The pion pole is located at
$$1G_1\left[J_{\pi \pi }^{(0)}(q^2)+J_{\pi \pi }^{(1)}(q^2)\right]=0,$$
(45)
where superscripts $`(0)`$ and $`(1)`$ denote LO and NLO terms. The coupling to quarks (Eq. (14)) is given by
$$\left(g_{\pi qq}^{(0)}+g_{\pi qq}^{(1)}\right)^2=\frac{d(J_{\pi \pi }^{(0)}+J_{\pi \pi }^{(1)})}{dq^2}|_{q^2=m_\pi ^2}.$$
(46)
From these we get the squared pion mass to NLO,
$$m_\pi ^2=m_\pi ^{(0)\mathrm{\hspace{0.17em}2}}\left[1+g_{\pi qq}^{(0)\mathrm{\hspace{0.17em}2}}\frac{dJ_{\pi \pi }^{(1)}}{dq^2}\right]_{q^2=m_\pi ^{(0)\mathrm{\hspace{0.17em}2}}}g_{\pi qq}^{(0)\mathrm{\hspace{0.17em}2}}J_{\pi \pi }^{(1)}(m_\pi ^{(0)\mathrm{\hspace{0.17em}2}}),$$
(47)
and the pion-quark coupling,
$$g_{\pi qq}^{(1)}=\frac{g_{\pi qq}^{(0)\mathrm{\hspace{0.17em}3}}}{2}\left[2m_\pi ^{(0)}m_\pi ^{(1)}\frac{d^2J_{\pi \pi }^{(0)}}{d(q^2)^2}+\frac{dJ_{\pi \pi }^{(1)}}{dq^2}\right]_{q^2=m_\pi ^{(0)\mathrm{\hspace{0.17em}2}}}.$$
(48)
### D Meson coupling to a current
In this section, we consider a NLO determination of the coupling between a meson and an external current, taking the pion decay constant as an example. The situation is more complicated in a nonlocal than a local model because there are two kinds of contribution to be considered, arising from the local and nonlocal parts of the symmetry current (see Sec. II).
The LO diagrams are shown in Fig. 1. At NLO several of the extra contributions can be identified straightforwardly, either by inserting the NLO self-energy on any one of the quark lines in these diagrams or by using the NLO contribution to the pion vertex function. Other contributions are analogous to those in the BSE at NLO. These involve either a $`t`$-channel meson exchange in the two-quark loop or intermediate two-meson states. They are shown in Fig. 4.
The remaining NLO contributions arise solely from the nonlocal piece of the current. The first of these are Fock terms arising from the exchange of two of the quark lines at a nonlocal current vertex, as shown in Fig. 5. Since they contain only a single colour trace, as indicated by separating the quark lines associated with each of the $`\overline{\psi }\mathrm{\Gamma }\psi `$ factors, this diagram is suppressed by one power of $`N_c`$ compared with the corresponding LO terms. Although, as noted above, there is an ambiguity in defining the transverse components of such Fock terms, this does not affect the calculation of the pion decay constant.
Finally there are two NLO contributions which can be thought of as arising from nonlocal current couplings to virtual mesons. They are shown in Fig. 6. The first (a) is somewhat similar to diagram 4(d), except that the coupling of the nonlocal current to the virtual mesons involves both the quark and the antiquark instead of a separate quark bubble. This is analogous to the two-body diagrams that contribute at LO to several of the electromagnetic amplitudes described in Ref.. The other diagram (b) can be thought of as coupling the current at a meson-quark vertex.
Although the NLO contributions to the pion decay constant involve rather a large number of diagrams, in practice the situation can be significantly simplified as a result of cancellations amongst them. A complete description of these cancellations would be somewhat lengthy and so we offer a brief summary in Appendix A. We have also checked the consistency of our treatment by verifying that the GMOR relation holds to NLO (see Appendix B).
## IV Numerical results
### A Model parameters
The multiple integrals involved in the NLO diagrams can be rather time-consuming to evaluate numerically to good accuracy. We have therefore chosen to fit the model parameters at LO and then to calculate the NLO changes to observables. In fact, for the parameter sets of most interest, we find that the corrections are small for the observables used to fix the parameters. Thus, our results should not be very different from those of a full NLO fit.
We fit the parameters at LO as in Ref.. In the simpler version of the model with no vector and axial couplings, we fit two parameters to the pion mass and decay constant. This leaves one free parameter which we characterize in terms of $`m_0(0)`$, the quark self-energy at zero momentum in the chiral limit. Details of the sets used are given<sup>§</sup><sup>§</sup>§Note that the parameters given differ slightly from those quoted in Ref. where a very similar fit was made at the same values of $`m_0(0)`$. This is simply because the calculations of Ref. were performed within the chiral expansion of the model. in Table I. Values of some quantities calculated at LO with these parameter sets are also given in the table. They are qualitatively similar to those obtained in the extended version of the model. For the parameter sets with $`m_0(0)=300`$ MeV and larger, the quark propagator has no poles on the real axis and so single quarks do not appear as physical states.
Although most of our work at NLO has used this simpler model, we also show a few results for the extended model including vector and axial couplings. For these we use a parameter set with $`m_0(0)=300`$ MeV, which is intermediate between the two cases studied at LO in Ref..
### B Quark condensate
The NLO term in the quark condensate can be evaluated as the sum of two contributions. These correspond to the two self-energy diagrams of Fig. 2 where the external quark line is formed into a closed loop with a local scalar insertion. By comparing the values for the condensate at NLO in Table II with the LO values in Table I, we see that the NLO contributions are very small. This may seem surprising since virtual pions would be expected to make a large (positive) contribution to the condensate. However it should be recalled that the meson “propagators” in the NLO diagrams also contain the Fock terms of the two-body interaction. These make a negative contribution to the condensate which tends to cancel the contributions of virtual mesons. This agrees with the findings of Ref. in a similar model. In the local version of the NJL model, the condensate is also little altered at NLO if proper time regularization of the quark loops is used, but in the $`O(4)`$ scheme there can be appreciable changes at NLO. Including only the pole pieces of the meson propagators, significant shifts in the condensate are also found with $`O(3)`$ regularization.
### C Quark propagator
The quark self-energy at NLO is given by Eqs. (29) and (31). Its evaluation requires integrating over the scattering matrix $`\widehat{T}`$. Evaluating the $`\overline{q}q`$ loop integrals $`J`$ for each point in the mesonic loop integral is very time consuming. We have therefore found it convenient to parameterize the $`J`$’s in terms of a set of smoothly-varying functions. Working in Euclidean space, and with the momentum routings of Eqs. (29) and (31), the $`J`$’s always have a spacelike momentum argument and so are smoothly-varying functions. We have found that they are well represented by expansions in Chebyshev polynomials of the variable $`x=\mathrm{exp}(q^2/\mathrm{\Lambda }^2)`$.
We write the inverse quark propagator as
$$S_F^1(p)=(1+a(p))\text{/}pb(p).$$
(49)
The functions $`a(p)`$ and $`b(p)`$ are shown in Figs. 7 and 8, for both spacelike and timelike momenta. They are plotted only up to the energy given by the real part of the complex pole in the LO quark propagator. Above that energy there are pseudo-threshold effects in the model associated with the continuation from Euclidean to Minkowski space.
At LO $`a(p)=0`$ and there is no wave function renormalization. The NLO contributions to the coefficient $`\text{/}p`$ range up to about 0.25, which is consistent with the expected magnitude of $`1/N_c`$ corrections. An interesting aspect of the results for $`a(p)`$ is the appearance of a sudden dip just before the pseudo-threshold energy. It would be interesting to examine the behaviour of the function above this energy, although that would require a detailed analysis of the quark pole structure at NLO, which is beyond the scope of the present work. Also plotted are the individual contributions to $`a(p)`$ which arise from dressing the quark line with virtual pions and with virtual sigma mesons. The pion cloud is obviously the dominant contributor. Since its propagator has a pole at small timelike momentum, one expects the $`T`$ matrix in the pseudoscalar channel to be large at modest values of spacelike momenta (the region which dominates the NLO integrals). Note also that the pion contributions contain an extra factor of three due to isospin.
The coefficient of the unit matrix in the quark self-energy at NLO receives contributions from both the tadpole and the meson-cloud diagrams of Fig. 2. The results in Fig. 8 show that the tadpole contribution is in fact rather small, $`Cm(0)`$. The meson-cloud contributions are rather more significant, increasing $`b(0)`$ by about $`25\%`$ (again consistent with expectations for a $`1/N_c`$ effect). The net NLO shift in the quark “mass” at zero momentum $`b(0)/(1+a(0))`$ is therefore fairly modest, $`15\%`$. Quark dressing by meson clouds in the local NJL model has been investigated in Ref.. It was found that the pion cloud tends to increase $`b(p)`$ but that this is partially cancelled by the sigma cloud. The nonlocal model studied here supports that conclusion and is able to place it on a firmer footing since there are no ambiguities associated with the meson loop regularizationNote that the tadpole contributions were not included in Ref.. Moreover the model meson propagators were approximated by the canonical forms for point-like bosons..
Although the tadpole diagrams in the NLO quark self-energy are not numerically significant, $`C`$ controls the NLO changes to the pion mass and decay constant. This can be seen in the chiral expansions of these quantities described in Appendix B. One might therefore wonder whether there is some physical reason why its chiral limit $`C_0`$ turns out to be small. For the parameter set with $`m_0(0)=200`$ MeV Table II shows that the pion tadpole makes a contribution of about $`45\%`$ of the LO chiral quark mass, but that this is cancelled to a large extent by the sigma tadpole. Such cancellations do not persist at larger $`m_0(0)`$, where the pion tadpole changes sign. However for these parameter sets both tadpole contributions are small. In these cases it is a cancellation between “true” virtual meson effects and Fock terms (similar to what happens in the quark condensate) which is responsible for the small net effect.
In the alternative scheme of Ref. , the $`1/N_c`$ contribution of the tadpole diagram in Fig. 2 is resummed to all orders when the SDE is solved self-consistently. In contrast, the momentum-dependent contribution of the meson-cloud diagram is treated perturbatively. This partial resummation is motivated by an expectation that contributions from the tadpole diagrams are larger than the other NLO effects. However we find that such an expectation is not fulfilled in the nonlocal model; in fact we find that the contribution of the tadpole diagram is small compared to that of the meson-cloud diagram. We thus see no compelling reason to resum this contribution, and so we can avoid possible problems with the self-consistent scheme which have been noted in the context of the local NJL model.
In the extended version of the model there are additional contributions to the NLO quark self-energy from tadpole and meson-cloud diagrams with other mesons. Although a calculation of the properties of these mesons at NLO is beyond the scope of the present work, it is straightforward to evaluate their contributions to the quark self-energy. As in the simpler model, we find that there is only a modest NLO change to the condensate and that the constant $`C`$ is small. Pseudoscalar-axial mixing makes a significant difference to the contribution from the pion cloud to the function $`b(p)`$, but this is largely cancelled by the vector mesons to leave a function very similar to that found in the simpler version of the model.
The function $`a(p)`$ in the extended model is shown in Fig. 9. Mixing in the pion channel and the clouds of the spin-1 states are again important. In $`a(p)`$ these effects reinforce each other and so the function is significantly larger than in the simpler model. An interesting difference between Figs. 7 and 9 is the absence in the latter case of a dip just before the LO pseudo-threshold energy. This is eliminated in the overall result due to the contributions from the longitudinal components of the vector and axial propagators. In particular, the steep drop in the pion contribution which occurs in the simpler model is removed by mixing with the axial channel. We note that the local NJL model has been used in Ref. to describe the $`\rho `$ meson at NLO, assuming that the effects of $`\pi a_1`$ mixing and intermediate spin-1 states could be neglected. The results presented here raise some doubts about those assumptions. Of course definitive statements about the quark propagator in the present model cannot be made without a full NLO analysis of meson properties in order to refit the model parameters. Since $`1/N_c`$ effects could alter meson masses significantly, the parameters used here may not constitute a reasonable choice for the extended model at NLO.
### D Pion properties
The denominator of the $`\overline{q}q`$ scattering amplitude in the pion channel, Eq. (45), is shown in Fig. 10. For low energies, less than about 250 MeV, the differences between the LO and NLO curves are very small. This implies that NLO contributions to both the pion mass and the pion-quark coupling are extremely small.
Although NLO contributions to the pion mass must cancel in the chiral limit, as described in Appendix B, such cancellations do not influence higher order terms in the chiral expansion. The smallness of the NLO terms in the pion propagator is another consequence of the cancellation between virtual mesons and Fock terms, as we saw above in the case of the quark propagator. In particular, all of the contributions to the pion mass and decay constant that are proportional to the NLO piece of the quark condensate or to $`C`$ are small. (More details of these quantities are given in the appendices.) This supports the usual LO treatment of the pion in this type of model and it also justifies our use of model parameters fitted to the LO pion properties in these NLO computations. Actual determinations of the NLO shifts in $`m_\pi `$ and $`f_\pi `$ have not been made in this work since they are sufficiently small (not more than a couple of MeV) that our numerical integration procedures would have to be refined in order to yield accurate values.
At higher energies, beyond those shown in Fig. 10, NLO contributions in the pion channel do start to become significant. When the $`\pi \sigma `$ channel opens, the denominator develops an imaginary part. Shortly before the pseudo-threshold energy (twice the real part of the complex energy of the pole in the quark propagator) is reached, the NLO terms become sufficiently important to change the qualitative behaviour of the pion Bethe-Salpeter amplitude. The slope of the denominator changes sign and there is even a second zero of Eq. (45). This is a potentially worrying feature since it corresponds to an unphysical pole in the amplitude, whose residue has the wrong sign. A similar pole has been found in the local NJL model. In both models the undesirable high-energy behaviour is caused by the insertions of NLO quark self-energies on the quark lines of the basic polarization loop integral. However it is perhaps worth noting that this behaviour is reminiscent of that exhibited by $`a(p)`$ at high energies in Fig. 7. As mentioned in Sec. IV C, this behaviour is not present in the model with vector and axial mesons and it is thus possible that the unphysical pion pole may be removed by extending the model to include interactions in these channels.
### E Sigma properties
The interpretation of the scalar, isoscalar sigma meson in dynamical quark models is complicated by the fact that it does not correspond to an experimentally well-determined resonance. The problem, not just for models of such a resonance but also for phenomenological determinations of its properties, is its very strong coupling to the two-pion channel. Indeed, while some phenomenological analyses favour a scalar, isoscalar resonance at around 1 GeV, more recently others have found resonances at 600 MeV or lower (for reviews and more complete lists of references see Refs.). The strong coupling to pions means that any resonance must form a very broad structure. Consequently the corresponding pole in the scattering amplitude must lie far from the real axis, which makes its exact position hard to determine in a model-independent way.
This strong coupling to two pions also means that a bare $`\overline{q}q`$ state in the scalar, isoscalar channel should not be directly compared with any phenomenological resonance. The effects of the two-pion channel need to be included. In the present model, the leading effects of this type arise from the second diagram of Fig. 3.
At LO in $`1/N_c`$, the nonlocal NJL model leads to a $`\overline{q}q`$ sigma meson with a mass of less than 500 MeV. This is similar to the local model, where the sigma lies on the $`\overline{q}q`$ threshold. Some of the properties of the state are listed in Table III. It is strongly coupled to two pions, although its low mass means that its decay width is less than about 160 MeV. Although wide, this is still narrower than phenomenological fits which tend to give widths of at least 300 MeV, even for low sigma masses.
The real part of the denominator of the scattering amplitude at NLO is shown in Fig. 11. Above the two-pion threshold, there is an imaginary part generated by the contribution to $`J_{\sigma \sigma }`$ from the second diagram of Fig. 3. This can be estimated from the LO decay amplitude using the Cutkosky cutting rules and it is shown in Fig. 12 along with the results from a direct numerical evaluation. Some care must be taken in the evaluation of this diagram owing to singularities of the scattering matrices in the integrand above the decay threshold. We have regulated Eq. (44) by replacing $`\widehat{T}_{\pi \pi }`$ with $`G_1(1G_1J_{\pi \pi }^{(0)}iϵ)^1`$ and linearly extrapolated our results to $`ϵ=0`$, based on several computations for $`ϵ10^3`$.
The values listed in Table III for the mass of the sigma at NLO are for the mass defined as the zero of the real part of the denominator. Although this is a commonly used choice, it might be better to use the real part of the complex energy of the pole in the scattering matrix, as discussed by Pennington. The choice used here is purely one of convenience. In any case the width of the state we find is not so large that we would expect substantial differences between the two definitions.
We find that the sigma mass is increased from its LO value by about 30%. Below the two-pion threshold the largest NLO contribution to $`J_{\sigma \sigma }`$ arises from virtual two-pion intermediate states. This is partially cancelled by the NLO quark self-energy insertions to leave a modest net increase in $`J_{\sigma \sigma }`$ for most parameter sets. On their own these effects of virtual mesons would tend to reduce the sigma mass. However, above the threshold, real two-pion states make a substantial contribution with the opposite sign, leading to the increase to the sigma mass. This is consistent with the qualitative expectation that such states should be important for modelling the sigma.
NLO corrections to the sigma mass have been calculated by Pallante using a derivative expansion of the bosonized NJL model. In that framework the corrections were found to be large and negative, prompting that author to speculate that the $`1/N_c`$ expansion might break down for the mass of this state. In the local NJL model, the large negative corrections to $`1G_1J_{\sigma \sigma }`$ persist in calculations to all orders in momentum if one uses Pauli-Villars regularization of the quark loops and a sharp cutoff for the meson loop momenta. The shifts in this case are sufficiently large to move the model sigma pole to negative $`q^2`$, producing an unstable vacuum. In both of these calculations, however, the results in the scalar sector were very sensitive to the additional cut-off parameter needed to regularize meson loops in the local model. In contrast, the magnitude of the NLO mass shift which we obtain in the nonlocal model is consistent with expectations for a $`1/N_c`$ correction. Thus it would appear that the details of the treatment of high-momentum states are important in NJL-type models.
The imaginary part of the polarization loop $`J_{\sigma \sigma }`$ is shown in Fig. 12 and compared with the estimate obtained by applying the Cutkosky cutting rules to graphs with two-pion intermediate states. The difference between these curves is caused by the analytic structure of the quark propagator, which contains additional, unphysical poles (Sec. II). It is reassuring to note that these extra poles are located sufficiently far from the low-energy region of interest that they do not have a large effect on the imaginary part. For other parameter sets, with large $`m_0(0)`$, the extra poles lie closer to the physical region and may have some influence on properties of the sigma meson at NLO. Since the model predictions may be unreliable in these cases, we have preferred to quote only upper bounds on the sigma mass for some parameter sets in Table III, limited by the position of the first pole in the LO quark propagator.
For energies up to $``$600 MeV, the three-quark integral in the $`\sigma \pi \pi `$ vertex function is only weakly dependent on the energy. Therefore the variation of $`\mathrm{Im}(J_{\sigma \sigma })`$ at such energies is primarily a consequence the available two-pion phase space. At higher energies the coupling to two pions becomes much weakerThis supports a suggestion made in Ref. that the weak coupling of the scalar channel to two pions above the sigma mass means that the broad width found for $`a_1\sigma \pi `$ in the model need not be inconsistent with the experimental observation of a small width for $`a_1\pi (\pi \pi )_s`$.. This is similar to the behaviour of the coupling of the sigma to two pions observed in the quark model studied in Ref. (where the scalar mass was taken as a free parameter and chosen with the intention of interpreting the model scalar resonance as a heavy, narrow state).
We note that a two-flavour NJL model of pions and sigma mesons is insufficient for a fully realistic description of the spectrum of scalar mesons. One must also take account of strangeness, radial excitations, mixing with $`\overline{q}q`$ molecules or glueballs, and decays to four pion states. A recent attempt to include such effects can be found in Ref..
## V Conclusions
We have investigated the effects of meson fluctuations on quark and meson properties in a nonlocal NJL model by treating it to NLO in $`1/N_c`$. The two-body interaction between quarks has a separable form, similar to that suggested by instanton liquid studies. At LO the model has been shown to provide successful descriptions of mesons and baryons. The interaction form factors regulate both quark and meson loops in a natural way. As a result, no new parameters emerge at NLO and the contributions are unambiguous. This makes the model a particularly suitable one in which to study NLO effects.
The contributions to the Schwinger–Dyson and Bethe–Salpeter equations at NLO contain the same structures as those in the local model, apart from the presence of interaction form factors. The NLO pieces include the dressing of quarks and mesons with virtual mesons and the coupling of mesons to two-meson states. They also include the exchange or Fock terms for the basic interaction between quarks.
With a nonlocal interaction between the quarks, there are two-body terms in the symmetry currents and these must be taken into account in the NLO contributions to the couplings of mesons to external currents. We have obtained the NLO pieces of the currents by first Fierz transforming our interaction and then applying the technique of Refs.. As a consistency check on our treatment we have verified that the GMOR relation holds at NLO.
In the simple version of the model with only scalar and pseudoscalar interactions we find that the NLO contributions to a number of quantities are small for the parameter sets of most interest (the ones with LO quark self-energies of about 300–500 MeV). These quantities include the quark condensate and the tadpole piece of the quark self-energy. The small size of the NLO pieces is a consequence of cancellations between the effects of virtual mesons and Fock terms, as also noted by Ripka. Similar cancellations mean that the NLO contributions to the pion mass and decay constant are also small. At NLO these quantities depend not just on the NLO pieces of the quark condensate and pion-quark coupling, but also on the tadpole piece of the quark self-energy. The fact that this NLO tadpole energy is small is crucial to the smallness of the corrections to pion properties. In other cases, such as the “meson cloud” pieces of the quark-self energy, NLO contributions are larger. However they typically contribute at the $`25\%`$ level and so give reason to hope that the $`1/N_c`$ expansion is a useful organizing principle for this model.
In this context, it is worth making a remark about the recent claim of Kleinert and Van den Bossche that meson fluctuation effects in the NJL model are large enough to restore manifest chiral symmetry, at least for small values of $`N_c`$. Although our approach is based on a loop expansion, one might have expected to see large decreases in the magnitudes of quantities like the quark condensate, the quark self-energy and the pion decay constant if mesonic fluctuations were driving the system towards symmetry restoration. In fact pionic fluctuations do make significant contributions in this direction, but they are largely cancelled by the Fock terms. In the language of Ref. , the Fock terms contribute at NLO in $`1/N_c`$ to an increased “stiffness” against mesonic fluctuations. This raises a question about the estimate of the critical stiffness in that work, which relies on an expression for the stiffness obtained from the LO pion propagator. This is in addition to the questions raised by other authors about the universality of the critical stiffness in $`3+1`$ dimensions and the need to choose an additional cut-off to regulate meson loops in the local NJL model.
The sigma meson is of particular interest since it can be excited by forces which act to restore manifest chiral symmetry, and its mass can be thought of as describing the forces against symmetry restoration. A light sigma seems to be a common feature in all NJL-type models at LO. We find that such a state is still present when NLO effects are included, although its mass is increased by about 30% (as one might expect from $`1/N_c`$ arguments). For our preferred parameter sets we obtain a sigma mass in range 600–650 MeV. This suggests that the light sigma meson is a real feature of this type of model and not just an artefact of LO treatments.
However, the phenomenological identification of a corresponding low-lying scalar, isoscalar resonance remains controversial. As discussed above, the problem is the strong coupling of such a state to the two-pion channel, leading to a very broad resonance. In our model we find a width of 100–160 MeV from the LO coupling to two pions. While large, this width is much smaller than most phenomenological determinations. We have not yet calculated NLO contributions to the $`\sigma \pi \pi `$ coupling, for which terms involving the vector mesons may prove to be important. In the extended model at NLO there will be $`t`$-channel exchange of a $`\rho `$ meson between the pions. The attractive force between the pions generated by this mechanism is known to be important for the sigma resonance and could increase its decay width.
Finally, we have made a preliminary study of NLO effects in the extended model with vector and axial couplings. We find that NLO contributions, in particular from the longitudinal vector and axial channels, can improve the high-energy behaviour of the quark propagator. In general NLO effects seem to be significantly larger than in the simple version of the model, and the convergence of the $`1/N_c`$ expansion may be less good as a result.
## Acknowledgements
This work was supported by the EPSRC and PPARC.
## A Pion decay constant at NLO
Although the NLO contributions to the pion decay constant arise from a rather large set of diagrams, it is possible to deduce some useful cancellations amongst them. A complete description of these cancellations would be somewhat lengthy and so we offer only a few brief comments here. Full details can be found in Ref..
Obviously the diagrams that couple the local axial current to the pion all contain an insertion of $`\gamma _\mu \gamma _5`$. Contracting this with $`q^\mu `$ to isolate the longitudinal component one can then substitute for $`\text{/}q\gamma _5`$ from the following identity:
$$\text{/}q\gamma _5=S^1(p_+)\gamma _5+\gamma _5S^1(p_{})+(m(p_+)+m(p_{}))\gamma _5.$$
(A1)
The contributions from the first two terms here can be simplified by cancelling the LO inverse propagators with quark propagators that appear in the loop integrals.
Moreover, for each local diagram there is a similar nonlocal one which contains an additional one-quark loop. (For example, compare Fig. 4a with 4c and Fig. 4b with 4d.) These can be combined with the local contributions, as was done at LO. Each of the NLO nonlocal diagrams can be written as a sum of products of two-quark and one-quark loops. In one of these terms, the one-quark loop can be simplified by using the ladder SDE to yield a factor proportional to $`m(0)m_c`$. There is then a cancellation between this nonlocal term and the part of the corresponding local diagram arising from the final term of Eq. (A1) which leaves only $`2m_c`$ surviving from the factor $`m(p_+)+m(p_{})`$ in Eq. (A1).
After exploiting Eq. (A1) as outlined above, several other useful cancellations can be identified. As an example of the procedure, consider the diagram shown in Fig. 4b. When Eq. (A1) is substituted for $`\text{/}q\gamma _5`$ in the triangular loop integral, the first two terms reduce to a sum of two-quark loops. Up to a factor of $`GJ`$ the resulting contributions are then similar in form to those of the diagram shown in Fig. 6a.
Hence, in the sum of the contributions from Fig. 4b and 6a, one can cancel one of the intermediate meson propagators to leave a product of a two-quark loop, a three-quark loop and one meson propagator. It turns out that this can then be profitably combined with the contribution of Fig. 6b, which has the same structure.
We give here the result of a full analysis of the NLO contributions to $`f_\pi `$ in the version of the model with only scalar and pseudoscalar interactions. Note that, rather than defining further loop integrals whose structure is very similar to those given above in Sec. III C, we simply indicate how the products of form factors need to be changed in those integrals. The final expression for $`f_\pi `$ at NLO is
$`f_\pi ^{(1)}(q^2)`$ $`=`$ $`{\displaystyle \frac{ig_{\pi qq}^{(0)}G_1}{2q^2}}J_{\pi \pi }^{(1)}(q^2){\displaystyle \frac{d^4k}{(2\pi )^4}\mathrm{Tr}\left[S(k)\right]f(k)\left(f(k+q)+f(kq)\right)}`$ (A12)
$`+{\displaystyle \frac{g_{\pi qq}^{(0)}m_c}{q^2}}\left(J_{\pi \pi }^{(1)}\text{from Eq. (}\text{33}\text{) with}f^2(p_+)f^2(p_{})f(p_+)f(p_{})\right)`$
$`+{\displaystyle \frac{g_{\pi qq}^{(0)}m_c}{q^2}}\left(J_{\pi \pi }^{(1)}\text{from Eq. (}\text{35}\text{) with}f^2(p_+)f^2(p_{})f(p_+)f(p_{})\right)`$
$`+{\displaystyle \frac{g_{\pi qq}^{(0)}m_c}{q^2}}(J_{\pi \pi }^{(1)}\text{from Eq. (}\text{44}\text{) with}f^2(p_+)f^2(p_{})f(p_+)f(p_{})`$
$`\text{in}\overline{L}_{\pi ,km})`$
$`i{\displaystyle \frac{Cg_{\pi qq}^{(0)}}{2q^2}}\left(1G_1J_{\pi \pi }^{(0)}(q^2)\right){\displaystyle \frac{d^4p}{(2\pi )^4}\mathrm{Tr}\left[S(p)S(p)\right]f^3(p)}`$
$`\times \left(f(p+q)+f(pq)\right)`$
$`{\displaystyle \frac{g_{\pi qq}^{(0)}}{2q^2}}\left(1G_1J_{\pi \pi }^{(0)}(q^2)\right){\displaystyle \underset{i,j}{}}{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\widehat{T}_{ij}(pk)}`$
$`\times \mathrm{Tr}\left[S(p)\mathrm{\Gamma }_iS(k)\overline{\mathrm{\Gamma }}_jS(p)\right]`$
$`\times f^3(p)\left(f(p+q)+f(pq)\right)f^2(k)`$
$`+{\displaystyle \frac{g_{\pi qq}^{(1)}}{g_{\pi qq}^{(0)}}}f_\pi ^{(0)}(q^2).`$
## B Chiral relations
A large number of diagrams are required in a NLO treatment of the nonlocal NJL model, as described in Sec. III. It is important to have some check that a consistent set has been identified and evaluated. To this end, we demonstrate that low-energy chiral constraints are satisfied in the model at NLO. (Again we consider only the simpler version of the model, without vector and axial interactions.)
First consider the pion mass. The chiral expansion of the pion denominator, Eq. (45), at LO is
$$1G_1J_{\pi \pi }^{(0)}(q^2)=G_1m_c\frac{\overline{\psi }\psi _0^{(0)}}{m_0(0)^2}G_1\frac{q^2}{g_{\pi qq0}^{(0)\mathrm{\hspace{0.17em}2}}}+𝒪(q^4,m_c^2),$$
(B1)
where the subscript $`0`$ is used to denote a quantity evaluated in the chiral limit. (We apologize for the fact that the notation inevitably becomes rather cumbersome at this point.) If one substitutes the LO chiral expansion into the NLO on-shell condition, it is immediately clear that the chiral expansion of $`J_{\pi \pi }^{(1)}(q^2)`$ must start at $`𝒪(q^2,m_c)`$ in order for the pion to remain a Goldstone boson in the chiral limit.
A proof of this statement in the local NJL model was presented in Ref.. Although one has to keep track of interaction form factors, the proof for the nonlocal model proceeds along very similar lines. In particular, when one takes the chiral limit of the contributions to $`J_{\pi \pi }^{(1)}`$ from Eqs. (33) and (44), each of the Dirac traces yields a factor that cancels with the denominator of one of the LO quark propagators. The three-quark loop in the contribution of the tadpole term to Eq. (33) can be reduced to a two-quark integral which can then be used to cancel the denominator of the $`\sigma (0)`$ propagator from Eq. (31). In the case of the contribution to $`J_{\pi \pi }^{(1)}`$ from Eq. (35), the Dirac trace takes an even more convenient form in the chiral limit. For both the sigma and pion exchanges, it factorizes into pieces which cancel the denominators from two LO quark propagators.
An important property noted in the proof of Ref. concerns the second diagram of Fig. 3 which involves a virtual pion and a virtual sigma meson. In the chiral limit the three-quark loop integrals are proportional to the difference between LO polarization loops: $`J_{\sigma \sigma }^{(0)}J_{\pi \pi }^{(0)}`$. This allows one to replace the product of the virtual meson propagators with their difference. It is then possible to combine the contribution from this diagram with other contributions, which contain only a single meson propagator. The same property holds in the nonlocal model<sup>\**</sup><sup>\**</sup>\**It is in fact straightforward to generalize the proof that $`J_{\pi \pi 0}^{(1)}(0)=0`$ to the extended version of the model, including other mesons. In so doing the products of meson propagators occurring in diagrams of the type shown on the right–hand side of Fig. 3 can be dealt with using the relations $`J_{\sigma \sigma }^{(0)}(q^2)J_{\pi \pi }^{(0)}(q^2)=J_{VV}^{(0)T}(q^2)J_{AA}^{(0)T}(q^2)=J_{VV}^{(0)L}(q^2)J_{AA}^{(0)L}(q^2)`$, where $`T`$ and $`L`$ denote transverse and longitudinal components respectively. Note, however, that in the extended model, one must take $`\pi a_1`$ mixing into account and so to complete the proof that the pion is a Goldstone boson at NLO, it is also necessary to establish that $`J_{A\pi 0}^{(1)}(0)=0`$..
Many of the simplifications that lead to the vanishing of $`J_{\pi \pi }^{(1)}`$ in the chiral limit can also be exploited in calculating its $`𝒪(m_c)`$ and $`𝒪(q^2)`$ terms. One finds that
$$J_{\pi \pi }^{(1)}=m_c\frac{\overline{\psi }\psi _0^{(1)}}{m_0^2(0)}m_c\frac{\overline{\psi }\psi _0^{(0)}}{m_0^2(0)}\frac{2C_0}{m_0(0)}q^2\frac{2g_{\pi qq0}^{(1)}}{g_{\pi qq0}^{(0)\mathrm{\hspace{0.17em}3}}}+𝒪(q^4,m_c^2)$$
(B2)
where the NLO contribution to the condensate is given by:
$$\overline{\psi }\psi ^{(1)}=i\mathrm{Tr}\frac{d^4p}{(2\pi )^4}S(p)\mathrm{\Sigma }^{(1)}(p)S(p).$$
(B3)
Substituting the above expansion into the on-shell condition (Eq. (45)) yields the pion mass to this order,
$$m_\pi ^2=\left(g_{\pi qq0}^{(0)}+g_{\pi qq0}^{(1)}\right)^2\frac{m_c[\overline{\psi }\psi _0^{(0)}+\overline{\psi }\psi _0^{(1)}]}{(m_0(0)+C_0)^2}+𝒪(m_c^2)+𝒪(N_c^2).$$
(B4)
We see that the chiral expansion of the pion Bethe–Salpeter amplitude at NLO retains a similar structure to that at LO. The changes include the expected ones from the NLO contributions to the quark condensate and the pion-quark coupling. However, the form of the dynamical mass scale that appears in Eq. (B4) is less obvious, the shift in this being given entirely by the contribution of the tadpole diagram, Eq. (31). This is despite the fact that the meson cloud of a dressed quark contributes a term of $`𝒪(m_c)`$ to the chiral expansion of $`J_{\pi \pi }^{(1)}`$ and also makes a finite contribution to the quark self-energy at zero momentum. The fact that the NLO shift in the pion mass is controlled by the coefficient of the tadpole diagram has also been noted in the local NJL model, in which the overall mass scale can be recognized as the stationary point of the one-meson-loop effective action.
From Eq. (B4) it can be seen that the Gell-Mann–Oakes–Renner relation will be satisfied if a version of the Goldberger–Treiman relation holds in the chiral limit at NLO,
$`f_{\pi 0}^{(0)}+f_{\pi 0}^{(1)}`$ $`=`$ $`{\displaystyle \frac{m_0(0)+C_0}{g_{\pi qq0}^{(0)}+g_{\pi qq0}^{(1)}}}`$ (B5)
$`=`$ $`{\displaystyle \frac{m_0(0)}{g_{\pi qq0}^{(0)}}}{\displaystyle \frac{g_{\pi qq0}^{(1)}}{g_{\pi qq0}^{(0)}}}{\displaystyle \frac{m_0(0)}{g_{\pi qq0}^{(0)}}}+{\displaystyle \frac{C_0}{g_{\pi qq0}^{(0)}}}+𝒪(N_c^2).`$ (B6)
The term of order $`N_c^0`$ on the right-hand side of this condition was shown in Ref. to be equal to the LO part of $`f_\pi `$.
In order to demonstrate that the NLO part of the decay constant reduces to the next two terms of Eq. (B6) in the chiral limit, first note that all of the contributions in Eq. (A12) are of zeroth order in the chiral expansion. Hence in all of the integrals the chiral limit may be taken directly, without making an expansion of the integrand. In the first term of Eq. (A12) the ladder SDE can be used to perform the integral, giving
$$2m_0(0)\frac{g_{\pi qq0}^{(1)}}{g_{\pi qq0}^{(0)\mathrm{\hspace{0.17em}2}}}+m_c\frac{g_{\pi qq0}^{(0)}}{q^2}\frac{\overline{\psi }\psi _0^{(1)}}{m_0(0)}m_c\frac{g_{\pi qq0}^{(0)}}{q^2}\frac{\overline{\psi }\psi _0^{(0)}}{m_0(0)}\frac{2C_0}{m_0(0)}$$
(B7)
where Eq. (B2) has been used for the chiral expansion of $`J_{\pi \pi }^{(1)}`$.
The next three terms in Eq. (A12) are explicitly proportional to $`m_c`$ and contain integrals which are very similar to those occurring in $`J_{\pi \pi }^{(1)}`$, but with different combinations of form factors. Procedures for manipulating such integrals in order to simplify their sum were described in Ref. and also above, where they were exploited in order to show that $`J_{\pi \pi }^{(1)}`$ vanishes in the chiral limit. Very similar simplifications can be made here, despite the different form factor structures that appear. These give rise to a nonvanishing contribution to $`f_{\pi 0}^{(1)}`$,
$$\frac{g_{\pi qq0}^{(0)}m_c}{q^2m_0(0)}\left(\overline{\psi }\psi _0^{(1)}+\frac{C_0}{m_0(0)}\overline{\psi }\psi _0^{(0)}\right).$$
(B8)
The fifth term in Eq. (A12) is proportional to the tadpole contribution $`C`$. Using Eq. (B1) to expand the factor of $`1G_1J_{\pi \pi }^{(0)}(q^2)`$, the chiral limit of this term is
$$\frac{g_{\pi qq0}^{(0)}}{q^2}C_0G_1J_{\sigma \sigma 0}^{(0)}(0)\left(\frac{m_c\overline{\psi }\psi _0^{(0)}}{m_0^2(0)}+\frac{q^2}{g_{\pi qq0}^{(0)\mathrm{\hspace{0.17em}2}}}\right).$$
(B9)
In the chiral limit, the integral in the sixth term of Eq. (A12) reduces to the one in the definition of $`C`$, Eq. (31), giving
$$\frac{g_{\pi qq0}^{(0)}}{q^2}C_0\left(1G_1J_{\sigma \sigma 0}^{(0)}(0)\right)\left(\frac{m_c\overline{\psi }\psi _0^{(0)}}{m_0^2(0)}+\frac{q^2}{g_{\pi qq0}^{(0)\mathrm{\hspace{0.17em}2}}}\right).$$
(B10)
The final term of Eq. (A12) is trivially
$$\frac{g_{\pi qq0}^{(1)}}{g_{\pi qq0}^{(0)}}\frac{m_0(0)}{g_{\pi qq0}^{(0)}}.$$
(B11)
Adding together Eqs. (B7) to (B11) reproduces exactly the $`𝒪(N_c^1)`$ terms of Eq. (B6), demonstrating that the Gell-Mann–Oakes–Renner relation holds in our treatment of the nonlocal NJL model at NLO.
## Tables
## Figures
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# References
Modern Physics Letters A,
c World Scientific Publishing Company
COMMENT ON EQUAL VELOCITY ASSUMPTION
FOR NEUTRINO OSCILLATIONS
L.B. OKUN<sup>*</sup><sup>*</sup>*E-mail: okun@heron.itep.ru, and I.S. TSUKERMANE-mail: zuckerma@heron.itep.ru
State Research Center – Institute for Theoretical and Experimental Physics
Moscow, 117259, Russia
Received (received date)
Revised (revised date)
The so-called equal velocity prescription for neutrino oscillations is forbidden kinematically.
In a recent article<sup>1</sup> it was suggested that when considering neutrino oscillations one can assume that neutrino mass eigenstates have equal velocities.
Independently and later the equal velocity scenario was suggested in Ref.<sup>2</sup>. The authors of Ref.<sup>2</sup> consider this scenario as “aesthetically the most pleasing”. They proclaimed it as their “preferred choice” in particular because it leads to the frequency of neutrino oscillations twice as large as the standard one.
Somewhat different approach is preferred by the authors of Ref.<sup>3</sup> who consider the equal velocity prescription on the same footing as the well known prescriptions of equal energy or equal momentum (see Refs.<sup>4-7</sup> and the literature therein). In particular, in Ref.<sup>3</sup> it is stated that none of the three prescriptions holds in pion decays or in any two-body decays.
The aim of this note is to stress that unlike the two “traditional” prescriptions, that of equal velocity is forbidden by simple kinematical considerations. This was shown in Ref.<sup>7</sup> which contains notes of a lecture given at the 1999 ITEP International Winter School. During that lecture Yu. Dokshitzer raised the issue of equal velocity case. The answer was that assuming $`v_1=v_2=v`$ one immediately arrives $`\gamma _1=\gamma _2=\gamma =1/\sqrt{1v^2}`$ and hence to $`E_1/E_2=m_1/m_2`$. This equality cannot hold, because $`E_1/E_21`$, while $`m_1/m_2`$ may be extremely small, or extremely large. We want to emphasize here that this reasoning is equally valid for two-body as well as multi-body decays.
Acknowledgments
We are grateful to M.I. Vysotsky for drawing our attention to Ref.<sup>1-2</sup>. This work was supported by RFBR grant 00-15-96562 and by A. von Humboldt award to one of us (L.O.).
References
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# Quantum Field Theoretic Treatment of the Non–Forward Compton Amplitude in the Generalized Bjorken Region
## 1 INTRODUCTION
Compton scattering of a virtual photon off a hadron $`\gamma _1^{}+p_1\gamma _{}^{}{}_{2}{}^{}+p_2`$ is an important processes in QCD: theoretically it can be treated in detail and experimentally it can be tested for a large variety of processes. In lowest approximation in the electromagnetic coupling it is described by
$`T_{\mu \nu }(p_+,p_{},q)=i{\displaystyle }d^4xe^{iqx}\times `$
$`p_2,S_2|T(J_\mu (x/2)J_\nu (x/2))|p_1,S_1,`$
where
$`p_\pm =p_2\pm p_1,q=\frac{1}{2}\left(q_1+q_2\right),`$
with $`q_1(q_2)`$ and $`p_1(p_2)`$ being the four–momenta of the incoming (outgoing) photon and hadron, respectively, and $`S_1,S_2`$ being the spins of the initial– and final–state hadron, where $`p_1+q_1=p_2+q_2`$. The generalized Bjorken region is the asymptotic domain being defined by
$`\nu =qp_+\mathrm{},q^2\mathrm{},`$
where the two scaling variables
$`\xi ={\displaystyle \frac{q^2}{qp_+}},\eta ={\displaystyle \frac{qp_{}}{qp_+}}={\displaystyle \frac{q_1^2q_2^2}{2\nu }}`$ (1.2)
are fixed. Of special experimental importance are the cases of deep inelastic scattering (DIS) described by the absorptive part of the forward Compton amplitude, $`\eta =0`$, and the deeply virtual Compton scattering (DVCS) with one real outgoing photon $`q_2^2=0`$ corresponding to $`\xi =\eta `$. In the generalized Bjorken region the amplitude (1) is dominated by the light–cone singularities which allows to apply the (non–local) operator product expansion of $`T(J_\mu (x/2)J_\nu (x/2))`$.
Our aim is a detailed quantum field theoretic investigation of that approach, cf. Refs. . In contrast to earlier considerations we take into account the explicit twist decomposition of the non–local vector operators, and their matrix elements, which thereby occur. The twist decomposition in the case of the quark–antiquark operators has been treated in . An extension to the gluon operators and more general multiparticle operators is given in . The relations between vector and scalar operators of twist 2 allows to express the final results with the help of matrix elements of the scalar operators only. In addition, this procedure leads to the derivation of new relations on the amplitude level which correspond in the case of forward scattering to the Callan-Gross and Wandzura-Wilczek relations. Furthermore, the electromagnetic current conservation may be shown to hold on the level of the twist 2 contributions.
## 2 LIGHT–CONE–EXPANSION
The Compton amplitude for the case of non–forward scattering, Eq. (1), in the generalized Bjorken region is dominated by the light–cone singularities. Therefore, the $`T`$–product of the electromagentic currents will be approximated by its non–local light–cone expansion – a summed–up form of the local light–cone expansion which allows for a quite compact representation of the resulting expressions.
Let us present a series of intuitive arguments leading to that approximation. We start from the renormalized time-ordered operator product ($`S`$ being the renormalized $`S`$–matrix):
$`\widehat{T}_{\mu \nu }(x)=iRT\left[J_\mu \left(x/2\right)J_\nu \left(x/2\right)S\right].`$
At first we consider this expression in the Born approximation
$`\widehat{T}_{\mu \nu }(x)=e^2{\displaystyle \frac{x^\lambda }{2\pi ^2(x^2iϵ)^2}}\times `$
$`\left[\overline{\psi }\left(\frac{x}{2}\right)\gamma _\mu \gamma _\lambda \gamma _\nu \psi \left(\frac{x}{2}\right)\overline{\psi }\left(\frac{x}{2}\right)\gamma _\nu \gamma _\lambda \gamma _\mu \psi \left(\frac{x}{2}\right)\right].`$
Here $`e`$ denotes the charge of the fermion field $`\psi `$. We dropped the flavor indices in the expressions considered. Reordering in the standard way the Dirac–structure we obtain
$`\widehat{T}_{\mu \nu }(x)=e^2{\displaystyle \frac{x^\lambda }{i\pi ^2(x^2iϵ)^2}}\times `$
$`\left[S_{\alpha \mu \lambda \nu }O^\alpha (\frac{x}{2},\frac{x}{2})i\epsilon _{\alpha \mu \lambda \nu }O_5^\alpha (\frac{x}{2},\frac{x}{2})\right],`$
where
$`S_{\alpha \mu \lambda \nu }=g_{\alpha \mu }g_{\lambda \nu }+g_{\lambda \mu }g_{\alpha \nu }g_{\mu \nu }g_{\lambda \alpha }.`$
The essential objects are the bilocal operators
$`O^\alpha (\frac{x}{2},\frac{x}{2})=`$ (2.2)
$`=\frac{i}{2}\left[\overline{\psi }\left(\frac{x}{2}\right)\gamma ^\alpha \psi \left(\frac{x}{2}\right)\overline{\psi }\left(\frac{x}{2}\right)\gamma ^\alpha \psi \left(\frac{x}{2}\right)\right],`$
$`O_5^\alpha (\frac{x}{2},\frac{x}{2})=`$ (2.3)
$`=\frac{i}{2}\left[\overline{\psi }\left(\frac{x}{2}\right)\gamma _5\gamma ^\alpha \psi \left(\frac{x}{2}\right)+\overline{\psi }\left(\frac{x}{2}\right)\gamma _5\gamma ^\alpha \psi \left(\frac{x}{2}\right)\right].`$
Expression (2) satisfies electromagnetic current conservation in the case of free fields, i.e. at zeroth order in QCD. We are interested, however, in the case of general fields $`\psi `$ and the twist–2 operators associated to them. We calculate these operators at leading order passing the following steps.
STEP 1: Use gauge invariant operators in place of (2.2), (2.3). This is achieved by including the phase factor $`U(y,z)=𝒫\mathrm{exp}(ig_z^yA_\mu 𝑑x^\mu )`$. The integration can be performed over a straight path connecting $`y`$ and $`z`$.
STEP 2: Perform the twist decomposition of these operators according to and restrict to the twist–2 (axial) vector operators only.
STEP 3: Take the operators as renormalized ones (at the light–cone) to all orders of QCD:
$`O_\alpha ^{\mathrm{tw2}}(\frac{\stackrel{~}{x}}{2},\frac{\stackrel{~}{x}}{2})=`$ (2.4)
$`=\frac{i}{2}RT\{[\overline{\psi }\left(\frac{\stackrel{~}{x}}{2}\right)\gamma _\alpha U(\frac{\stackrel{~}{x}}{2},\frac{\stackrel{~}{x}}{2})\psi \left(\frac{\stackrel{~}{x}}{2}\right)`$
$`\overline{\psi }\left(\frac{\stackrel{~}{x}}{2}\right)\gamma _\alpha U(\frac{\stackrel{~}{x}}{2},\frac{\stackrel{~}{x}}{2})\psi \left(\frac{x}{2}\right)]^{\mathrm{tw2}}S\},`$
$`O_{5\alpha }^{\mathrm{tw2}}(\frac{\stackrel{~}{x}}{2},\frac{\stackrel{~}{x}}{2})=`$ (2.5)
$`=\frac{i}{2}RT\{[\overline{\psi }\left(\frac{\stackrel{~}{x}}{2}\right)\gamma _5\gamma _\alpha U(\frac{\stackrel{~}{x}}{2},\frac{\stackrel{~}{x}}{2})\psi \left(\frac{\stackrel{~}{x}}{2}\right)`$
$`+\overline{\psi }\left(\frac{\stackrel{~}{x}}{2}\right)\gamma _5\gamma _\alpha U(\frac{\stackrel{~}{x}}{2},\frac{\stackrel{~}{x}}{2})\psi \left(\frac{\stackrel{~}{x}}{2}\right)]^{\mathrm{tw2}}S\},`$
where
$`\stackrel{~}{x}=x+\zeta \left[\sqrt{(x\zeta )^2x^2\zeta ^2}(x\zeta )\right]/\zeta ^2`$
and $`\zeta `$ is a subsidiary vector. Therefore our final expression to be considered in the following reads
$`\widehat{T}_{\mu \nu }^{\mathrm{tw2}}(x)=e^2{\displaystyle \frac{\stackrel{~}{x}^\lambda }{i\pi ^2(x^2iϵ)^2}}\times `$
$`\left[S_{\mu \lambda \nu }^\alpha O_\alpha ^{\mathrm{tw2}}(\frac{\stackrel{~}{x}}{2},\frac{\stackrel{~}{x}}{2})i\epsilon _{\mu \lambda \nu }^\alpha O_{5\alpha }^{\mathrm{tw2}}(\frac{\stackrel{~}{x}}{2},\frac{\stackrel{~}{x}}{2})\right].`$
## 3 TWIST DECOMPOSITION
We are confronted now with the twist decomposition of non–local operators. The concept of twist was originally introduced in and successfully applied in . The twist decomposition of simple non–local operators was studied in . for the first time. A unique group theoretical procedure, based on the decomposition of related local tensor operators into irreducible ones with respect to the Lorentz group $`O(3,1)`$ has been introduced recently and successfully applied to non–local tensor operators up to second rank .
As an example let us consider the following (uncentered, unsymmetrized) quark operator
$`O_\mathrm{\Gamma }(0,\kappa x)=\left[\overline{\psi }(0)\mathrm{\Gamma }U(0,\kappa x)\psi (\kappa x)\right],`$
where $`\mathrm{\Gamma }=\{1,\gamma _5;\gamma _\alpha ,\gamma _\alpha \gamma _5;\sigma _{\alpha \beta }\}`$. Its expansion into local operators reads
$`O_\mathrm{\Gamma }(0,\kappa x)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\kappa ^n}{n!}}x^{\mu _1}x^{\mu _2}\mathrm{}.x^{\mu _n}\times `$
$`[\overline{\psi }(0)\mathrm{\Gamma }D_{\mu _1}(y)D_{\mu _2}(y)\mathrm{}.D_{\mu _n}(y)\psi (y)]|_{y=0},`$
where $`D_\mu (y)`$ denotes the covariant derivative taken at $`y`$. Now, the local operators whose tensor structure is determined by $`\mathrm{\Gamma }`$, e.g., $`\alpha \beta (\mu _1\mathrm{}\mu _n)`$, and being totally symmetric with respect to $`\mu _1\mathrm{}\mu _n`$ are to be decomposed into irreducible tensors. These tensors should be traceless and their symmetry behavior is uniquely determined by Young patterns $`(m_1,m_2,\mathrm{}m_r)`$ whose $`i`$-th row has length $`m_i`$. In the (pseudo) scalar case the only allowed Young pattern is $`(n)`$, in the (axial) vector case there are two Young patterns $`(n+1)`$ and $`(n,1)`$, whereas in the antisymmetric resp. symmetric tensor case the Young patterns $`(n+1,1)`$ and $`(n,1,1)`$ resp. $`(n+2),(n+1,1)`$ and $`(n,2)`$ appear. The complete decomposition of the local tensors into irreducible ones besides the leading twist part contains also irreducible tensors of higher twist being related to the trace terms (e.g., in the symmetric tensor case contributions up to twist 6 occur, cf. ). The next step to be performed consists in resuming, according to (3), the towers (with respect to $`n`$) of the local operators with the same twist to non–local operators of definite twist which are tensorial harmonic functions. Let us remark that according to this definition the twist decomposition depends on the basic point $`y=0`$ where the local expansion is made. Finally, these non–local operators are to be projected onto the light–cone in order to obtain the twist decomposition of the light–ray operators we are seeking for.
Now we list those results of which are relevant for the present consideration:
(a) Twist–2 scalar quark operators for arbitrary positions $`\kappa _1x`$ and $`\kappa _2x`$ are given by
$`O^{\mathrm{tw2}}(\kappa _1x,\kappa _2x)=O(\kappa _1x,\kappa _2x)`$
$`+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle _0^1}𝑑t({\displaystyle \frac{1t}{t}})^{k1}{\displaystyle \frac{(x^2)^k\mathrm{}^k}{4^kk!(k1)!}}O(\kappa _1tx,\kappa _2tx),`$
with (compare Eq. (2.2))
$`O(\kappa _1x,\kappa _2x)=`$ (3.3)
$`=\frac{i}{2}RT\{[\overline{\psi }(\kappa _1x)(x\gamma )U(\kappa _1x,\kappa _2x)\psi (\kappa _2x)`$
$`\overline{\psi }(\kappa _2x)(x\gamma )U(\kappa _2x,\kappa _1x)\psi (\kappa _1x)]S\}.`$
The operator (3) being of leading twist, i.e., with all traces being subtracted, obeys the following important relation
$`\mathrm{}O^{\mathrm{tw2}}(\kappa _1x,\kappa _2x)=0.`$ (3.4)
Obviously, on the light–cone, $`x\stackrel{~}{x}`$, both operators, (3) and (3.3), coincide.
(b) Twist–2 vector quark operators are shown to be determined through the twist–2 scalar quark operator by
$`O_\alpha ^{\mathrm{tw2}}(\kappa _1x,\kappa _2x)={\displaystyle _0^1}𝑑\tau _\alpha O^{\mathrm{tw2}}(\kappa _1\tau x,\kappa _2\tau x),`$ (3.5)
which, after projecting onto the light–cone, is crucial for the reduction of the matrix elements of the operators (2.4, 2.5) to those of the corresponding (pseudo)scalar operators. The operator (3.5) satisfies the relations
$`\mathrm{}O_\alpha ^{\mathrm{tw2}}(\kappa _1x,\kappa _2x)=0,`$ (3.6)
$`^\alpha O_\alpha ^{\mathrm{tw2}}(\kappa _1x,\kappa _2x)=0.`$ (3.7)
The second of these relations is crucial for current conservation. Let us remark that on the light–cone these relations have to be written by using the interior derivative , namely
$`_\alpha d_\alpha (1+\stackrel{~}{x}\stackrel{~}{})\stackrel{~}{}_\alpha \frac{1}{2}\stackrel{~}{x}_\alpha \stackrel{~}{}^2\mathrm{with}d^2=0.`$
Quite analogous relations hold for the pseudoscalar as well as axial vector operators.
## 4 MATRIX ELEMENTS OF TWIST–2 OPERATORS
Our aim is to obtain an expression for the non–forward Compton amplitude. Up to now we considered the representation of the $`T`$–product of currents containing twist–2 non–local (axial) vector operators. The next step will be to perform matrix elements of that $`T`$–product which, because of Eq. (3.5), can be traced back to matrix elements of the (pseudo) scalar operators (3). Let us consider these matrix elements first. They decompose into two parts having a Dirac and a Pauli structure, respectively:
$`e^2p_2,S_2\left|O^{\mathrm{tw2}}(x/2,x/2)\right|p_1,S_1`$
$`=i\overline{u}(p_2,S_2)(\gamma x)u(p_1,S_1)\times `$
$`{\displaystyle Dze^{ixp(z)/2}f(z_1,z_2,p_ip_jx^2,p_ip_j,\mu _R^2)}`$
$`+i\overline{u}(p_2,S_2)(x\sigma p_{})u(p_1,S_1)\times `$
$`{\displaystyle Dze^{ixp(z)/2}g(z_1,z_2,p_ip_jx^2,p_ip_j,\mu _R^2)},`$
where $`(x\sigma p_{})x^\alpha \sigma _{\alpha \beta }p_{}^\beta `$ and $`p(z)p_1z_1+p_2z_2`$, $`\mu _R`$ denotes the renormalization scale and
$`Dz=\frac{1}{2}dz_1dz_2\theta (1z_1)\theta (1+z_1)\theta (1z_2)\theta (1+z_2).`$
The kinematic decomposition given above follows if one takes into account that the spinors $`u(p_i,S_i)`$ describing the hadrons satisfy the free Dirac equation. The functions $`f(z_1,z_2,(p_ip_j)x^2,(p_ip_j),\mu _R^2)`$ are the parton distribution amplitudes and $`z_i`$ are the momentum fractions. For brevity we drop the remaining variables. In the present approach we, moreover, set $`(p_i.p_j)0`$. Under these assumptions the relation (3.4) is also valid for the matrix elements:
$`\mathrm{}p_2,S_2|O^{\mathrm{tw2}}(x/2,x/2)|p_1,S_1=`$
$`p_2,S_2|\mathrm{}O^{\mathrm{tw2}}(x/2,x/2)|p_1,S_1=0.`$
Now, let us reconstruct the vector operator using Eq. (3.5). We obtain
$`e^2p_2,S_2\left|O_\alpha ^{\mathrm{tw2}}(x/2,x/2)\right|p_1,S_1`$ (4.3)
$`=i{\displaystyle }Dze^{ixp(z)/2}F(z_1,z_2)\times `$
$`[\overline{u}(p_2,S_2)\gamma _\alpha u(p_1,S_1)`$
$`\frac{i}{2}p_\alpha (z)\overline{u}(p_2,S_2)(\gamma x)u(p_1,S_1)]`$
$`+i{\displaystyle }Dze^{ixp(z)/2}G(z_1,z_2)\times `$
$`[\overline{u}(p_2,S_2)\sigma _{\alpha \beta }p_{}^\beta u(p_1,S_1)`$
$`\frac{i}{2}p_\alpha (z)\overline{u}(p_2,S_2)(x\sigma p_{})u(p_1,S_1)],`$
where
$`F(z_1,z_2)={\displaystyle _0^1}{\displaystyle \frac{d\lambda }{\lambda ^2}}f({\displaystyle \frac{z_1}{\lambda }},{\displaystyle \frac{z_2}{\lambda }}),`$ (4.4)
and an analogous representation connects $`G`$ to $`g`$. Moreover, similar representations between the corresponding functions $`f_5`$$`F_5`$$`g_5`$ and $`G_5`$ are valid for the operators containing $`\gamma _5`$. Note that also here the necessary conditions
$`\mathrm{}p_2,S_2|O_{(5)\alpha }^{\mathrm{tw2}}(x/2,x/2)|p_1,S_1=0`$ (4.5)
$`^\alpha p_2,S_2|O_{(5)\alpha }^{\mathrm{tw2}}(x/2,x/2)|p_1,S_1=0`$ (4.6)
are satisfied. The relations (4.4) are of central importance since they form the theoretical basis of the Wandzura–Wilczek relations and allow expressions for the non–forward Compton amplitude based on expectation values of scalar operators only.
## 5 CURRENT CONSERVATION
Current conservation is a very important criterion for the relevance of the derived expressions. Formally we have to start with Eq. (1), and apply Eq. (2) and the expressions for the matrix elements (4.3). In the case of forward scattering current conservation holds .
For non–forward scattering one is confronted with the following problem. Let us consider the asymmetric current product with respect to $`x=0`$
$`\widehat{T}_{\mu \nu }^{\mathrm{tw2}}(\kappa _1x,\kappa _2x)=e^2{\displaystyle \frac{x^\lambda }{i\pi ^2(x^2iϵ)^2}}\times `$ (5.1)
$`\left[S_{\mu \lambda \nu }^\alpha O_\alpha ^{\mathrm{tw2}}(\kappa _1x,\kappa _2x)i\epsilon _{\mu \lambda \nu }^\alpha O_{5\alpha }^{\mathrm{tw2}}(\kappa _1x,\kappa _2x)\right],`$
where $`\kappa _1\kappa _2=1`$. It is easy to convince oneself that, independent of the values of $`\kappa _i`$,
$`_x^\mu \widehat{T}_{\mu \nu }(\kappa _1x,\kappa _2x)=0=_x^\nu \widehat{T}_{\mu \nu }(\kappa _1x,\kappa _2x),`$ (5.2)
holds because the relations (3.6) and (3.7) are satisfied. These relations, ensuring tracelessness of the vector operators of definite twist, are not changed by perturbation theory and renormalization.
This proves conservation of the first (or second) current if $`\kappa _1=1,\kappa _2=0`$ (or $`\kappa _1=0,\kappa _2=1`$) is chosen. However, if conservation of both electromagnetic currents simultaneously shall be proven we have to study
$`iRT[J_\mu (x)J_\nu (y)S]^{\mathrm{tw2}}=e^2{\displaystyle \frac{\xi ^\lambda }{i\pi ^2(\xi ^2iϵ)^2}}\times `$
$`[S_{\mu \lambda \nu }^\alpha O_\alpha ^{\mathrm{tw2}}(\eta +\frac{\xi }{2},\eta \frac{\xi }{2})`$
$`iϵ_{\mu \lambda \nu }^\alpha O_{5\alpha }^{\mathrm{tw2}}(\eta +\frac{\xi }{2},\eta \frac{\xi }{2})],`$ (5.3)
where $`\eta =(x+y)/2`$ by convention denotes the reference point for the twist decomposition, and $`\xi =xy`$ approaches the light–cone. Because of
$`O_{(5)\alpha }^{\mathrm{tw2}}(\eta +\frac{\xi }{2},\eta \frac{\xi }{2}):=e^{i\eta P}O_{(5)\alpha }^{\mathrm{tw2}}(\frac{\xi }{2},\frac{\xi }{2})e^{i\eta P},`$
where $`P_\mu `$ is the momentum operator, we find that Eqs. (5.2) hold with respect to the variable $`\xi `$ which is part of $`x`$ and $`y`$. Therefore, applying both derivations, either $`\frac{}{x^\mu }=\frac{1}{2}\frac{}{\eta ^\mu }+\frac{}{\xi ^\mu }`$ or $`\frac{}{y^\mu }=\frac{1}{2}\frac{}{\eta ^\mu }\frac{}{\xi ^\mu }`$, to the expression (LABEL:Txy) there remains, in both cases, a non–vanishing part which is proportional to $`[P_\mu ,O_{(5)\alpha }^{\mathrm{tw2}}(\eta +\frac{\xi }{2},\eta \frac{\xi }{2})]`$.
This shows that the proof of current conservation essentially depends on the choice of the reference point for the twist definition: translation of the non–linear operator also shifts that reference point. It prevented us from proving conservation of both currents simultaneously. This intrinsic problem of all the twist definitions could be very important when non–leading twist contributions are considered.
Let us turn to the case of twist 2 now. For the explicit calculations of the resulting expressions we use a normalized helicity basis, cf. , with $`\epsilon _{0\mu }^{(i)}=q_{i\mu }/\sqrt{|q_i^2|}`$ for $`q_i^2<0`$ resp. $`\epsilon _{0\mu }^{(i)}=q_{i\mu }/(\sqrt{2}|q_{0i}|)`$ for $`q_i^2=0`$. We have shown in Ref. , that the current violating contributions are of $`O(\nu ^{1/2}\times \mathrm{OME})`$ or higher order. These terms are of higher twist and have to be dealt with the operator matrix elements of the higher twist operators. <sup>1</sup><sup>1</sup>1C. Weiss has proven recently using the representation of that the terms $`O(\nu ^{1/2})`$ cancel with corresponding terms due to twist–3 operators.
## 6 INTEGRAL RELATIONS
Here we present the final expression for the Compton amplitude. The calculations are given in . We split the amplitude into its symmetric part and the antisymmetric part $`T^{\mu \nu }=T_\mathrm{s}^{\mu \nu }+T_{\mathrm{as}}^{\mu \nu }`$. In the case of forward scattering the former one corresponds to the unpolarized and the latter one to the polarized contribution. First we consider the symmetric part:
$`T_\mathrm{s}^{\mu \nu ,\mathrm{tw2}}={\displaystyle \frac{1}{\nu }}{\displaystyle _1^1}dt{\displaystyle \frac{1}{\xi +ti\epsilon }}\times `$
$`\{[2(g^{\mu \nu }(qp_+)(q^\mu p_+^\nu +q^\nu p_+^\mu ))f_1(t,\eta )`$
$`+p_+^\mu p_+^\nu f_2(t,\eta )]\overline{u}(p_2,S_2)(\gamma q)u(p_1,S_1)`$
$`+[2(g^{\mu \nu }(qp_+)(q^\mu p_+^\nu +q^\nu p_+^\mu ))g_1(t,\eta )`$
$`+p_+^\mu p_+^\nu g_2(t,\eta )]\overline{u}(p_2,S_2)(q\sigma p_{})u(p_1,S_1)\}`$
$`+\mathrm{non}\mathrm{leading}\mathrm{terms},`$
where $`t=z_++\eta z_{},z_\pm =\frac{1}{2}(z_2\pm z_1)`$. Here the partition functions $`f_i`$ and $`g_i`$ are ‘one–variable’ distribution amplitudes which are defined by
$`f_{(5)}(t,\eta )`$ $`=`$ $`{\displaystyle 𝑑z_{}f_{(5)}(z_+=t\eta z_{},z_{})},`$ (6.2)
$`g_{(5)}(t,\eta )`$ $`=`$ $`{\displaystyle 𝑑z_{}g_{(5)}(z_+=t\eta z_{},z_{})},`$ (6.3)
from the ‘two–variable’ distribution amplitudes used in the representation (4) of the matrix elements of the (pseudo) scalar operators. Unlike the case of forward scattering these functions do not depend on scaling variables only but besides of the scaling variable $`\eta `$ which describes non–forwardness of the combination of momentum fractions $`t`$. The following new relations are obtained between the amplitude–functions $`f_i(g_i)`$, see :
$`f_2(t,\eta )`$ $`=`$ $`2tf_1(t,\eta )2f(t,\eta ),`$ (6.4)
$`g_2(t,\eta )`$ $`=`$ $`2tg_1(t,\eta )2g(t,\eta ).`$ (6.5)
These relations are structurally similar to the Callan–Gross relation for forward scattering. There, by virtue of the optical theorem, $`\frac{1}{\xi +t+iϵ}i\pi \delta (t+\xi )`$ and for $`p_2p_1=p`$ it follows $`tz_+z`$, $`t`$ is turned into a scaling variable. In the above expressions furthermore
$`\overline{u}(p_2,S_2)(\gamma q)u(p_1,S_1)`$ $``$ $`2pq,`$
$`\overline{u}(p_2,S_2)(q\sigma p_{})u(p_1,S_1)`$ $``$ $`0`$
holds in the latter case.
The result for the antisymmetric part is more complicated:
$`T_{(\lambda _1)(\lambda _2)}`$ $`=`$ $`ϵ_{\mu (\lambda _1)}^{(2)}ϵ_{\nu (\lambda _2)}^{(1)}T_{\mathrm{as}}^{\mu \nu ,\mathrm{tw2}}`$ (6.6)
$`=`$ $`iϵ^{\mu \rho \nu \sigma }ϵ_{\mu (\lambda _1)}^{(2)}ϵ_{\nu (\lambda _2)}^{(1)}B_{\rho \sigma }`$
with
$`B^{\rho \sigma }={\displaystyle \frac{q^\rho }{\nu ^2}}[{\displaystyle _1^1}dt{\displaystyle \frac{1}{\xi +ti\epsilon }}\times `$
$`\{(f_{5,1}(t,\eta )+f_{5,2}(t,\eta ))\nu S_\sigma ^{12}`$
$`+\left(g_{5,1}(t,\eta )+g_{5,2}(t,\eta )\right)\nu \mathrm{\Sigma }_\sigma ^{12}`$
$`+f_{5,2}(t,\eta )p_+^\sigma (qS^{12})+g_{5,2}(t,\eta )p_+^\sigma (q\mathrm{\Sigma }^{12})\}].`$
Here, we used the following abbreviations for the spinor structure
$`S_\sigma ^{12}`$ $`=`$ $`\frac{1}{2}\overline{u}(p_2,S_2)\gamma _5\gamma _\sigma u(p_1,S_1),`$
$`\mathrm{\Sigma }_\sigma ^{12}`$ $`=`$ $`\frac{1}{2}\overline{u}(p_2,S_2)\gamma _5\sigma _{\sigma \rho }p_{}^\rho u(p_1,S_1).`$
In the case of forward scattering $`p_2p_1=p`$, $`S_\sigma ^{12}S_\sigma `$ and $`\mathrm{\Sigma }_\sigma ^{12}0`$ holds, where $`S_\sigma `$ is the spin vector introduced for forward scattering. We obtain the following relations for the amplitude functions, cf. :
$`f_{5,1}(t,\eta )`$ $``$ $`f_5(t,\eta ),`$ (6.7)
$`f_{5,2}(t,\eta )`$ $`=`$ $`f_5(t,\eta )+{\displaystyle _t^{\mathrm{sgn}t}}𝑑z{\displaystyle \frac{f_5(z,\eta )}{z}},`$ (6.8)
$`g_{5,1}(t,\eta )`$ $``$ $`g_5(t,\eta ),`$ (6.9)
$`g_{5,2}(t,\eta )`$ $`=`$ $`g_5(t,\eta )+{\displaystyle _t^{\mathrm{sgn}t}}𝑑z{\displaystyle \frac{g_5(z,\eta )}{z}}.`$ (6.10)
Again $`f_{5,i}`$ and $`g_{5,i}`$ depend on the momentum fraction $`t`$ and a scaling variable. These relations generalize the Wandzura–Wilczek relation of deep inelastic scattering to non–forward amplitudes.
## 7 CONCLUSIONS
We studied the structure of the virtual Compton amplitude for deep–inelastic non–forward scattering $`\gamma ^{}+p\gamma ^{}+p^{}`$ in lowest order in QED in the massless limit. In the generalized Bjorken region $`(qp_+),q^2\mathrm{}`$ the twist–2 contributions to the Compton amplitude were calculated using the non–local operator product expansion. The twist separation and the relations between twist–2 vector operators and twist 2 scalar operators are essential for the current conservation at the level of twist–2 and, moreover, all parton distributions are connected with the matrix elements of the scalar twist–2 operators.
The relations between the twist–2 contributions of the unpolarized and polarized amplitude functions were derived. They are the non–forward generalizations of the CallanGross and WandzuraWilczek relations for unpolarized and polarized deep–inelastic forward scattering. The relations for the Dirac and Pauli parts are of the same form.
Acknowledgement. This work was supported in part by EU contract FMRX-CT98-0194 (DG 12–MIHT). B.G. and D.R. would like to thank DESY for financial support.
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# Maxwell stress tensor and the Casimir effect
## 1 Introduction
According to Casimir , the (macroscopically) observable vacuum energy of a quantum field is the regularized difference between the zero point energies with and without the external conditions demanded by the particular physical situation at hand. In the case of the quantized electromagnetic field confined between two infinite parallel conducting plates separated by a distance $`a`$, Casimir’s conception of the vacuum energy leads to a force per unit area between the plates given by
$$\frac{F}{A}=\frac{\pi \mathrm{}c}{240a^4}.$$
(1)
Until 1997 only one experiment involving Casimir’s original setup had been performed . Recently, however, the experimental observation of this tiny force with metallic surfaces was significantly improved by the experiments due to Lamoreaux and to Mohideen and Roy .The concept of an observable vacuum energy can be extended to all quantum fields and several types of boundary conditions and/or applied external fields. A review of all these Casimir effects can be found in, for example, Mostepanenko and Trunov or in Plunien et al. .
The local approach to the electromagnetic Casimir effect was initiated by Brown and Maclay who calculated the renormalized stress energy tensor between two parallel perfectly conducting plates by means of Green functions techniques . An interesting approach to the standard Casimir effect is the one due to Gonzales . This author pointed out that the apparently non-objectionable definition of the vacuum energy given above could easily lead to conceptual errors and stressed the fact that in any Casimir interaction calculation, contributions from both sides of the material surfaces involved must be taken into account. This is so because the vacuum pressure always pushes the material surfaces involved, therefore, the repulsiveness or atractiveness of the Casimir force depends on the discontinuity of the relevant component of the quantized Maxwell stress tensor at the location of the surface.
The purpose of this paper is to pursue this line of reasoning by analyzing the stresses on parallel material surfaces due to the vacuum distortions caused by the presence of these surfaces through the quantized Maxwell stress tensor. Though the approach chosen here has many points in common with Ref. cited above, it is a different alternative in the sense that it relies on objects known as correlators, which are vacuum expectation values of products of field components taken at the same point in space and time. These correlators contain all the information we need on the local behavior of the vacuum expectation values of elements of Maxwell stress tensor, in particular, their behavior near both sides of the material surface in question, a feature crucial to the obtention of the correct result. The local behavior of the Maxwell tensor, or of the relativistic symmetrical stress-energy tensor, is extremely important because as shown by, for example, Deutsch and Candelas with the help of Green functions technique, as we approach the boundaries we find strong divergencies that cannot be removed by renormalization. Besides reviewing the obtention of the electromagnetic Casimir for the standard case of two perfectly conducting parallel plates, we will also consider a pair of parallel plates, one of them perfectly permeable. This setup was first proposed by Boyer who analyzed them from the viewpoint of random electrodymanics and it is the simplest example of a repulsive Casimir force. For both cases, the conducting plate and Boyer’s setup we also construct the symmetrical stress tensor.
It is well-known that an atom can be attracted or repelled by a material surface, thus probing locally the vacuum distortions, for these reasons we take advantage of the knowledge of the above correlators and include a derivation of the Casimir-Polder interaction between an atom and a material surface. We will employ gaussian units and set $`c=\mathrm{}=1`$.
## 2 Maxwell stress tensor and the electromagnetic field correlators
Our aim in this section is to obtain an expression for the quantum version of the electromagnetic force per unit area that acts on plane material surface. Material surface here means a perfectly conducting square surface($`ϵ\mathrm{}`$) or a perfectly permeable one ($`\mu \mathrm{})`$ whose linear dimension $`L`$ is much larger than others relevant dimensions envolved such as the distance between two of those surfaces. The physical interaction between any of the two types of surfaces considered here and the vacuum electromagnetic field is mimicked by the imposition of appropriate boundary conditions on the electromagnetic field on the location of the material surface. The Cartesian components of the Maxwell stress tensor in Gaussian units are given by
$$T_{ij}=\frac{1}{4\pi }\left(E_iE_j\frac{1}{2}\delta _{ij}𝐄^2+B_iB_j\frac{1}{2}\delta _{ij}𝐁^2\right)$$
(2)
where $`i,j=x,y,z`$. Suppose that the material surface is placed perpendicularly to the $`𝒪𝒵`$ axis. Upon quantizing the electromagnetic field we can write the quantum version of (2). For instance,
$$\widehat{T}_{zz}_0=\frac{1}{8\pi }\left[\widehat{E}_z^2_0\widehat{E}_{}^2_0+\widehat{B}_z^2_0\widehat{B}_{}^2_0\right],$$
(3)
where $`\widehat{E}_{}^2=\widehat{E}_x^2+\widehat{E}_y^2`$ and $`\widehat{B}_{}^2=\widehat{B}_x^2+\widehat{B}_y^2`$; $`\widehat{O}_00\left|\widehat{O}\right|0`$ denotes a vacuum expectation value. The other Cartesian components of this tensor can be obtained in an analogous way. The quantum macroscopic force$`\widehat{𝐅}_0`$ on the material surface can be evaluated by integrating the quantum version of the classical result
$$𝐅=_{}\stackrel{~}{𝐓}\widehat{𝐧}𝑑a,$$
(4)
were $`\widehat{𝐧}`$ is outwards normal at $``$ and $``$ is any region containing the material surface. Classically, (4) can be obtained by integrating the Lorentz force per unit volume acting on charge and current distributions and eliminating the sources in favor of the fields. From a quantum point of view, we see that the problem of evaluating the pressure the material surface due to the distorted zero point oscillations of the electromagnetic field is reduced to the evaluation of the vacuum expectation value of the quantum operators $`\widehat{E}_i(𝐫,t)\widehat{E}_j(𝐫,t)`$, $`\widehat{B}_i(𝐫,t)\widehat{B}_j(𝐫,t),`$ and $`\widehat{E}_i(𝐫,t)\widehat{B}_j(𝐫,t)`$. The evaluation of these correlators depends on the specific choice of the boundary conditions. A regularization recipe will also be necessary, for these objects are mathematically ill-defined. For the setup envolving conducting plates these correlators were evaluated by Lütken and Ravndal , see also . They can be also obtained from the coincidence limit of the photon propagator between conducting plates evaluated by Bordag et al . In the next section we will display a method of evaluating these correlators by means of analitycal continuation techniques similar to the ones employed by Lütken e Ravndal . We will also show how to obtain the corresponding results for another unusual but intersting setup .
## 3 Correlators for Casimir’s setup
Consider an experimental setup consisting in two infinite perfectly conducting parallel plates ($`ϵ\mathrm{}`$) kept at a fixed distance $`a`$ from each other. We will choose the coordinates axis in such a way that the $`𝒪𝒵`$ direction is perpendicular to the plates. One of the plates will be placed at $`z=0`$ and the othrt one at $`z=a`$. The field must satisfy the following boundary conditions on the plates: the tangential components $`E_x`$ e $`E_y`$ of the electric field and the normal component $`B_z`$ of the magnetic field must be zero on the plates. Since there are no real charges or currents it will be convenient to work in the Coulomb gauge in which $`𝐀(𝐫,t)=0`$, and $`\mathrm{\Phi }=0`$, thus $`𝐄(𝐫,t)=𝐀(𝐫,t)/t`$ and $`𝐁(𝐫,t)=\times 𝐀(𝐫,𝐭)`$. These physical boundary conditions combined with the choice of gauge allow us to rewrite the boundary conditions in terms of the components of the vector potential $`𝐀(𝐫,t)`$ in the following way: at $`z=0`$ we will have,
$$A_x(x,y,0,t)=0;A_y(x,y,0,t)=0;\frac{}{z}A_z(x,y,0,t)=0,$$
(5)
and at $`z=a`$ ,
$$A_x(x,y,a,t)=0;A_y(x,y,a,t)=0;\frac{}{z}A_z(x,y,a,t)=0.$$
(6)
The vector potential operator $`\widehat{𝐀}(𝐫,t)`$ that satisfies the wave equation, the Coulmb gauge and the boundary conditions can be written as
$`\widehat{𝐀}(𝐫,t)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\left({\displaystyle \frac{\pi }{a}}\right)^{\frac{1}{2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}\prime \prime }{}}}{\displaystyle }{\displaystyle \frac{d^2\kappa }{\sqrt{\omega }}}\{\widehat{a}^{(1)}(\kappa ,n)\widehat{\kappa }\times \widehat{𝐳}\mathrm{sin}\left({\displaystyle \frac{n\pi z}{a}}\right)`$ (7)
$`+`$ $`\widehat{a}^{(2)}(\kappa ,n)[\widehat{\kappa }{\displaystyle \frac{in}{\omega a}}\mathrm{sin}\left({\displaystyle \frac{n\pi z}{a}}\right)\widehat{𝐳}{\displaystyle \frac{\kappa }{\omega }}\mathrm{cos}\left({\displaystyle \frac{n\pi z}{a}}\right)]\}e^{i(\kappa \rho \omega t)}+h.c,`$
where $`\kappa =(k_x,k_y)`$ and $`\rho `$ is the position arrow on the $`𝒳𝒴`$ plane. The normal frquencies are given by
$$\omega =\omega (\kappa ,n)=\sqrt{\kappa ^2+n^2\frac{\pi ^2}{a^2}},$$
(8)
with $`k_x,k_y`$ e $`n1`$. The symbol $`^\mathrm{"}`$ indicates that the term corresponding to $`n=0`$ for normaliztion reasons must be multiplied by $`1/2`$ . The Fourier coefficients $`\widehat{a}^{(\lambda )}(\kappa ,n)`$ where $`\lambda =1,2`$ is the polarization index, are operators in the photon ocupation number space and satisfy
$$[\widehat{a}^{(\lambda )}(\kappa ,n),\widehat{a}^{(\lambda ^{})}(\kappa ^{},n^{})]=\delta _{\lambda \lambda ^{}}\delta _{nn^{}}\delta \left(\kappa \kappa ^{}\right).$$
(9)
It is convenient to write the vector potential in the general form
$$\widehat{𝐀}(𝐫,t)=\underset{n=0}{\overset{\mathrm{}\prime \prime }{}}d^2\kappa \underset{\lambda =1}{\overset{2}{}}\widehat{a}^{(\lambda )}(\kappa ,n)𝐀_{\kappa n}^{(\lambda )}(𝐫)e^{i\omega (\kappa ,n)t}+h.c,$$
(10)
where $`𝐀_{\kappa n}^{(\lambda )}(𝐫)`$ are the modal functions. The modal functions for each polarization state must obey Helmholtz equation and the boundary conditions given above. In our case the modal functions are
$$𝐀_{\kappa n}^{(1)}(𝐫)=\frac{1}{\pi }\left(\frac{\pi }{a}\right)^{\frac{1}{2}}\frac{1}{\sqrt{\omega }}\mathrm{sin}\left(\frac{n\pi z}{a}\right)e^{i\kappa \rho }\widehat{\kappa }\times \widehat{𝐳},$$
(11)
and
$$𝐀_{\kappa n}^{(2)}(𝐫)=\frac{1}{\pi }\left(\frac{\pi }{a}\right)^{\frac{1}{2}}\frac{1}{\sqrt{\omega }}\left[\widehat{\kappa }\frac{in\pi }{a\omega }\mathrm{sin}\left(\frac{n\pi z}{a}\right)\widehat{𝐳}\frac{\kappa }{\omega }\mathrm{cos}\left(\frac{n\pi z}{a}\right)\right]e^{i\kappa \rho }.$$
(12)
The next step is to evaluate the electric field operator $`\widehat{𝐄}(𝐫,t)`$. Recalling that $`\widehat{a}^{(\lambda )}(\kappa ,n)|0=0,`$ we first write the correlators $`<\widehat{E}_i(𝐫,t)\widehat{E}_j(𝐫,t)_0`$ in the general form
$$<\widehat{E}_i(𝐫,t)\widehat{E}_j(𝐫,t)_0=\underset{\alpha }{}E_{i\alpha }(𝐫)E_{j\alpha }^{}(𝐫),$$
(13)
where we have introduced the modal functions $`E_{i\alpha }(𝐫)`$ for the electric field. In our case (11) and (12) yield
$$𝐄_{i\kappa n}^{(1)}(𝐫)=\frac{i}{\pi }\left(\frac{\omega (\kappa ,n)\pi }{a}\right)^{\frac{1}{2}}\mathrm{sin}\left(\frac{n\pi z}{a}\right)e^{i\kappa \rho }(\widehat{\kappa }\times \widehat{𝐳})_i,$$
(14)
and,
$$𝐄_{i\kappa n}^{(2)}(𝐫)=\frac{i}{\pi }\left(\frac{\omega (\kappa ,n)\pi }{a}\right)^{\frac{1}{2}}\left[\kappa _i\frac{in\pi }{a\omega (\kappa ,n)}\mathrm{sin}\left(\frac{n\pi z}{a}\right)\widehat{𝐳}_i\frac{\kappa }{\omega (\kappa ,n)}\mathrm{cos}\left(\frac{n\pi z}{a}\right)\right]e^{i\kappa \rho },$$
(15)
respectively. Taking (14) and (15) into (13), we write $`\widehat{\kappa }_i=\mathrm{cos}\varphi \delta _{ix}+\mathrm{sin}\varphi \delta _{iy}`$, $`\widehat{𝐳}_i=\delta _{iz}`$ e $`(\widehat{𝐳}\times \widehat{\kappa })_i=\mathrm{sin}\varphi \delta _{ix}\mathrm{cos}\varphi \delta _{iy}`$, where $`\varphi `$ is the azimuthal angle on the $`𝒳𝒴`$ plane and we have performed all angular integrals. In this way we end up with
$`\widehat{E}_i(𝐫,t)\widehat{E}_j(𝐫,t)_0=\left({\displaystyle \frac{2}{\pi }}\right)\left({\displaystyle \frac{\pi }{a}}\right){\displaystyle \frac{\delta _{ij}^{}}{2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}\prime \prime }{}}}\mathrm{sin}^2\left({\displaystyle \frac{n\pi z}{a}}\right){\displaystyle _0^{\mathrm{}}}𝑑\kappa \kappa \omega (\kappa ,n)`$ (16)
$`+`$ $`\left({\displaystyle \frac{2}{\pi }}\right)\left({\displaystyle \frac{\pi }{a}}\right)\left({\displaystyle \frac{\pi }{a}}\right)^2{\displaystyle \frac{\delta _{ij}^{}}{2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}\prime \prime }{}}}n^2\mathrm{sin}^2\left({\displaystyle \frac{n\pi z}{a}}\right){\displaystyle _0^{\mathrm{}}}𝑑\kappa \kappa \omega ^1(\kappa ,n)`$
$`+`$ $`\left({\displaystyle \frac{2}{\pi }}\right)\left({\displaystyle \frac{\pi }{a}}\right)\delta _{ij}^{}{\displaystyle \underset{n=0}{\overset{\mathrm{}\prime \prime }{}}}\mathrm{cos}^2\left({\displaystyle \frac{n\pi z}{a}}\right){\displaystyle _0^{\mathrm{}}}𝑑\kappa \kappa ^3\omega ^1(\kappa ,n),`$
where $`\delta _{ij}^{}:=\delta _{ix}\delta _{jx}+\delta _{iy}\delta _{jy}`$ e $`\delta _{ij}^{}:=\delta _{iz}\delta _{jz}`$. Equation (16) is a formal expression for the correlator $`E_i(𝐫,t)E_j(𝐫,t)_0`$, since it is mathematically ill-defined unless a regularization recipe is prescribed. We will regularize the integrals in (16) with the help of analytical continuation methods.
Consider, for instance, the first integral on the r.h.s. of (16) and let us rewrite it as follows
$$_0^{\mathrm{}}𝑑\kappa \kappa \left(\kappa ^2+\frac{n^2\pi ^2}{a^2}\right)^{1/2}_0^{\mathrm{}}𝑑\kappa \kappa \left(\kappa ^2+\frac{n^2\pi ^2}{a^2}\right)^{1/2s}.$$
The first term in (16) can be rewritten as
$$T_1=\left(\frac{2}{\pi }\right)\left(\frac{\pi }{a}\right)\frac{\delta _{ij}^{}}{2}\underset{n=0}{\overset{\mathrm{}\mathrm{"}}{}}\mathrm{sin}^2\left(\frac{n\pi z}{a}\right)_0^{\mathrm{}}𝑑\kappa \kappa \left(\kappa ^2+\frac{n^2\pi ^2}{a^2}\right)^{\frac{1}{2}s},$$
(17)
the second as
$$T_2=\left(\frac{2}{\pi }\right)\left(\frac{\pi }{a}\right)\left(\frac{\pi }{a}\right)^2\frac{\delta _{ij}^{}}{2}\underset{n=0}{\overset{\mathrm{}\mathrm{"}}{}}n^2\mathrm{sin}^2\left(\frac{n\pi z}{a}\right)_0^{\mathrm{}}𝑑\kappa \kappa \left(\kappa ^2+\frac{n^2\pi ^2}{a^2}\right)^{\frac{1}{2}s},$$
(18)
and the third one as
$$T_3=\left(\frac{2}{\pi }\right)\left(\frac{\pi }{a}\right)\delta _{ij}^{}\underset{n=0}{\overset{\mathrm{}\mathrm{"}}{}}\mathrm{cos}^2\left(\frac{n\pi z}{a}\right)_0^{\mathrm{}}𝑑\kappa \kappa ^3\left(\kappa ^2+\frac{n^2\pi ^2}{a^2}\right)^{\frac{1}{2}s}.$$
(19)
Let us assume that $`\mathrm{}s`$ is large enough to give precise mathematical meaning to these integrals. After evaluating them and make use of the analytical continuation of the results we will take the limit $`s0`$. Let us see, for instance, what happens with $`T_1`$. Making use of the following representation of Euler beta function
$$_0^{\mathrm{}}𝑑xx^{\mu 1}\left(x^2+c^2\right)^{\nu 1}=\frac{B}{2}(\frac{\mu }{2},1\nu \frac{\mu }{2})c^{\mu +2\nu 2},$$
(20)
where $`B(x,y)=\mathrm{\Gamma }(x)\mathrm{\Gamma }(y)/\mathrm{\Gamma }(x+y)`$, and that holds for $`\mathrm{}\left(\nu +\frac{\mu }{2}\right)<1`$ and $`\mathrm{}\mu >0`$, we obtain
$$_0^{\mathrm{}}𝑑\kappa \kappa \left(\kappa ^2+\frac{n^2\pi ^2}{a^2}\right)^{1/2s}=\frac{1}{2}\left(\frac{n\pi }{a}\right)^{32s}\frac{\mathrm{\Gamma }(s3/2)}{\mathrm{\Gamma }(s1/2)}=\frac{1}{(2s3)}\left(\frac{n\pi }{a}\right)^{32s}.$$
(21)
Taking this result into $`T_1`$ we obtain
$$T_1=\left(\frac{1}{2s3}\right)\left(\frac{\pi }{a}\right)^{32s}\frac{\delta _{ij}^{}}{2a}\left[\zeta _R(2s3)\underset{n=0}{\overset{\mathrm{}}{}}n^{32s}\mathrm{cos}\left(\frac{2n\pi z}{a}\right)\right],$$
(22)
where $`\zeta _R(z)`$ is the well-known Riemann zeta function. Taking the limit $`s0`$, we have
$$T_1=\frac{1}{3\pi }\left(\frac{\pi }{a}\right)^4\frac{\delta _{ij}^{}}{2}\left[\frac{1}{120}+\frac{1}{8}\underset{n=1}{\overset{\mathrm{}}{}}\frac{d^3}{d\xi ^3}\mathrm{sin}\left(2n\xi \right)\right],$$
(23)
where we made use of the fact that $`\zeta _R(3)=1/120`$, defined $`\xi :=\pi z/a`$, and wrote
$$n^3\mathrm{cos}\left(2n\xi \right)=\frac{1}{8}\times \frac{d^3}{d\xi ^3}\mathrm{sin}\left(2n\xi \right).$$
(24)
The sum on the R.H.S. of (23) can be regularized in many ways. A quick though non-rigorous way is to write
$$T_1=\frac{1}{3\pi }\left(\frac{\pi }{a}\right)^4\frac{\delta _{ij}^{}}{2}\left[\frac{1}{120}+\frac{1}{8}\frac{d^3}{d\xi ^3}\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{sin}\left(2n\xi \right)\right],$$
(25)
and express the summand in terms of exponential functions of imaginary argument thereby transforming each one of the sums into the euclidean space. In this way we obtain
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{sin}\left(2n\xi \right)`$ $`=`$ $`{\displaystyle \frac{1}{2i}}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}\left(i2n\xi \right){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}\left(i2n\xi \right)\right)`$ (26)
$`=`$ $`{\displaystyle \frac{1}{2i}}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}\left(2n\xi _E\right){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{exp}\left(2n\xi _E^{}\right)\right)`$
where we have made the substitution $`i\xi \xi _E`$ in the first sum and $`i\xi \xi _E^{}`$ in the second. Each one of the sums above can be easily performed and the result is
$$\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{sin}\left(2n\xi \right)=\frac{1}{2}\mathrm{cot}\left(\xi \right)$$
(27)
It follows that
$$T_1=\frac{1}{3\pi }\left(\frac{\pi }{a}\right)^4\frac{\delta _{ij}^{}}{2}\left[\frac{1}{120}+\frac{1}{8}\frac{d^3}{d\xi ^3}\frac{1}{2}\mathrm{cot}\left(\xi \right)\right].$$
(28)
Treating the two other terms in (16) in a similar manner we obtain
$$T_2=\frac{1}{\pi }\left(\frac{\pi }{a}\right)^4\frac{\delta _{ij}^{}}{2}\left[\frac{1}{120}+\frac{1}{8}\frac{d^3}{d\xi ^3}\frac{1}{2}\mathrm{cot}\left(\xi \right)\right],$$
(29)
and
$$T_3=\frac{4}{3\pi }\left(\frac{\pi }{a}\right)^4\frac{\delta _{ij}^{}}{2}\left[\frac{1}{120}\frac{1}{8}\frac{d^3}{d\xi ^3}\frac{1}{2}\mathrm{cot}\left(\xi \right)\right].$$
(30)
Notice that some care must be taken when we aply this procedure to the third term. This is so because the term corresponding to $`n=0`$ in $`T_3`$ is not zero. In fact its contribution is:
$$T_3(n=0)=\left(\frac{2}{\pi }\right)\left(\frac{\pi }{a}\right)\delta _{ij}^{}_0^{\mathrm{}}𝑑\kappa \kappa ^{22s},$$
(31)
which diverges when the regularization is removed. How ever this term is non-physical and can be safely ignored. Finally, collecting all partial results we have
$`E_i(𝐫,t)E_j(𝐫,t)_0`$ $`=`$ $`T_1+T_2+T_3`$ (32)
$`=`$ $`\left({\displaystyle \frac{\pi }{a}}\right)^4{\displaystyle \frac{2}{3\pi }}\left[\left(\delta ^{}+\delta ^{}\right)_{ij}{\displaystyle \frac{1}{120}}+\delta _{ij}F(\xi )\right].`$
The function $`F\left(\xi \right)`$ is defined by
$$F\left(\xi \right):=\frac{1}{8}\frac{d^3}{d\xi ^3}\frac{1}{2}\mathrm{cot}\left(\xi \right),$$
(33)
and its expansion about $`\xi =0`$ is given by
$$F\left(\xi \right)\frac{3}{8}\xi ^4+\frac{1}{120}+O\left(\xi ^2\right).$$
(34)
Near $`\xi =\pi `$ (which corresponds to $`z=a`$) we make the replacement $`\xi \xi \pi `$. Notice that due to the behavior of $`F\left(\xi \right)`$ near $`\xi =0,a`$, strong divergences predominate in the behavior of the correlators near the plates.
By applying the exactly the same procedure we obtain the magnetic field correlators
$$B_i(𝐫,t)B_j(𝐫,t)_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }\left[\left(\delta ^{}+\delta ^{}\right)_{ij}\frac{1}{120}\delta _{ij}F(\xi )\right].$$
(35)
A direct evaluation also shows that the correlators $`<E_i(𝐫,t)B_j(𝐫,t)_0`$ are zero.
## 4 Correlators for Boyer’s setup
The other setup we are interested in is that one in which a perfectly conducting plate is placed at $`z=0`$ and perfectly permeable plate is placed at $`z=a`$. This setup was analyzed for the first time by Boyer in the contetxt of stochastic electrodynamic and it is the simplest case of a repulsive Casimir effect that can be found in the literature. The boundary conditions now are: *(a)* the tangential components $`E_x`$ and $`E_y`$ of the electric field as well as the normal component $`B_z`$ of the magnetic field must vanish on the surface of the plate at $`z=0`$; *(b)* the tangential components of $`B_x`$ e $`B_y`$ of the magnetic field as well as normal conponent $`E_z`$ of the electric field must vanish on the surface of the plate at $`z=a`$. These boundary conditions translated in terms of the components of the vector potential read
$$A_x(x,y,0,t)=0;A_y(x,y,0,t)=0;\frac{}{z}A_z(x,y,0,t)=0,$$
(36)
at $`z=0`$, and at $`z=a`$
$$\frac{}{x}A_x(x,y,a,t)=0;\frac{}{y}A_y(x,y,a,t)=0;A_z(x,y,a,t)=0.$$
(37)
The appropriate vector potential operator $`\widehat{𝐀}(𝐫,t)`$ is given by
$`\widehat{𝐀}(𝐫,t)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\left({\displaystyle \frac{\pi }{a}}\right)^{\frac{1}{2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle }{\displaystyle \frac{d^2\kappa }{\sqrt{\omega }}}\{\widehat{a}^{(1)}(\kappa ,n)\widehat{\kappa }\times \widehat{𝐳}\mathrm{sin}\left[(n+{\displaystyle \frac{1}{2}}){\displaystyle \frac{\pi z}{a}}\right]`$ (38)
$`+`$ $`\widehat{a}^{(2)}(\kappa ,n)[\widehat{\kappa }{\displaystyle \frac{i(n+\frac{1}{2})}{\omega a}}\mathrm{sin}\left[(n+{\displaystyle \frac{1}{2}}){\displaystyle \frac{\pi z}{a}}\right]\widehat{𝐳}{\displaystyle \frac{\kappa }{\omega }}\mathrm{cos}\left[(n+{\displaystyle \frac{1}{2}}){\displaystyle \frac{\pi z}{a}}\right]]\}`$
$`\times e^{i(\kappa \rho \omega t)}+\text{ }h.c.,`$
where as before $`\kappa =(k_x,k_y)`$ and $`\rho `$ is the position vector on the $`𝒳𝒴`$ plane. The normal frequencies are given
$$\omega (\kappa ,n)=\sqrt{\kappa ^2+\left(n+\frac{1}{2}\right)^2\frac{\pi ^2}{a^2}},$$
(39)
with $`k_x,k_y`$ and $`n1`$. Notice that contrary to the case of two conducting plates normalization does not require the we multiply the term corresponding to $`n=0`$ by $`1/2`$. The electric and magnetic field correlators for Boyer’s setup can be evaluated with the same technique employed before . In fact, it is not hard to convince ourselves that it is sufficient to perform the substitution $`nn+1/2`$ and follow the same steps as before to obtain
$$E_i(𝐫,t)E_j(𝐫,t)_0=T_1+T_2+T_3,$$
(40)
where, for example
$$T_1=\frac{\mathrm{\Gamma }\left(s\frac{3}{2}\right)}{\mathrm{\Gamma }\left(s\frac{1}{2}\right)}\left(\frac{\pi }{a}\right)^{32s}\frac{\delta _{ij}^{}}{2a}\left\{\underset{n=0}{\overset{\mathrm{}}{}}\left(n+\frac{1}{2}\right)^{32s}\frac{1}{2}\left[1\mathrm{cos}\left(2\left(n+1/2\right)\frac{\pi z}{a}\right)\right]\right\},$$
(41)
The terms $`T_2`$ e $`T_1`$ show a similar structure. The main difference with respect to the two conducting plate case is that now we have to deal with the Hurwitz zeta function $`\zeta _H(z,q)`$ which has a series representation given by
$$\zeta _H(z,q)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{\left(n+q\right)^z},$$
(42)
with $`\mathrm{}z>1`$, and $`q0,1,2,\mathrm{}.`$. In our case we must set
$$\zeta _H(2s3,\frac{1}{2})=\underset{n=0}{\overset{\mathrm{}}{}}\left(n+\frac{1}{2}\right)^{32s}.$$
(43)
It follows that in the limit $`s0`$ we have
$$T_1=\left(\frac{\pi }{a}\right)^4\frac{1}{3\pi }\frac{\delta _{ij}^{}}{2}\left\{\zeta _H(3,\frac{1}{2})+\underset{n=0}{\overset{\mathrm{}}{}}\left(n+\frac{1}{2}\right)^3\mathrm{cos}\left(2\left(n+\frac{1}{2}\right)\frac{\pi z}{a}\right)\right\}.$$
(44)
In the same way we obtain for $`T_2`$ e $`T_1`$ the results
$$T_2=\left(\frac{\pi }{a}\right)^4\frac{1}{\pi }\frac{\delta _{ij}^{}}{2}\left\{\zeta _H(3,\frac{1}{2})\underset{n=0}{\overset{\mathrm{}}{}}\left(n+\frac{1}{2}\right)^3\mathrm{cos}\left(2\left(n+\frac{1}{2}\right)\frac{\pi z}{a}\right)\right\}.$$
(45)
and
$$T_3=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }\frac{\delta _{ij}^{}}{2}\left\{\zeta _H(3,\frac{1}{2})+\underset{n=0}{\overset{\mathrm{}}{}}\left(n+\frac{1}{2}\right)^3\mathrm{cos}\left(2\left(n+\frac{1}{2}\right)\frac{\pi z}{a}\right)\right\}.$$
(46)
Notice that this time we do not have the divergent contribution corresponding to $`n=0`$ as in the case of the conducting plates. Adding the three terms we have
$`\widehat{E}_i(𝐫,t)\widehat{E}_j(𝐫,t)_0`$ $`=`$ $`\left({\displaystyle \frac{\pi }{a}}\right)^4{\displaystyle \frac{2}{3\pi }}\{(\delta _{ij}^{}+\delta _{ij}^{})\zeta _H(3,{\displaystyle \frac{1}{2}})`$ (47)
$`+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(n+{\displaystyle \frac{1}{2}})^3\mathrm{cos}\left(2(n+{\displaystyle \frac{1}{2}}){\displaystyle \frac{\pi z}{a}}\right)\}.`$
The numerical value $`\zeta _H(3,\frac{1}{2})`$ can be obtained from
$$\zeta _H(n,q)=\frac{B_{n+1}(q)}{n+1},$$
(48)
where $`n𝐍`$ and $`B_{n+1}(q)`$ is a Bernoulli polynomial defined by
$$B_n(x)=\underset{p=0}{\overset{n}{}}\frac{n!}{p!(np)!}B_px^{nk},$$
(49)
where $`B_p`$ is a Bernoulli number. The relevant polynomial here is:
$$B_4(x)=x^42x^3+x^2\frac{1}{30}.$$
(50)
With $`B_4(1/2)=(7/8)\times (1/30)`$, it follows that $`\zeta _H(3,\frac{1}{2})=(7/8)(1/120)`$. The sum can be regularized with the same technique employed before. In fact, we can define the function $`G(\xi )`$ by
$`G(\xi )`$ $`:`$ $`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(n+{\displaystyle \frac{1}{2}}\right)^3\mathrm{cos}\left[2\left(n+{\displaystyle \frac{1}{2}}\right)\xi \right]`$ (51)
$`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(2n+1\right)^3\mathrm{cos}\left[\left(2n+1\right)\xi \right],`$
where as before $`\xi :=z\pi /a`$. We can write
$$\left(2n+1\right)^3\mathrm{cos}\left[\left(2n+1\right)\xi \right]=\frac{d^3}{d\xi ^3}\mathrm{sin}\left[\left(2n+1\right)\xi \right]$$
(52)
and formally we have
$$G(\xi )=\frac{1}{8}\frac{d^3}{d\xi ^3}\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{sin}\left[\left(2n+1\right)\xi \right].$$
(53)
Writing $`\mathrm{sin}\left[\left(2n+1\right)\xi \right]`$ in terms of exponentials of imaginary argument and passing to the euclidean space we obtain after some simple manipulations
$$G\left(\xi \right)=\frac{1}{8}\frac{d^3}{d\xi ^3}\frac{1}{2\mathrm{sin}\left(\xi \right)}.$$
(54)
Collecting all partial results we finally obtain for $`\widehat{E}_i(𝐫,t)\widehat{E}_j(𝐫,t)_0`$ the result
$$\widehat{E}_i(𝐫,t)\widehat{E}_j(𝐫,t)_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }\left[\left(\frac{7}{8}\right)\frac{\left(\delta ^{}+\delta ^{}\right)_{ij}}{120}+\delta _{ij}G\left(\xi \right)\right].$$
(55)
Proceeding in the same way in the evaluation of $`\widehat{B}_i(𝐫,t)\widehat{B}_j(𝐫,t)_0`$ we obtain
$$\widehat{B}_i(𝐫,t)\widehat{B}_j(𝐫,t)_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }\left[\left(\frac{7}{8}\right)\frac{\left(\delta ^{}+\delta ^{}\right)_{ij}}{120}\delta _{ij}G\left(\xi \right)\right].$$
(56)
Observe that near $`\xi =0`$ the function $`G\left(\xi \right)`$ behaves as
$$G\left(\xi \right)=\frac{3}{8}\xi ^4\frac{7}{8}\frac{1}{120}+O\left(\xi ^2\right),$$
(57)
But near $`\xi =\pi `$ its behavior is slightly different
$$G\left(\xi \right)=\frac{3}{8}\left(\xi \pi \right)^4+\frac{7}{8}\frac{1}{120}+O\left[\left(\xi \pi \right)^2\right].$$
(58)
Again, a direct calculation shows that $`\widehat{E}_i(𝐫,t)\widehat{B}_j(𝐫,t)_0=0`$ for this case.
As before the divergent behavior of the correlators near the plates we are interested in is an effect of the distortions of the electromagnetic oscillations with respect to a situation where the plates are not present. This fact has received the attention of several authors, see for example .
## 5 The Casimir effect for conducting plates
In order to apply the above results to Casimir’s original setup we consider for convenience three parallel perfectly conducting plates perpendicular to the $`𝒪𝒵`$ axis at $`z=0`$, $`z=a`$ and $`z=\mathrm{}`$.
The quantum version of Maxwell tensor is
$$4\pi \widehat{T}_{ij}_0=\widehat{E}_i\widehat{E}_j_0\frac{1}{2}\delta _{ij}\widehat{𝐄}^2_0+\widehat{B}_i\widehat{B}_j_0\frac{1}{2}\delta _{ij}\widehat{𝐁}^2_0$$
(59)
Making use of the correlator given by (32) we obtain the following partial results
$$\widehat{E}_x^2(z,t)_0=\widehat{E}_y^2(z,t_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }[\frac{1}{120}+F\left(\xi \right)],$$
(60)
$$\widehat{E}_z^2(z,t_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }[\frac{1}{120}+F\left(\xi \right)],$$
(61)
also $`\widehat{E}_x(z,t)\widehat{E}_y(z,t)_0=\widehat{E}_x(z,t)\widehat{E}_z(z,t)_0=\widehat{E}_y(z,t)\widehat{E}_z(z,t)_0=0`$. In the same way, making use of (35) we obtain
$$\widehat{B}_x^2(z,t)_0=\widehat{B}_y^2(z,t_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }[\frac{1}{120}F\left(\xi \right)],$$
(62)
$$\widehat{B}_z^2(z,t_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }[\frac{1}{120}F\left(\xi \right)],$$
(63)
and also $`\widehat{B}_x(z,t)\widehat{B}_y(z,t)_0=\widehat{B}_x(z,t)\widehat{B}_z(z,t)_0=\widehat{B}_y(z,t)\widehat{B}_z(z,t)_0=0`$. The components of the quantum version of Maxwell tensor can be easily evaluated. For instance
$`8\pi \widehat{T}_{zz}(z,t)_0`$ $`=`$ $`\widehat{E}_z^2(z,t)_0\widehat{E}_x^2(z,t)_0\widehat{E}_y^2(z,t)_0`$ (64)
$`+\widehat{B}_z^2(z,t)_0\widehat{B}_x^2(z,t)_0\widehat{B}_y^2(z,t)_0.`$
performing the necessary substitutions we have
$$\widehat{T}_{zz}_0=\frac{\pi ^2}{240a^4}.$$
(65)
In the same way
$`\widehat{T}_{xx}_0`$ $`=`$ $`\widehat{T}_{yy}_0={\displaystyle \frac{1}{8\pi }}\left(\widehat{E}_z^2(z,t)_0+\widehat{B}_z^2(z,t)_0\right)`$ (66)
$`=`$ $`{\displaystyle \frac{\pi ^2}{720a^4}}`$
Notice how conveniently the divergent parts near tha plates cancel out yielding finite results.
Let us obtain now the Casimir force per unit area between the conducting plates. Consider Figure (1) and the plate at $`z=a`$. The Casimir force per unit area on this plate is
$`{\displaystyle \frac{F_z}{A}}`$ $`=`$ $`T_{zz}(za^L)+T_{zz}(za^R)`$ (67)
$`=`$ $`{\displaystyle \frac{\pi ^2}{240a^4}}+{\displaystyle \frac{\pi ^2}{240\left(\mathrm{}a\right)^4}},`$
where $`za^{L,R}`$ means that $`z`$ tends to $`a`$ from the left/right. Taking the limit $`\mathrm{}\mathrm{}`$ we obtain the expected result for the Casimir force per unit area
$$\frac{F_z}{A}=\frac{\pi ^2}{240a^4}.$$
(68)
The minus sign shows that the resulting pressure pushes towards the region between the plates. If simultaneously we take the limits $`\mathrm{},a\mathrm{}`$ keeping the distance $`\mathrm{}a`$ constant. The pressure changes its sign but it still pushes the plate at $`z=a`$ towards the one at $`z=\mathrm{}`$.
In order to calculate the renormalized symmetrical stress-energy tensor $`\widehat{\mathrm{\Theta }}^{\mu \nu }\left(z\right)_0^{ren}`$ we evaluate the energy density $`\rho \left(z\right)\widehat{\mathrm{\Theta }}^{00}\left(z\right)_0^{_{ren}}`$ in the region between the plates as well as the saptial components $`\widehat{\mathrm{\Theta }}^{xx}\left(z\right)_0^{ren}`$, $`\widehat{\mathrm{\Theta }}^{yy}\left(z\right)_0^{ren}`$, and $`\widehat{\mathrm{\Theta }}^{zz}\left(z\right)_0^{ren}`$. The energy density is given by
$$\rho (𝐫,t)=\frac{1}{8\pi }\left(\widehat{𝐄}^2(𝐫,t)_0+\widehat{𝐁}^2(𝐫,t)_0\right)$$
(69)
making use of the correlators given by (32) and (35) we obtain
$$\rho (a)=\frac{\pi ^2}{720a^4}.$$
(70)
This result is due to the fact that the divergent pieces in (32) and (35) cancel out yielding a finite result for the vacuum energy density. Recalling that $`\mathrm{\Theta }^{ij}(z)=T_{ij}(z)`$, see , with help of (2), (32) and (35) the remanescent components of the symmetrical stress-energy tensor are easily obtained. The final result is
$$\widehat{\mathrm{\Theta }}^{\mu \nu }(z)_0^{_{ren}}=\frac{\pi ^2}{720a^4}\text{diag}(1,1,1,3),$$
(71)
which is perfect agreement with Brown and Maclay’s results . Notice also that $`\widehat{\mathrm{\Theta }}_\mu ^\mu (z)_0^{_{ren}}=g_{\mu \nu }\widehat{\mathrm{\Theta }}^{\mu \nu }(z)_0^{_{ren}}=0`$, with $`g_{\mu \nu }=\text{diag}(1,1,1,1)`$.
## 6 The Casimir effect for one conducting plate and an infinitely permeable one
Let us consider now the setup proposed by Boyer which consists of a perfectly conducting plate placed perpendicularly to the $`𝒪𝒵`$ axis at $`z=0`$ and another infinitely permeable one parallel to the first placed at $`z=a`$. The boundary conditions on the conducting plate are as before $`E_x=E_y=0`$ and $`B_z=0`$, and for the infinitely permeable plate: $`B_x=B_y=0`$ and $`E_z=0`$. Making use of the correlator given by (55) the following partial results:
$$\widehat{E}_x^2(z,t)_0=\widehat{E}_y^2(z,t_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }[\frac{7}{8}\times \frac{1}{120}+G\left(\xi \right)],$$
(72)
$$\widehat{E}_z^2(z,t_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }[(\frac{7}{8})\times \frac{1}{120}+G\left(\xi \right)],$$
(73)
and $`\widehat{E}_x(z,t)\widehat{E}_y(z,t)_0=\widehat{E}_x(z,t)\widehat{E}_z(z,t)_0=\widehat{E}_y(z,t)\widehat{E}_z(z,t)_0=0`$. By the same token making use of the correlator given by (56) we obtain
$$\widehat{B}_x^2(z,t)_0=\widehat{B}_y^2(z,t_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }[\frac{7}{8}\times \frac{1}{120}G\left(\xi \right)],$$
(74)
$$\widehat{B}_z^2(z,t_0=\left(\frac{\pi }{a}\right)^4\frac{2}{3\pi }[(\frac{7}{8})\times \frac{1}{120}G\left(\xi \right)],$$
(75)
and also
$$\widehat{B}_x(z,t)\widehat{B}_y(z,t)_0=\widehat{B}_x(z,t)\widehat{B}_z(z,t)_0=\widehat{B}_y(z,t)\widehat{B}_z(z,t)_0=0$$
(76)
Proceeding as in the case of the conducting plates we obtain the following results for the componenets of the quantum version of Maxwell tensor
$$\widehat{T}_{xx}_0=\widehat{T}_{yy}_0=\frac{7}{8}\times \frac{\pi ^2}{720a^4},$$
(77)
and
$$\widehat{T}_{zz}_0=\left(\frac{7}{8}\right)\times \frac{\pi ^2}{240a^4}.$$
(78)
Notice that it is sufficient to multiply the results obtained for Casimir’s setup by the factor $`(7/8)`$ in order to obtain the results corresponding to Boyer’s setup.
In order to obtain the Casimir force per unit area for this setup it is convenient to place a third conducting plate at $`z=\mathrm{}`$. Then the Casimir force per unit area on the plate at $`z=a`$ will be given by
$`{\displaystyle \frac{F_z}{A}}`$ $`=`$ $`T_{zz}(za^L)+T_{zz}(za^R)`$ (79)
$`=`$ $`\left({\displaystyle \frac{7}{8}}\right)\times {\displaystyle \frac{\pi ^2}{240a^4}}+\left({\displaystyle \frac{7}{8}}\right)\times {\displaystyle \frac{\pi ^2}{240\left(\mathrm{}a\right)^4}}.`$
Taking the limit $`\mathrm{}\mathrm{}`$ we obtain a Casimir force which pushes the plate at $`z=a`$ towards the region $`z>a`$ given by
$$\frac{F_z}{A}=\frac{7}{8}\times \frac{\pi ^2}{240a^4}$$
(80)
This is the result obtained by Boyer for this setup using stochastic electrodynamic methods and it is one of the simplest example of a repulsive Casimir force.
In order to evaluate the symmetrical stress-energy tensor we first evaluate the Casimir energy density. Making use of (55) e (56) we have
$$\rho =\frac{7}{8}\times \frac{\pi ^2}{240a^4}.$$
(81)
As in the case of the conducting plates a simple calculation shows that the stress energy tensor for Boyer’s setup is given by
$$\widehat{\mathrm{\Theta }}^{\mu \nu }(z)_0^{_{ren}}=\frac{7}{8}\times \frac{\pi ^2}{720a^4}\text{diag}(1,1,1,3).$$
(82)
As before $`\widehat{\mathrm{\Theta }}_\mu ^\mu (z)_0^{_{ren}}=g_{\mu \nu }\widehat{\mathrm{\Theta }}^{\mu \nu }(z)_0^{_{ren}}=0`$.
## 7 The interaction between an atom and two material surfaces
In 1948, Casimir and Polder taking into account a suggestion made by experimentalists evaluated the interaction potential between two eletrical polarizable molecules separated by a distance $`r`$ including the effects due to the finiteness of the speed of propagation of the electromagnetic interaction, i.e.: of the retardment. Casimir and Polder showed that the retardment causes the interaction potential to change from a $`r^6`$ power law to a $`r^7`$ power law. In the same paper, Casimir and Polder also analyzed the interaction between an atom and a conducting wall e showed the interaction potetntial in this case varies according to a $`r^4`$, where now $`r`$ is the distance between the atom and the wall. Here we will show how it is possible to reobtain with the help of the correlators given by (32) and (35),the piece of Casimir and Polder’s result for the atom-wall interaction that depends on the distortion of the vaccum oscillations of the electromagnetic field caused by the presence of the wall.
From a classical point of view the induced eletrical polarization density $`𝐏`$ can be thought of as a function of the electric amd magnetic fields $`𝐄`$ and $`𝐁`$. In many cases only the dependence on eletric field is relevant. It can be shown that under conditions for which the effects of the retardment must be taken into account, i.e.: of the finiteness of the speed of light it is sufficient to consider the static eletrical polarizability $`\alpha \left(0\right)`$ only, see for instance reference . If the electric field changes by $`\delta 𝐄`$, the interaction between the polarizable body and the electric field will change according to $`\delta V=𝐏\left[𝐄\right]\delta 𝐄=\alpha \left(0\right)𝐄\delta 𝐄`$. Therefore, if the field changes from zero to a finite value $`𝐄`$, the interaction energy is $`V_E=\alpha \left(0\right)𝐄^2/2`$. In the quantum version of this interaction potential we must replace $`𝐄^2`$ by its vacuum expectation value, $`\widehat{𝐄}^2_0`$. The same arguments hold when we consider the magnetization $`𝐌`$. the interaction potential between a magnetically polarizable atom and the magnetic field is given by $`V_M=\beta \left(0\right)𝐁^2/2`$, where $`\beta \left(0\right)`$ is the static magnetic polarizability. The correlators are given by (32) and (35) which allow us to obtain in a straightforward way expressions for the interaction potential energy between an electrically or magnetically polarizable atom (placed between the plates) and the conducting plates.
Let us consider first an atom electrically polarizable placed at a distance $`z`$ from the conducting plate placed at $`z=0`$. The interaction potential between is given by
$$V_E\left(z\right)=\frac{1}{2}\alpha \left(0\right)\widehat{𝐄}^2\left(z\right)_0,$$
(83)
where $`\alpha \left(0\right)`$ is the static polarizability of the molecule. Making use of (32) and (35) we can evaluate $`\widehat{𝐄}^2\left(z\right)_0`$ and using the above equation we obtain
$$V_E\left(z\right)=\frac{\alpha \left(0\right)\pi ^3}{3a^4}\left[3F\left(\frac{\pi z}{a}\right)\frac{1}{120}\right].$$
(84)
Making use of (34) and taking the limit $`a\mathrm{}`$ we obtain the interaction potential between a polarizable atom and a conducting plate
$$V_E\left(z\right)=\frac{3\alpha \left(0\right)}{8\pi z^4}.$$
(85)
If a magnetically atom or molecule is placed between conducting plates the interaction potential will be given by
$$V_M\left(z\right)=+\frac{\beta \left(0\right)\pi ^3}{3a^4}\left[3F\left(\frac{\pi z}{a}\right)+\frac{1}{120}\right],$$
(86)
If the atom or molecule is simultaneously electric and magntically polarizable the interaction potential will be simply $`V\left(z\right)=V_E\left(z\right)+V_M\left(z\right)`$, that is
$$V\left(z\right)=\left(\alpha \left(0\right)\beta \left(0\right)\right)\frac{\pi ^3}{a^4}F\left(\frac{\pi z}{a}\right)+\left(\alpha \left(0\right)+\beta \left(0\right)\right)\frac{\pi ^3}{360a^4}.$$
(87)
The single conducting plate limit $`(a\mathrm{})`$ of (87) is easily obtained with the help of (34). The result is:
$$V\left(z\right)\frac{3}{8\pi z^4}\left(\alpha \left(0\right)\beta \left(0\right)\right),$$
(88)
which is in agreement with .
The polarizable atom or molecule can be also placed between a conducting plate at $`z=0`$ and a permeable one at $`z=a`$. In this case, making use of (55) e (56) a straightforward calculation leads to the following result
$$V\left(z\right)=\left(\alpha \left(0\right)\beta \left(0\right)\right)\frac{\pi ^3}{a^4}G\left(\frac{\pi z}{a}\right)+\left(\alpha \left(0\right)+\beta \left(0\right)\right)\left(\frac{7}{8}\right)\frac{\pi ^3}{360a^4}.$$
(89)
There are now two single plate limits to be considered. Near the conducting plate at $`z=0`$ the potential is given by (88), but near the perfectly permeable plate at $`z=a`$, the potential is repulsive and given by
$$V\left(z\right)+\frac{3}{8\pi \left(za\right)^4}\left(\alpha \left(0\right)\beta \left(0\right)\right),$$
(90)
where we made use of (58). Notice that here, as mentioned before, we dealt with a part of the interaction beteween an atom and two or one walls. The contribution of the interaction between the electric/magnetic dipole moment and its images was neglected. Therefore, the results refer only to the contribution of the quantum vacuum distorted by one or two walls to the total interaction potential. Keeping this in mind we can state that the Casimir-Polder interaction shows certain aspects of the quantum structure of the vacuum inbetween ans near the surfaces in question.
## 8 Conclusions
In this paper we have shown how to employ the equal time and space electromagnetic field correlators evaluated between parallel material surfaces to rederive results concernig the Casimir energy and pressure amd the symmetrical traceless stress energy tensor. This is a local alternative to the Green function technique. We have shown that for the cases we had in mind here finite results are obtained only when we consider what happens on both sides of the surface boundary. This consideration provided the mechanism by which precise cancellations occurred and finite results were obtained. This is in agreement with, for example, Ref. and should be considered as a concrete example of the behavior of quantized fields near and on boundary surfaces. As a byproduct of these calculations we have also analyzed the Casimir and Polder interaction between an atom and parallel material surfaces.
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# Detection of the old stellar component of the major Galactic bar
## 1 Introduction
The position of the Sun in the Galactic disc makes it difficult to study the large-scale structures in the inner Galaxy. The interstellar extinction means that visible wavelengths cannot penetrate more than a few kpc into the plane, so until large-area IR surface brightness maps and point source surveys became available little could be said on the distribution of stars in the inner regions. Even in the infrared there remains the problem that a single line of sight will contain sources from many Galactic components, which makes the results ambiguous. Even the Galaxy’s Hubble type, of importance in understanding how the Milky Way relates to other galaxies, is still uncertain.
Recently it has become accepted that the Galaxy is a barred spiral, however there is still some considerable argument over what the bar actually is. There is general agreement that the near end of the Galactic bar is in the first quadrant, but angles between the bar and the line of sight between the Sun and the Galactic Centre (GC) vary between about 10 and 75. The most popular angle based on studies of the distribution of stars is between 10 and 30 (Dwek 1995; Freudenreich 1998; López-Corredoira et al. 2000 ). However, these studies are all based on off-plane bulge sources and so although the feature is usually referred to as the bar, it could equally well be described as a triaxial bulge. Binney et al. (1991) studied the flow of gas in the inner few degrees and also showed that the gas-flow lines followed a bar-like potential with a position angle of 16.
Other authors have studied other regions in the plane and arrived at a larger angle for the bar. Sevenster et al. (1999) examined the kinematics of OH–IR stars and find a bar with an angle of 44. Peters (1975) on the basis of HI maps and Nakai (1992) using large scale CO maps have both suggested a bar near 45 . The most extreme bar angle yet suggested is 75 (Hammersley et al. 1994, hereafter Paper 1). This was based on the detection a large star formation region between $`l`$=27 and 21 and then showing that a stick-like bar with a half length of about 4 kpc would readily explain the form of the $`COBE`$ surface brightness maps in the Galactic plane. Other authors, however, have preferred to ascribe these features to a ring or blobs of star formation in spiral arms or trailing from the triaxial bulge (e.g. Freudenreich 1998).
If there really were a major bar, then as well as having major SFRs at both extremes it could also have a high concentration of older stars along its whole length, if the bar is a long-lived feature. Such a clustering would form a strong giant branch in a colour–magnitude diagram for any line of sight crossing the bar. The best way to search for an old population deep within the Galactic plane is to use near-IR colour–magnitude diagrams. Infrared wavelengths penetrate the interstellar dust far more easily than optical wavelengths. Furthermore, early K giants, which have by far the highest space density for the giants, have a very restricted range of absolute magnitudes. This means that they naturally form a dense clump on an infrared HR diagram. At the distance of interest these giants will have a $`K`$ magnitude of 12 to 14 making them simple to detect in a few seconds on even a small telescope.
In this paper we assume that the distance to the GC is 8 kpc, and that all positions quoted are on the Galactic plane unless otherwise stated. We will also assume absolute magnitudes for K2-3III of $`M_K`$=$``$1.65 $`M_H`$=$``$1.5 and $`M_J`$=$``$0.89 as in Table 2 of Wainscoat et al. (1992) and interstellar extinctions as given in Reike & Lebofsky (1985) $`A_K`$=0.112$`A_V`$, $`A_H`$=0.175$`A_V`$ and $`A_J`$=0.282$`A_V`$.
## 2 Observations
During 1999 June observations of a series of $`20\times 12`$ arcmin fields were observed along the Galactic plane between $`l`$=0 and $`l`$=37, using CAIN, the facility IR camera on the 1.5 m TCS (Observatorio del Teide, Tenerife). The seeing was typically 1<sup>′′</sup> and data were obtained only in photometric conditions. Star counts were then produced to $`J`$=17, $`H`$=16.5 and $`K_\mathrm{s}`$=15.2. Figure 1 show the $`H`$ vs. $`JH`$ colour–magnitude diagrams for the fields on the plane at $`l`$=32, 27, 20 and 10. The diagonal arrows shows the effect of 10 magnitudes of interstellar extinction in the visible.
## 3 Analysis
The main sequence stars have relatively low intrinsic luminosities, so only the closer ones are detected and there is little interstellar reddening. They form the triangular clump at $`JH`$ = 0.5. To the right of the main sequence there is a diagonal curving stripe running from top left to bottom right. This is formed by the K giants in the disc with increasing distance from the Sun. The middle solid line running diagonally downwards from left to right on the $`l`$=32 and 27 plots, shows the position of a K3III star on this diagram. This assumes that the extinction at any point follows a double exponential with a scale height of 50 pc and scale length of 3500 pc, as described in Wainscoat et al. (1992). A change in absolute magnitude displaces the line vertically whereas a change in the extinction changes the slope of the line in the diagram. Therefore, the fact that this line fits the stripe means that there cannot be significant errors in the assumed absolute magnitudes and extinction. The main sequence stars and the K-giant stripe can also be clearly seen in the $`l`$=27, 20 and 10 plots. A cluster of stars along the line of sight will form a giant branch extending upwards and slightly to the right (e.g. the bulge at $`l`$=10 at $`JH`$=2.8). On the $`l`$=27 plot two giant branches have been plotted for distances of 4.5kpc with an extinction of $`A_V`$=4.7mags and 6.6kpc with $`A_V`$=8.7mags.
In the middle of the $`l`$=27 stripe (m<sub>H</sub>=13.3,$`JH`$=1.3) there is a large clump of sources with a giant branch extending almost vertically to brighter magnitudes. This clump can also be seen at $`l`$=20 although somewhat more reddened ($`JH`$=2) and this extra extinction causes the giant branch to become less distinct. Extinction is patchy even on small angular scales and this leads to variations in the extinction to the individual sources detected. Therefore, if the average extinction rises the scatter in the magnitudes and colours will also increase. At $`l`$=10 the giant branch can be seen at $`JH`$=2.8, which is due to the bulge. The distance to the bulge plus extinction is too great to see the clump of K giants on the stripe.
In order to gain an idea of the number of sources in the clump, the K giants in the $`l`$=32 and 27 regions were isolated. The sources with a $`JH`$ within 0.3 mag of the predicted K3III line were extracted (the dashed lines in Fig. 1) and the $`K_\mathrm{s}`$ star counts for only these source plotted in Figure 2. $`K_\mathrm{s}`$ was used as it is less effected by extinction than $`H`$ or $`J`$. The sources used here have to be detected in all three bands and the limiting magnitude for this plot is actually set by $`J`$ and not $`K_\mathrm{s}`$, which is why the limit is around 13.5 and not 15.2. For magnitudes brighter than about m$`_{K_\mathrm{s}}`$=11.5 there is little difference between the two plots but then by m$`_{K_\mathrm{s}}`$=12.8 there are three times as many sources at $`l`$=27. The clump sources are extremely well defined; the FWHM is under 1 magnitude and the spread that is seen can be attributed principally to the luminosity function (LF) with a small part to the differences in extinction and errors in the photometry. This leaves very little spread due to the distance through the feature, which implies that the feature is compact along the line of sight. As the LF is not precisely known, it is impossible to determine the distance through the clump from these data alone, however it must be significantly less than 2 kpc as alone this would lead to a spread in magnitudes of 0.4 mag.
Figure 3 shows the $`K_\mathrm{s}`$ differential star counts for six regions after first removing the dwarfs and then subtracting the non-dwarf counts at $`l`$=32 from the resulting counts in each region. The regions have been offset by 30,000 counts to separate the plots and the zero for each region is shown by a dotted line. The dwarfs were removed using their position on the $`H`$ vs. $`JH`$ diagram in order to increase the contrast of any features in the inner Galaxy. There are no major clumps if the HR diagram at $`l`$=32 and the counts follow a smooth curve with magnitude (Fig. 2). By subtracting the non-dwarf $`l`$=32 counts, the steep rise in counts with increasing magnitude is reduced and the differences between the regions become far clearer. All the giants rather than just the K giants (as in Figure 2) were plotted as the $`J`$-band sensitivity is not sufficient to follow these stars to the distance of the bulge, also in many areas the K giant stripe becomes distorted by extinction.
The line of sight at $`l`$=30 runs into the dust lane inside the Scutum spiral arm whilst the $`l`$=32 region is tangential to the stars in the arm. Therefore, the $`l`$=30 region has more extinction than $`l`$=32 and hence the curve runs well below zero. In all of the other regions the plots are more or less zero between m$`{}_{K}{}^{}\mathrm{s}`$=10 and 12, which is surprising given the large difference in longitudes, however it does indicate that the regions can be directly compared. The peak at $`l`$=27, m$`_{K_\mathrm{s}}`$=12.8 shown in Figs. 2 and 3 can be clearly seen ($`l`$=27 was used since at $`l`$=25 there is a known dust lane which severely distorts the counts). This peak is also seen at $`l`$=20, 15, 10 and 5 but the magnitude of the peak goes fainter as the longitude decreases. At $`l`$=27, 20 and 15 the bulge should not provide any significant counts (e.g. Freudenreich 1998; López-Corredoira et al. 2000), although by $`l`$=10 and 5 the bulge is becoming important and this is seen in the increased number of sources and the increased width of the peak. Therefore, the clump at $`l`$=27 cannot belong to the bulge, but to a feature that runs into the bulge.
The $`JH`$ position of the clump (Fig. 1) is at 1.32 which indicates an extinction along the line of sight to the feature at $`l`$=27 of A<sub>V</sub>=6 $`\pm 1`$ magnitudes. The peak of the K giants in Figure 3 is at m$`_{K_\mathrm{s}}`$=12.8 and assuming an absolute magnitude of $`1.65`$ for the K giants and that $`A_{K_\mathrm{s}}`$=0.11$`A_V`$ gives the best distance from the Sun to the feature at $`l`$=27 of 5.7$`\pm `$0.7 kpc. The error is dominated by uncertainty in the absolute magnitude and extinction. The stated error was determined from the accuracy of the fit of the line to the peak of the K gaint stripe in Figure 1, which is 0.20 mags at the distance of the cluster. At $`l`$=20 the peak at $`K_\mathrm{s}`$ is about 0.5 mag fainter than at $`l`$=27 which, when extinction is taken into account, makes it about 0.5 kpc further away (Figs. 1 and 3).
As cross-check for the distance to the clump at $`l`$=27 Figure 4 shows the $`H`$ vs. $`JH`$ diagram for the region at $`l`$=2 $`b`$=$`2^{}`$. The bulge is clearly seen forming a giant branch extending upwards and slightly to the right near $`JH`$ =1. However, as the line of sight runs below the Galactic plane, the extinction is far less than in Figure 1. The K giants again form a very strong clump at m<sub>H</sub>=13.5, $`JH`$ =0.95. Therefore, assuming that these sources have a similar absolute magnitude as the clump at $`l`$=27, this implies that the clump at $`l`$=27 is at a distance of about 0.69 that of the bulge, i.e. 5.5 kpc. Conditions in the inner disc will be different from those in the bulge, hence this calculation can only be approximate but it is in agreement with the above distance.
Figure 5 gives part of the luminosity function for the clump sources in stars mag<sup>-1</sup>pc<sup>-3</sup> assuming that the distance through the clump is 500 pc, if the distance turns out to be larger then the density will be proportionally less. It was obtained by isolating the sources between the giant branches shown in the $`l`$=27 plot in Figure 1. The sources were then de-reddened assuming a distance of 5700 pc and the standard extinction values, the apparent magnitudes were then converted into absolute magnitudes. The process was repeated for $`l`$=32 and the result was then subtracted from that for $`l`$=27 to remove the disc and leave the the clump sources. No attempt has been made to correct for the distance through the clump or error in the photometry. For comparison the disc at the position of the clump would have a density between 5 and 10 times lower whereas the bulge at this location would have a density over 100 times lower (e.g. Wainscoat et al., 1992, López-Corredoira et al. ,2000). The bulge reaches this sort of density at only a few hundred pc from the GC. Given that there is also a large number of very luminous young sources at $`l`$=27 (Paper 1), then this must be one of the most luminous and densely populated parts of the Galaxy after the GC itself.
## 4 Discussion
The $`K_\mathrm{s}`$ star counts clearly show a feature where the number of old sources is similar between $`l`$=27 and $`l`$=5, but whose distance from the Sun increases with decreasing longitude. Paper 1 discusses the possible causes for the very high numbers of very luminous stars also seen at $`l`$=27 and $`l`$=21 and shows that these stars are related with the peaks in the COBE surface maps. Like the old star these young stars are seen at $`l`$=27 but not at $`l`$=32. The distance to the clump of K giants at $`l`$=27 is about 5.7 kpc, which is very close to the value determined in Garzón et al. (1997) for the luminous stars at $`l`$=27. This implies that there is an old population co-existing with, or very close to, a young population. Hence, the peaks in the $`COBE`$ surface brightness maps near $`l`$=27, are important for understanding the structures in the inner Galaxy and should not be dismissed as patchy star formation.
Paper 1 concludes that the most probable explanation for the young stars seen at $`l`$=27 is a bar, but that a ring cannot be ruled out. The above data discounts the hypothesis that the old sources seen here belong to a ring
* A ring, which is seen tangentially at $`l`$=27, should have far more sources at $`l`$=27 than at $`l`$=15, but this is not seen.
* The ring would be significantly closer at $`l`$=15 than at 27, hence the sources should appear brighter at $`l`$=15, not fainter as seen here.
A long lived bar, however, would naturally produce all of the features seen, in particular that the distance to the feature will increase with decreasing longitude if the near end of the bar is in the first quadrant.
The distance to the clump at $`l`$=27 fixes the position angle at 43$`{}_{}{}^{}\pm 7^{}`$ and a half-length of about 4 kpc. This is clearly at odds with Paper 1, which gives an angle of 75. Such a large bar angle would mean that the distance to the bar would be almost the same between $`l`$=27 and $`l`$=5 so the peak in Figure 4 would be more or less in the same position and not show the systematic change that is seen. The 75 position angle was determined principally on star formation region the far end of the bar corresponding to a peak at $`l`$=$``$22. A 43 bar angle would put the far end of the bar near $`l`$=$``$12 and indeed there is a peak in $`COBE`$ 2.2 $`\mu `$m surface brightness maps at this position (see Figs. 1 and 6 of Paper 1). It should be noted that the peak in the $`COBE`$ 2.2 $`\mu `$m surface brightness maps at $`l`$=$``$22 still needs to be explained and it is possible that there is a ring-like structure with tangential points at $`l`$=27 and $`l`$=$``$22, however a discussion of this is beyond the scope of this letter.
Many authors have examined the bulge off the plane and found that it is triaxial (e.g. Dwek 1995; Freudenreich 1998; López-Corredoira et al. 2000), however the position angle is small, typically about 15. This is about 30 different from the angle derived here for the bar. Although feature seen at $`l`$=27 cannot be the bulge as the latter does not give significant counts beyond about $`l`$=12, the bar apparently does runs into the bulge near $`l`$=12. Therefore, although geometrically the two features are distinct, currently it is not clear if the are dynamically linked and hence different aspects of the same phenomenon or where there are two independent bar-like features in the inner Galaxy, the triaxial bulge which dominates the inner few hundred pc and a major bar which extends to about 4 kpc.
## 5 Conclusions
There is a major old component in the $`l`$=27 line of sight at a distance of about 5.7 kpc from the Sun. The component is seen at $`l`$=20 and 15 and merges with the bulge inwards of this. The $`l`$=27 to 21 region is already known to have a very high density of young stars and hence all of the expected properties of a bar with a position angle of around 43 are present. The distribution of old stars is therefore very similar to that suggested for CO in Nakai (1992). The position angle is, however, significantly different from that of the triaxial bulge leading to the possibility that the Milky Way is a double-barred spiral galaxy.
## 6 Acknowledgments
The TCS is operated on the island of Tenerife by the Instituto de Astrofísica de Canarias at the Spanish Observatorio del Teide of the Instituto de Astrofísica de Canarias. We would like to thank the referee for some valuable suggestions.
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# 1 Introduction
## 1 Introduction
Yang-Mills instantons on noncommutative geometry have been subject to some interest nowadays in connection with the recent developments in superstrings theory and compactification of matrix model of M-theory . Noncommutative instantons are involved in the study of D(p-4)/Dp brane systems (p$`3`$) of superstrings; in particular in the ADHM construction of the D0/D4 system and in the determination of the vacuum solutions of the Higgs branch of supersymmetric gauge theories with eight supercharges .
In string theory, noncommutative geometry appears from the study of quantum properties of D-branes coupled to the closed string background fields $`g_{\mu \nu }`$ and $`B_{\mu \nu }`$. The boundary conditions of open strings of D-branes interpolate between Neumann and Dirichlet conditions and depend on the value of the NS-NS $`B_{\mu \nu }`$ field . For large values of B (or equivalently B constant if one uses the Seiberg-Witten limit taking the closed string metric $`g_{\mu \nu }`$ to zero) the action of the system is dominated by its boundary terms describing a boundary world-sheet conformal field theory. In this case the correlation functions of the one dimensional boundary string fields $`x^M(t)`$ satisfy nontrivial commutation relations leading to a noncommutative space-time. This noncommutative space has a canonical geometry extending the usual phase space geometry of quantum mechanics . The coordinates $`x^M(t)`$ do no longer commute and the usual product of functions is replaced by the star product of Moyal Bracket . Soon after this development several studies have been devoted to develop the quantum field theory on noncommutative spaces and many partial results have been obtained .
The aim of this paper is to contribute to these efforts by considering the problem of noncommutative instantons in harmonic space (NHS). Our main motivations for this are: (1) harmonic space is a natural space where the problem of solving the self-dual Yang-Mills constraints may be done in straightforward way due to the important role played by harmonic analyticity discovered by Galperin et al in the mid-eighties. This latter has known many successful applications as in the off-shell formulation of extended supersymmetric and supergravity theories , in hyperKahler metrics building and in the study of Yang-Mills self-dual solutions . (2) Standard noncommutative instantons analysis shows that the self-dual constraint eqs are non-linear and hence difficult to solve exactly . The known non trivial solutions are obtained by perturbative analysis around the ordinary geometry. But here we will develop a different method based on a noncommutative harmonic space leading to non perturbative explicit solutions. We will see that NHS method offers a powerful manner to go beyond the standard perturbative solutions.
The presentation of this paper is as follows. In section 2, we review the main lines of the problem of solving self-dual Yang-Mills constraints eqs in ordinary HS. In section 3, we study the Yang-Mills instantons on NHS. We first construct the NHS spaces and show the existence of two subspaces NHS($`\eta `$,0) and NHS(0,$`\theta `$) depending on the values of the deformation parameters $`\theta `$ and $`\eta `$. Then we establish the connection between these noncommutative geometries and superstrings theory in presence of background fields $`g_{\mu \nu }`$ and $`B_{\mu \nu }`$. Next we focus our attention on NHS($`\eta `$,0) and study the perturbative solutions of the self-dual constraint eqs while the derivation of the exact solutions of these constraints is given in section 4. Finally, section 5 will be devoted to the conclusion.
## 2 Self-Dual Yang-Mills Constraints
As promised in the introduction, we review here briefly the main lines of one of possible resolution ways of the problem of finding self-dual Yang-Mills (YM) solutions in ordinary $`𝐑^4`$. The method we will present is a powerful way exhibiting manifestly the SU(2) symmetry rotating the three Kahler structures of the hyperKahler moduli space of instantons. It allows in addition to reformulate the self-dual YM constraint eqs as integrability conditions for the existence of analytic homeomorphisms of patches of $`𝐂^2𝐑^4`$. This method permits also to construct explicit solutions of the self-dual constraints by working in harmonic space based on first realising geometrically the SU(2)$`𝐒^3`$ symmetry of the instanton hyperKahler moduli space and second using the concept of harmonic analyticity extending the usual complex analyticity where only one Kahler structure is manifestly exhibited.
To fix the ideas, we will start by describing the fundamentals of harmonic analyticity for YM instantons. Recall that such analyticity was successfully exploited in different occasions. In particular in the off-shell superspace formulation of supersymmetric field theories with eight supercharges , $`d=4`$ $`𝒩=2`$ and $`d=2`$ $`𝒩=4`$ (conformal) supergravity theories and in the hyperKahler metrics building .
### 2.1 Harmonic Analyticity
Since the analysis we will present hereafter is valid for generalized self-dual YM fields in ordinary $`𝐑^{4n}`$, $`n=1,2,\mathrm{}`$, we shall give the machinery for $`𝐑^{4n}`$ and consider the particular leading case whenever it is necessary. Roughly speaking $`𝐑^{4n}`$ is a real commutative Euclidean flat space whose local coordinates $`\{x^M,M=1,\mathrm{},4n\}`$ obey generally the following natural identities
$`\overline{x^M}`$ $`=x^M`$
$`x^Mx^Nx^Nx^M`$ $`=0`$
$`x^Mx^M`$ $`=\mathrm{\Lambda }_N^Mx^N,`$
which define the reality, commutativity and homogeneity conditions respectively. The $`\mathrm{\Lambda }_N^M`$’s are SO($`4n`$) Lorentz matrices in the $`\mathrm{𝟒}𝐧`$ vector representation. If $`𝐑^{4n}`$ is endowed by a complex structure, the natural coordinates are the usual complex variables $`z^\alpha `$ and of $`z_{\overline{\alpha }}`$ of $`𝐂^{2n}`$. They transform in the fundamental $`\mathrm{𝟐}𝐧`$ and $`\overline{\mathrm{𝟐}𝐧}`$ representations of U($`2n`$)=U(1)$`\times `$SU($`2n`$) $``$ SO($`4n`$) and related to the $`x^M`$’s as
$`z^\alpha `$ $`=`$ $`x^\alpha +ix^{\alpha +2n};`$
$`z_{\overline{\alpha }}`$ $`=`$ $`x^\alpha ix^{\alpha +2n},\alpha =1,\mathrm{},2n.`$ (2.2)
Reality, commutativity and homogeneity of $`𝐑^{4n}`$ then become in $`𝐂^{2n}`$
$`[z^\alpha ,z^\beta ]=[z_{\overline{\alpha }},z_{\overline{\beta }}]=[z^\alpha ,z_{\overline{\beta }}]=0`$
$`z^\alpha =u_{.\beta }^\alpha z^\beta ,z_{\overline{\alpha }}^{}=z_{\overline{\beta }}u_{.\overline{\alpha }}^{\overline{\beta }}`$ (2.3)
where $`u_{.\overline{\beta }}^{\overline{\alpha }}\mathrm{𝟐}𝐧`$ and $`u_{.\overline{\alpha }}^{\overline{\beta }}\overline{\mathrm{𝟐}𝐧}`$ with the property $`z_{\overline{\alpha }}=(z^\alpha )^+`$. Though this complex frame exhibits manifestly one complex structure, one can go ahead and study the problem of solving the YM self-duality constraint eqs on $`𝐂^2`$. One recalls for instance the ADHM construction of YM instantons which find actually many applications in D-brane physics. Using this complex frame one can do even more by considering the problem of YM fields on $`𝐂^{2n},n1`$ and study the solving of the generalized self-dual constraint eqs. Such a kind of problem was considered in many occasions before as in YM and gravitational theories on $`𝐂^p`$ in connection with the Yang-Lee theory and (hyper)Kahler geometry respectively . In the present study, we shall not use the complex frame as given in (2.1,2.1). What we will do instead is to use the SU(2) symmetry factor of the SU(2)$`\times `$ SP($`n`$) homogeneity group of $`n`$ dimensional hyperKahler manifolds to introduce a new local frame of $`𝐂^{2n}`$ where the three complex structures are manifestly exhibited. Breaking the SO($`4n`$) Lorentz group of $`𝐑^{4n}`$ down to SU(2)$`\times `$ SP($`n`$) by reindexing the $`x^M`$ variables as $`x^{i\alpha }`$, where now the double index $`(i,\alpha )`$ transform in the $`(\mathrm{𝟐},\mathrm{𝟐}𝐧)`$ representation of SU(2)$`\times `$ SP($`n`$), eqs (2.1) read then as
$`\overline{x}^{i\alpha }`$ $`=\mathrm{\Omega }_{\alpha \beta }\epsilon _{ij}x^{j\beta },`$
$`[x^{i\alpha },x^{j\beta }]`$ $`=0`$
$`x^{i\alpha }x^{i\alpha ^{}}`$ $`=u_j^iC_\beta ^\alpha x^{j\beta },`$
where $`u_j^i`$ and $`C_\beta ^\alpha `$ are respectively $`2\times 2`$ and $`2n\times 2n`$ matrices in the fundamental representations of SU(2) and SP($`n`$), while $`\mathrm{\Omega }`$ and $`\epsilon `$ satisfy the following properties $`\mathrm{\Omega }^{\alpha \gamma }\mathrm{\Omega }_{\gamma \delta }=\delta _\delta ^\gamma ,\mathrm{\Omega }^{\alpha \beta }=\mathrm{\Omega }^{\beta \alpha },\mathrm{\Omega }^{\alpha ,\alpha +n}=1;1\alpha n`$ and $`\epsilon ^{ik}\epsilon _{kj}=\delta _j^i,\epsilon ^{ij}=\epsilon ^{ji},\epsilon ^{12}=\epsilon _{21}=1`$.
To write down the generalized self-dual YM constraint eqs using the local coordinates system $`\{x^{i\alpha }\}`$, it is enough to consider the gauge covariant derivatives $`𝒟_{i\alpha }`$ and the gauge curvatures $`F_{i\alpha j\beta }`$. Like for YM theory on $`𝐑^{4n}`$ with gauge group G, we have
$`𝒟_{i\alpha }`$ $`=`$ $`{\displaystyle \frac{}{x^{i\alpha }}}+A_{i\alpha }`$ (2.5)
$`=`$ $`_{i\alpha }+A_{i\alpha },`$
where $`A_{i\alpha }=A_{i\alpha }(x^{j\beta })`$ are the gauge potential component fields valued in the Lie algebra $`𝐠`$ of the gauge group G. The field strengths $`F_{i\alpha j\beta }`$ are given by the commutators of the $`𝒟_{i\alpha }`$’s
$$[𝒟_{i\alpha },𝒟_{j\beta }]=iF_{i\alpha j\beta }.$$
(2.6)
Taking the tensor product of the fundamental representation of $`SU(2)\times SP(n)`$, one can put the above eqs into a convenient form by decomposing $`F_{i\alpha j\beta }`$ curvatures as shown here below
$$[𝒟_{i\alpha },𝒟_{j\beta }]=i\epsilon _{ij}F_{(\alpha \beta )}+iF_{(ij)[\alpha \beta ]},$$
(2.7)
where ( ) and \[ \] denote complete symmetrisation and antisymmetrisation respectively. Having given the curvatures $`F_{i\alpha j\beta }`$, we turn now to introduce the generalized self-duality constraint eqs of YM fields on $`4n`$ dimensional space with SU(2)$`\times `$ SP($`n`$) homogeneity symmetry. They are defined as
$$2iF_{(ij)[\alpha \beta ]}=([𝒟_{i\alpha },𝒟_{j\beta }]+[𝒟_{j\alpha },𝒟_{i\beta }])=0,$$
(2.8)
or equivalently
$$[𝒟_{i\alpha },𝒟_{j\beta }]=i\epsilon _{ij}F_{(\alpha \beta )}.$$
(2.9)
Note that for $`n=1`$, the homogeneity group of $`𝐑^4`$ is SU(2)$`\times `$ SP(1)$``$ SO(4). In this case, $`F_{(ij)[\alpha \beta ]}=\mathrm{\Omega }_{\alpha \beta }F_{(ij)}`$ and so eq (2.7) becomes
$$[𝒟_{i\alpha },𝒟_{j\beta }]=i\epsilon _{(ij)}F_{(\alpha \beta )}+i\mathrm{\Omega }_{\alpha \beta }F_{(ij)}.$$
(2.10)
In this special situation one may impose not only the self-dual constraint eqs, $`F_{(ij)}=0`$, but also anti self-dual ones associated with the condition $`F_{(\alpha \beta )}=0`$.
What we want to do now is to show that in some special situations one may solve explicitly the generalized self-dual constraints (2.8,2.9) by using harmonic analyticity. The method may be summarized as follows: First, instead of using the local coordinates systems $`\{x^{i\alpha }\}`$ parameterizing the $`4n`$ dimensional homogeneous space $`𝒫/SU(2)\times SP(n)`$, where $`𝒫`$ stands for the Poincaré group with SU(2) $`\times `$ SP($`n`$) as its Lorentz subgroup, we rather use the harmonic space $`𝒫`$/SP($`n`$) parametrized by $`\{x^{\pm \alpha },u_i^\pm \}`$ .
On this new space we have the following:
(i) the $`u_i^\pm `$ variables parameterize the unit $`𝐒^3`$ sphere. In terms of these variables, the SU(2) Lorentz sub-symmetry of the $`4n`$ space is generated by the differential operators
$`D^{++}`$ $`=u^+/u^{}`$
$`D^{}`$ $`=u^{}/u^+`$
$`D^0`$ $`=u^+/u^+u^{}/u^{},`$
which obey the familiar su(2) algebra
$`[D^{++},D^{}]`$ $`=D^0`$
$`[D^0,D^{++}]`$ $`=2D^{++}`$
$`[D^0,D^{}]`$ $`=2D^0`$
(ii) the $`x_{}^{+}{}_{}{}^{\alpha }`$ and $`x_{}^{}{}_{}{}^{\alpha }`$ bosonic coordinates still carry an su(2) isospin $`\frac{1}{2}`$ representation since
$`[D^{++},x_{}^{+}{}_{}{}^{\alpha }]`$ $`=0`$
$`[D^{++},x_{}^{}{}_{}{}^{\alpha }]`$ $`=x_{}^{+}{}_{}{}^{\alpha }`$
$`[D^0,x_{}^{\pm }{}_{}{}^{\alpha }]`$ $`=\pm x_{}^{\pm }{}_{}{}^{\alpha }.`$
Eqs (2.1) mean that $`x_{}^{+}{}_{}{}^{\alpha }=u_i^+x^{i\alpha }`$ and $`x_{}^{}{}_{}{}^{\alpha }=u_i^{}x^{i\alpha }`$. Note that, as far as the $`x^\pm `$’s are concerned, the usual reality condition $`\stackrel{~}{x^{i\alpha }}=\epsilon _{ij}\mathrm{\Omega }_{\alpha \beta }x^{j\alpha }`$ is replaced in HS by
$`\stackrel{~}{x_{}^{+}{}_{}{}^{\alpha }}`$ $`=`$ $`\mathrm{\Omega }_{\alpha \beta }x_{}^{+}{}_{}{}^{\beta }`$
$`\stackrel{~}{x_{}^{}{}_{}{}^{\alpha }}`$ $`=`$ $`\mathrm{\Omega }_{\alpha \beta }x_{}^{}{}_{}{}^{\beta }.`$ (2.14)
This conjugaison has been introduced first in and turns out to be the natural conjugation in HS.
(iii) The gauge covariant derivatives in HS language are defined as
$`𝒟_\alpha ^+=/x_{}^{}{}_{}{}^{\alpha }+iA_\alpha ^+=_\alpha ^++iA_\alpha ^+`$
$`𝒟_\alpha ^{}=/x_{}^{+}{}_{}{}^{\alpha }+iA_\alpha ^{}=_\alpha ^{}+iA_\alpha ^{}.`$ (2.15)
They obey the following constraint eqs
$`[D^{++},𝒟_\alpha ^+]`$ $`=0`$
$`[D^{},𝒟_\alpha ^+]`$ $`=𝒟_\alpha ^{}`$
$`[D^0,𝒟_\alpha ^\pm ]`$ $`=\pm 𝒟^\pm `$
which ensure that $`𝒟_\alpha ^\pm `$ are SU(2) doublets; i.e. $`𝒟_\alpha ^\pm =u_i^\pm 𝒟_\alpha ^i`$.
(iv) In HS, the gauge curvatures read as
$`[𝒟_\alpha ^+,𝒟_\beta ^+]`$ $`=`$ $`iF_{[\alpha \beta ]}^{++}`$
$`[𝒟_\alpha ^+,𝒟_\beta ^{}]`$ $`=`$ $`iF_{(\alpha \beta )}+iF_{[\alpha \beta ]}^+`$ (2.17)
They verify the constraint eqs
$`[D_{}^{++}{}_{}{}^{2},F_{[\alpha \beta ]}^+]`$ $`=[𝒟^{++},F_{[\alpha \beta ]}^{++}]=0`$ (2.18)
$`[D^{++},F_{(\alpha \beta )}]`$ $`=0.`$
Eqs (2.1) mean amongst others that the symmetric part $`F_{(\alpha \beta )}`$ of the curvature is a SU(2) singlet as it does not depend on the $`u^\pm `$’s while the antisymmetric part is a SU(2) triplet, i.e.
$$F_{[\alpha \beta ]}^{++}=u_{(i}^+u_{j)}^+F_{[\alpha \beta ]}^{(ij)}.$$
(2.19)
The second step consists of writing the generalized self-duality conditions (2.8,2.9) in the HS as
$$[𝒟_\alpha ^+,𝒟_\beta ^+]=0.$$
(2.20)
These constraints, which take a simple form in HS, can be treated as integrability conditions for the existence of an analytic harmonic subspace (AHSS) of the HS space. On this subspace, the solving of $`F_{[\alpha \beta ]}^{++}=0`$ may be easily worked out. The key idea of the AHSS resolution method is to say that under a gauge transformation<sup>1</sup><sup>1</sup>1This transformation is given in , one may set $`A_\alpha ^+=0`$ so that the gauge covariant derivatives $`D_\alpha ^+`$ reduces to a flat space derivative $`^+=/x_{}^{}{}_{}{}^{\alpha }`$. The price one should pay for this operation is that the su(2) harmonic derivatives $`D^{++}`$, $`D^{}`$ and $`D^0`$ eqs (2.1) acquire now connections as shown here below
$`D^{++}𝒟^{++}`$ $`=D^{++}+iV^{++}`$
$`D^{}𝒟^{}`$ $`=D^{}+iV^{}`$
$`D^0𝒟^0`$ $`=D^0,`$
with the requirement
$`[𝒟^{++},𝒟^{}]`$ $`=D^0`$ (2.22)
$`[D^0,𝒟^{\pm \pm }]`$ $`=\pm 2𝒟^{\pm \pm },`$
or equivalently
$`[D^0,V^{++}]`$ $`=2V^{++},`$
$`[D^0,V^{}]`$ $`=2V^{}.`$
$`D^{++}V^{}D^{}V^{++}+i[V^{++},V^{}]`$ $`=0,`$
Before going ahead, let us make some remarks regarding this tricky approach. $`V^{++}`$ and $`V^{}`$ are harmonic gauge connections; they are not independent as they are related through eqs (2.1). Up on solving them, one may usually express $`V^{}`$ as a function of $`V^{++}`$. We will show later on how this is done in practice. For the moment let us note that $`V^{++}`$ turns out to be the basic object in solving the generalized self-duality constraint eqs. This property may be easily seen by rewriting the full set of the generalized self-duality constraints in HS. We have:
$`[D_\alpha ^+,D_\beta ^+]`$ $`=0`$
$`[D_\alpha ^+,D_\beta ^{}]`$ $`=iF_{(\alpha \beta )}`$
$`[D^0,D^\pm ]`$ $`=D_\alpha ^\pm ,`$
and
$`[𝒟^{++},D_\alpha ^+]`$ $`=0`$ (2.25)
$`[𝒟^{},D_\alpha ^+]`$ $`=D_\alpha ^{}`$
$`[𝒟^{++},𝒟^{}]`$ $`=D^0`$
$`[𝒟^0,𝒟^{\pm \pm }]`$ $`=\pm 2𝒟^{\pm \pm }.`$
To derive the solutions of these constraints eqs, let us rewrite them in a more convenient way. Using harmonic analyticity which allows to take $`𝒟_\alpha ^+`$ as $`/x_{}^{}{}_{}{}^{\alpha }`$; i.e. $`A_\alpha ^+=0`$, and replacing $`𝒟^{++}`$, $`𝒟^{}`$ and $`𝒟^0`$ by their expressions (2.1), the generalized self-duality constraint eqs may be represented as
$`[D^0,V^{++}]`$ $`=2V^{++}(a)`$
$`[_\alpha ^+,V^{++}]`$ $`=0(b)`$
$`D^{++}V^{}D^{}V^{++}+i[V^{++},V^{}]`$ $`=0(c)`$
$`[_\alpha ^+,V^{}]`$ $`=A_\alpha ^{}(d)`$
$`[_\alpha ^+,A_\beta ^{}]`$ $`=F_{(\alpha \beta )}.(e)`$
These relations define the form of the self-duality YM constraint eqs we are looking for.
### 2.2 Solving The Constraints
To solve the previous eqs, we shall proceed step by step by working out the solution of each relation of the system (2.1). These details will be useful in the study of the instantons in NHS. The first relation (2.1.a) tells us that $`V^{++}`$ is a harmonic function on the $`𝐒^2SU(2)/U(1)`$ sphere and so it may be expanded in a harmonic series in term of the $`u^\pm `$’s as:
$`V^{++}(x^+,x^{},u^\pm )`$ $`=`$ $`V^{++}(x^+,x^{})+u_{(i}^+u_{j)}^{}V_{}^{++}{}_{}{}^{(ij)}(x^+,x^{})`$ (2.27)
$`+u_{(i_1}^+u_{i_2}^+u_{j_1}^{}u_{j_2)}^{}V_{}^{++}{}_{}{}^{(i_1i_2j_1j_2)}(x^+,x^{})+\mathrm{},`$
where the fields $`V_{}^{++}{}_{}{}^{(i_1i_2\mathrm{}i_n,j_1j_2\mathrm{}j_n)}(x^+,x^{})`$ are coefficients of the harmonic expansion on $`𝐒^2`$. It turns out that except the leading term $`V^{++}(x^+,x^{})`$ where no explicit dependence on the $`u^\pm `$’s appear, all the extra terms of (2.27) may be ignored due to the existence of U(1) symmetry see (4). Forgetting about these terms, and putting (2.27) in (2.1.b), one sees that the potential $`V^{++}`$ should be an analytic function on $`𝐒^2`$ depending on the $`x^+`$ variable only. This means that the general solution of (2.1) for a gauge group G of generators $`\{T^a\}`$ should read as
$`(V^{++})^a`$ $`=`$ $`\mathrm{tr}(T^aV^{++})`$ (2.28)
$`=`$ $`_{\alpha \beta }^ax_{}^{+}{}_{}{}^{\alpha }x_{}^{+}{}_{}{}^{\beta },`$
where $`_{\alpha \beta }^a`$ are constant coefficients scaling as the inverse of the square length dimension. Note that $`_{\alpha \beta }^a`$ carry two SP($`n`$) indices $`\alpha ,\beta `$. Commutativity of $`x_{}^{+}{}_{}{}^{\alpha }`$ and $`x_{}^{+}{}_{}{}^{\beta }`$ requires $`_{\alpha \beta }^a=_{\beta \alpha }^a`$ and so transform under the adjoint representation of the SP($`n`$) Lorentz subgroup. Note also that in the case where the gauge group G is itself a SP($`n`$) symmetry, one may take as a solution for the SP($`n`$) potentials $`V_{}^{++}{}_{\alpha \beta }{}^{}`$ the following remarkable one
$$(V^{++})_{\alpha \beta }=\frac{i}{\rho ^2}\mathrm{\Omega }_{\alpha \gamma }\mathrm{\Omega }_{\beta \delta }x_{}^{+}{}_{}{}^{\gamma }x_{}^{+}{}_{}{}^{\delta }$$
(2.29)
where $`\rho `$ is a dimensionful parameter interpreted in $`𝐑^4`$ as the size of the SU(2) instanton. Putting this solution back into the constraint (2.1.c), one gets after some straightforward algebra the following expression for the SP($`n`$) gauge potentials:
$$(V^{})_{\alpha \beta }=\frac{i}{(x^2+\rho ^2)}\mathrm{\Omega }_{\alpha \gamma }\mathrm{\Omega }_{\beta \delta }x_{}^{}{}_{}{}^{\gamma }x_{}^{}{}_{}{}^{\delta },$$
(2.30)
where we have used the convention notation $`x^2=\mathrm{\Omega }_{\alpha \beta }x_{}^{+}{}_{}{}^{\alpha }x_{}^{}{}_{}{}^{\beta }=x_{}^{+}{}_{}{}^{\alpha }x_{}^{}{}_{\alpha }{}^{}`$. From (2.30), one gets the gauge potential $`A_\alpha ^{}`$ by help of (2.1.d). It reads as:
$$(A_\nu ^{})_{\alpha \beta }=\frac{i}{(x^2+\rho ^2)}\left[\mathrm{\Omega }_{\nu (\alpha }\mathrm{\Omega }_{\beta )\gamma }x_{}^{}{}_{}{}^{\gamma }+\frac{1}{(x^2+\rho ^2)}\mathrm{\Omega }_{\nu \sigma }\mathrm{\Omega }_{\alpha \gamma }\mathrm{\Omega }_{\beta \delta }x_{}^{+}{}_{}{}^{\sigma }x_{}^{}{}_{}{}^{\gamma }x_{}^{}{}_{}{}^{\delta }\right]$$
(2.31)
The curvatures $`(F_{\mu \nu })_{\alpha \beta }`$ are immediately obtained by solving (2.1.e) using the relation of $`A_\alpha ^{}`$ given here above. We find:
$`(F_{\mu \nu })_{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{(x^2+\rho ^2)}}\left[\mathrm{\Omega }_{\mu (\alpha }\mathrm{\Omega }_{\beta )\nu }+{\displaystyle \frac{1}{(x^2+\rho ^2)}}\mathrm{\Omega }_{\mu (\alpha }\mathrm{\Omega }_{\beta )\gamma }\mathrm{\Omega }_{\nu \sigma }x_{}^{+}{}_{}{}^{\sigma }x_{}^{}{}_{}{}^{\gamma }\right]`$ (2.32)
$`{\displaystyle \frac{i}{(x^2+\rho ^2)^2}}\mathrm{\Omega }_{\mu \rho }\mathrm{\Omega }_{\nu (\alpha }\mathrm{\Omega }_{\beta )\gamma }x_{}^{+}{}_{}{}^{\rho }x_{}^{}{}_{}{}^{\gamma }`$
$`{\displaystyle \frac{2i}{(x^2+\rho ^2)^3}}\mathrm{\Omega }_{\mu \rho }\mathrm{\Omega }_{\nu \sigma }\mathrm{\Omega }_{\alpha \gamma }\mathrm{\Omega }_{\beta \delta }x_{}^{+}{}_{}{}^{\rho }x_{}^{+}{}_{}{}^{\sigma }x_{}^{}{}_{}{}^{\gamma }x_{}^{}{}_{}{}^{\delta }`$
Finally the Lagrangian density of the instanton is:
$$\mathrm{tr}\left[(F_{\mu \nu })^2\right]=\frac{\rho ^4}{(x^2+\rho ^2)^4}.$$
(2.33)
In the limit $`\rho 0`$, eq (2.33) has a singularity near the origin and so the density tr$`(F_{\mu \nu })^2`$ behaves as a Dirac delta function.
## 3 Yang-Mills Instantons on Noncommutative $`𝐑_\theta ^4`$
We start this section by considering the extension of the harmonic space analyticity idea in order to study the problem of Y-M instantons on noncommutative $`𝐑_\theta ^4`$. To do so, we have a priori different choices of non commutative local coordinates. For example one can use directly the usual real coordinates system $`\{x^M\}`$, where the $`x^M`$’s transforming in the 4 vector representation of SO(4) satisfy:
$$[x^M,x^N]=i\theta ^{MN}.$$
(3.1)
Then extends the standard analysis of YM instantons by incorporating this noncommutativity feature. An other way is to break the SO(4) Lorentz group down to U(2)$`=`$U(1)$`\times `$SU(2) and use the complex coordinates frame $`\{z^\mu =x^\mu +ix^{\mu +2};iz^{\overline{\mu }}=x^\mu ix^{\mu +2};\mu =1,2\}`$ satisfying the following commutation relations
$`[z^\mu ,z^\nu ]`$ $`=`$ $`\theta ^{\mu \nu }`$
$`[z^{\overline{\mu }},z^{\overline{\nu }}]`$ $`=`$ $`\theta ^{\overline{\mu }\overline{\nu }}`$
$`[z^\mu ,z^{\overline{\nu }}]`$ $`=`$ $`\theta ^{\mu \overline{\nu }}`$
$`[z^{\overline{\mu }},z^\nu ]`$ $`=`$ $`\theta ^{\overline{\mu }\nu }`$ (3.2)
This coordinate system, which is useful whenever one can impose complex analyticity, is nowadays intensively used in the noncommutative ADHM formulation of instantons . But we will use neither the first method nor the second. What we will do instead is to develop a non commutative harmonic space analysis by extending the study of subsection 2.2. Thus taking the local coordinates of $`𝐑^4`$ as $`\{x^{i\alpha };i=1,2;\alpha =1,2\}`$. The space of noncommutativity reads as
$$[x^{i\alpha },x^{j\beta }]=\mathrm{\Omega }^{\alpha \beta }\eta ^{(ij)}+\epsilon ^{ij}\theta ^{(\alpha \beta )},$$
(3.3)
where $`\eta ^{(ij)}`$ and $`\theta ^{(\alpha \beta )}`$ transform respectively under the $`(1,0)`$ and $`(0,1)`$ representations of the SU(2)$`\times `$SP(1) Lorentz group. Now introducing the harmonic variables $`u_i^\pm `$, the noncommutative harmonic space is defined as
$`[x_{}^{+}{}_{}{}^{\alpha },x_{}^{+}{}_{}{}^{\beta }]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta ^{++}`$
$`[x_{}^{}{}_{}{}^{\alpha },x_{}^{}{}_{}{}^{\beta }]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta ^{}`$
$`[x_{}^{+}{}_{}{}^{\alpha },x_{}^{}{}_{}{}^{\beta }]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta i\theta ^{(\alpha \beta )}`$
$`[x_{}^{}{}_{}{}^{\alpha },x_{}^{+}{}_{}{}^{\beta }]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta +i\theta ^{(\alpha \beta )},`$ (3.4)
where we have used
$`\eta ^{++}`$ $`=u_{(i}^+u_{j)}^+\theta ^{(ij)}`$
$`\eta ^{}`$ $`=u_{(i}^{}u_{j)}^{}\theta ^{(ij)}`$
$`\eta `$ $`=u_{(i}^+u_{j)}^{}\theta ^{(ij)}.`$
Observe that according to the values of $`\eta ^{(ij)}`$ and $`\theta ^{(\alpha \beta )}`$ which together carry the six degrees of freedom of the six-dimensional antisymmetric representation of SO(4), one may distinguish four kinds of harmonic subspaces: (1) the ordinary one considered in section 2 corresponding to $`\theta ^{(ij)}=0`$; $`\eta ^{(ij)}=0`$ and three noncommutative ones given by (2) $`\eta ^{(ij)}0`$; $`\theta ^{(\alpha \beta )}=0`$, (3) $`\eta ^{(ij)}=0`$; $`\theta ^{(ij)}0`$ and (4) $`\eta ^{(ij)}0`$; $`\theta ^{(\alpha \beta )}0`$. We shall refer to these four kinds of harmonic spaces as HS($`\eta ,\theta `$) which roughly speaking may be viewed as a two deformation parameters of the standard harmonic space. Note moreover that these spaces are intimately related with the problem of finding (anti) self-dual YM instantons on noncommutative geometry. In this regards it is not difficult to see that the noncommutative harmonic subspaces HS($`\eta `$,0) and HS(0,$`\theta `$) are respectively associated with self-dual and anti self-dual YM instantons of the noncommutative geometry. The point is to use the Seiberg-Witten realisation according to which in the $`\alpha ^{}`$ zero slope limit the parameters $`\eta ^{(ij)}`$ and $`\theta ^{(\alpha \beta )}`$ are proportional to the inverse of the NS-NS antisymmetric $`B_{MN}`$ field leaving on an Euclidean D3-Brane world volume. Put differently, if we decompose $`B_{i\alpha j\beta }`$ as $`\epsilon _{ij}B_{(\alpha \beta )}+\mathrm{\Omega }_{\alpha \beta }B_{(ij)}`$; we have
$$\eta ^{(ij)}B_{(ij)}/B^2;\theta ^{(\alpha \beta )}B_{(\alpha \beta )}/B^2.$$
(3.6)
These eqs show clearly that $`\eta ^{(ij)}`$ is related to the self-dual part of the $`B`$ field while $`\theta ^{(\alpha \beta )}`$ is given by the anti self-dual one. In what follows we shall consider the example of noncommutative HS($`\eta `$,0) and study the problem of finding noncommutative YM instantons by using the idea of harmonic analyticity developed previously. Similar analysis may be done for the anti self-dual constraints by using NHS(0, $`\theta `$).
### 3.1 Noncommutative YM instantons in NHS($`\eta `$,0)
The commutation relations of the local coordinates $`\{x_{}^{+}{}_{}{}^{\alpha },x_{}^{}{}_{}{}^{\beta },u_i^\pm \}`$ defining the noncommutative HS($`\eta `$,0) are immediately obtained from (3) by setting $`\theta ^{(\alpha \beta )}=0`$ which gives
$`[x_{}^{+}{}_{}{}^{\alpha },x_{}^{+}{}_{}{}^{\beta }]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta ^{++}`$
$`[x_{}^{+}{}_{}{}^{\alpha },x_{}^{}{}_{}{}^{\beta }]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta `$ (3.7)
$`[x_{}^{}{}_{}{}^{\alpha },x_{}^{}{}_{}{}^{\beta }]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta ^{}.`$
Such commutation relations have however a nice interpretation in D-brane physics in presence of closed string background fields $`g_{i\alpha j\beta }`$ and $`B_{i\alpha j\beta }`$. Indeed if we consider a D3-brane with Euclidean $`𝐑^4`$ world volume in a constant antisymmetric $`B`$ field and world sheet action
$$𝒮=\frac{1}{4\pi \alpha ^{}}_\mathrm{\Sigma }\left(g_{i\alpha j\beta }\delta ^{AB}2i\pi B_{i\alpha j\beta }\epsilon ^{AB}\right)_Ax^{i\alpha }_Bx^{j\beta },$$
(3.8)
consistency requires that the open string fields $`x^{i\alpha }`$ should obey the following boundary conditions
$$\left(g_{i\alpha j\beta }_nx^{j\beta }+2i\alpha ^{}B_{i\alpha j\beta }_tx^{j\beta }\right)|_\mathrm{\Sigma }=0,$$
(3.9)
where $`_n`$ and $`_t`$ are the normal and tangential derivatives along the world sheet boundary $`\mathrm{\Sigma }`$ respectively. In practice we will use the classical approximation to open string where $`\mathrm{\Sigma }`$ may be taken as a disc and hence can be conformally mapped to the upper half plane parametrized by $`z=\tau +i\sigma `$ and $`=\frac{1}{2}(\frac{}{\tau }i\frac{}{\sigma })`$,0 $`\sigma 0`$. In this world sheet coordinates, the boundary conditions (3.9) becomes
$$\left[g_{i\alpha j\beta }\left(\overline{}\right)x^{j\beta }+2i\alpha ^{}B_{i\alpha j\beta }(+\overline{})x^{j\beta }\right]|_{z=\overline{z}}=0.$$
(3.10)
Eq (3.10) describes an interpolation from Neumann to Dirichlet boundary conditions. In the Seiberg Witten limit where $`g_{i\alpha j\beta }\alpha _{}^{}{}_{}{}^{2}\epsilon _{ij}\mathrm{\Omega }_{\alpha \beta }`$ with $`\alpha ^{}ϵ^{\frac{1}{2}}`$ goes to zero, Eq (3.10) becomes Dirichlet boundary conditions where at each boundary the open string world sheet is attached to a single point (zero brane) in the D3 brane. In this limit , the action (3.8) reduces to
$$𝒮=\frac{i}{2}_𝐑𝑑\tau \left(B_{ij}\mathrm{\Omega }^{\alpha \beta }x^{i\alpha }\frac{dx^{j\beta }}{d\tau }+B_{\alpha \beta }\epsilon _{ij}x^{i\alpha }\frac{dx^{j\beta }}{d\tau }\right).$$
(3.11)
Now taking the anti self-dual part $`B_{(\alpha \beta )}`$ of the antisymmetric NS-NS field to zero, one can calculate the propagators of the boundary fields $`x^{i\alpha }(\tau )`$. From the eqs of motion of the one dimensional fields $`x^{k\gamma }(\tau )`$, namely
$$i\epsilon ^{ki}B_{(ij)}\frac{dx^{j\gamma }}{d\tau }=0,$$
(3.12)
one can easily check that the propagators of the $`x^{i\alpha }(\tau )`$’s read as
$$x^{i\alpha }(\tau _1)x^{j\beta }(\tau _2)=\frac{i}{2}\eta ^{(ij)}\mathrm{\Omega }^{\alpha \beta }\epsilon (\tau _1\tau _2),$$
(3.13)
where $`\epsilon (\tau _1\tau _2)`$ is the Heveaside function $`\epsilon (\tau )=1`$ for $`\tau >0`$ and $`\epsilon (\tau )=1`$ for $`\tau <0`$ and where
$$\eta ^{(ij)}B_{(ij)}=1$$
(3.14)
In harmonic space $`\{x_{}^{+}{}_{}{}^{\alpha },x_{}^{}{}_{}{}^{\alpha },u^\pm \}`$, the boundary conformal field theory we have discussed is described by the following action
$$𝒮=\frac{i}{2}_{𝐑\times 𝐒^\mathrm{𝟐}}𝑑\tau 𝑑u\mathrm{\Omega }_{\alpha \beta }\left[\left(B^{++}x_{}^{}{}_{}{}^{\alpha }\frac{dx_{}^{}{}_{}{}^{\beta }}{d\tau }B^{}x_{}^{+}{}_{}{}^{\alpha }\frac{dx_{}^{+}{}_{}{}^{\beta }}{d\tau }\right)B\left(x_{}^{+}{}_{}{}^{\alpha }\frac{dx_{}^{}{}_{}{}^{\beta }}{d\tau }+x_{}^{}{}_{}{}^{\alpha }\frac{dx_{}^{+}{}_{}{}^{\beta }}{d\tau }\right)\right],$$
(3.15)
where integration with respect to the $`u`$’s keeps only SU(2) singlets and where we have used
$`B^{++}`$ $`=u_{(i}^+u_{j)}^+B^{(ij)}`$
$`B^{}`$ $`=u_{(i}^{}u_{j)}^{}B^{(ij)}`$
$`B`$ $`=u_{(i}^+u_{j)}^{}B^{(ij)}.`$
The propagators in the harmonic space are
$`x_{}^{+}{}_{}{}^{i\alpha }(\tau _1)x_{}^{+}{}_{}{}^{j\beta }(\tau _2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{\Omega }^{\alpha \beta }\eta ^{++}\epsilon (\tau _{12}),`$
$`x_{}^{}{}_{}{}^{i\alpha }(\tau _1)x_{}^{}{}_{}{}^{j\beta }(\tau _2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{\Omega }^{\alpha \beta }\eta ^{}\epsilon (\tau _{12}),`$
$`x_{}^{+}{}_{}{}^{i\alpha }(\tau _1)x_{}^{}{}_{}{}^{j\beta }(\tau _2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{\Omega }^{\alpha \beta }\eta \epsilon (\tau _{12}),`$
$`x_{}^{}{}_{}{}^{i\alpha }(\tau _1)x_{}^{+}{}_{}{}^{j\beta }(\tau _2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{\Omega }^{\alpha \beta }\eta \epsilon (\tau _{12}),`$ (3.17)
Computing the commutators of the conformal field operators $`x_{}^{\pm }{}_{}{}^{\alpha }(\tau )`$ by using the short distance products (3.1), we find the following relations
$`[x_{}^{+}{}_{}{}^{i\alpha }(\tau _1),x_{}^{+}{}_{}{}^{j\beta }(\tau _2)]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta ^{++},`$
$`[x_{}^{}{}_{}{}^{i\alpha }(\tau _1),x_{}^{}{}_{}{}^{j\beta }(\tau _2)]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta ^{},`$
$`[x_{}^{+}{}_{}{}^{i\alpha }(\tau _1),x_{}^{}{}_{}{}^{j\beta }(\tau _2)]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta ,`$
$`[x_{}^{}{}_{}{}^{i\alpha }(\tau _1),x_{}^{+}{}_{}{}^{j\beta }(\tau _2)]`$ $`=`$ $`i\mathrm{\Omega }^{\alpha \beta }\eta ,`$ (3.18)
which are similar to the commutation relations of the noncommutative harmonic space NHS($`\eta `$,0) given by (3). In the remainder of this section we want to discuss some features of the algebra of functions on NHS($`\eta `$,0). In conformal field theory on harmonic space one distinguishes different ground state vertex operators;
$$\mathrm{exp}(\pm ip^+x^{});\mathrm{exp}(\pm ip^{}x^+).$$
(3.19)
They satisfy the short distance products
$`:e^{+ip^{}x^+}(1)::e^{+iq^{}x^+)}(2):`$ $`=`$ $`e^{\frac{1}{2}\eta ^{++}p^{}q^{}}:e^{i(p^{}+q^{})x^+}(2):`$
$`:e^{ip^+x^{}}(1)::e^{iq^+x^{})}(2):`$ $`=`$ $`e^{\frac{1}{2}\eta ^{}p^+q^+}:e^{i(p^++q^+)x^{}}(2):`$
$`:e^{+ip^{}x^+}(1)::e^{iq^+x^{})}(2):`$ $`=`$ $`e^{\frac{1}{2}\eta p^{}q^+}:e^{+i(p^{}x^+q^+x^{})}(2):`$
$`:e^{ip^+x^{}}(1)::e^{iq^{}x^+)}(2):`$ $`=`$ $`e^{\frac{1}{2}\eta p^+q^{}}:e^{i(p^+x^{}q^{}x^+)}(2):.`$ (3.20)
In noncommutative harmonic space language, the above short distance products coincide with the usual star product on noncommutative geometry. Thus we have for instance the identification
$$\underset{\tau 0}{lim}e^{\pm ip^{}x^+}(\tau ).e^{\pm iq^{}x^+}(0)e^{\pm ip^{}x^+}e^{\pm iq^{}x^+}.$$
(3.21)
More generally given two harmonic space functions $`f^{(q_1)}(x^+,x^{},u)`$ and $`g^{(q_2)}(y^+,y^{},u)`$, the star product of these functions is defined as
$$f^{(q_1)}(x)g^{(q_2)}(y)=\mathrm{exp}\frac{i}{2}(\mathrm{\Omega }^{\alpha \beta }M_{\alpha \beta })f^{(q_1)}(x)g^{(q_2)}(y),$$
(3.22)
where
$$M_{\alpha \beta }=\left(\eta ^{++}\frac{}{x^{+\alpha }}\frac{}{y^{+\beta }}+\eta ^{}\frac{}{x^\alpha }\frac{}{y^\beta }\right)\eta \left(\frac{}{x^{+\alpha }}\frac{}{y^\beta }+\frac{}{x^\alpha }\frac{}{y^{+\beta }}\right).$$
(3.23)
At first order in $`\eta ^{++}`$, $`\eta ^{}`$ and $`\eta `$, (3.22) reduces to
$`f^{(q_1)}g^{(q_2)}`$ $`=`$ $`f^{(q_1)}g^{(q_2)}+{\displaystyle \frac{i}{2}}\mathrm{\Omega }^{\alpha \beta }[(\eta ^{++}_\alpha ^{}f^{(q_1)}_\beta ^{}g^{(q_2)}+\eta ^{}_\alpha ^+f^{(q_1)}_\beta ^+g^{(q_2)})`$ (3.24)
$`\eta (_\alpha ^{}f^{(q_1)}_\beta ^+g^{(q_2)}+_\alpha ^+f^{(q_1)}_\beta ^{}g^{(q_2)})]+𝒪(2)`$
### 3.2 Self-Dual Yang-Mills Constraints in NHS($`\eta `$,0)
Noncommutative YM theory in Euclidean four dimensional space is formulated in a similar way as YM theory in ordinary $`𝐑^4`$ except that the gauge group matrix multiplication is now replaced by the star product (3.22). For instance the transformations of the gauge field $`\widehat{A}_M`$ and the field strength $`\widehat{F}_{MN}`$ of noncommutative YM theory are
$`\delta _{\widehat{\tau }}\widehat{A}_\mu `$ $`=`$ $`_\mu \widehat{\tau }i\left(\widehat{\tau }\widehat{A}_\mu \widehat{A}_\mu \widehat{\tau }\right)`$
$`\widehat{F}_{MN}`$ $`=`$ $`_M\widehat{A}_N_N\widehat{A}_Mi\left(\widehat{A}_M\widehat{A}_N\widehat{A}_N\widehat{A}_M\right)`$ (3.25)
$`\delta _{\widehat{\tau }}\widehat{F}_{MN}`$ $`=`$ $`i\left(\widehat{\tau }\widehat{F}_{MN}\widehat{F}_{MN}\widehat{\tau }\right).`$
Remark that the gauge parameters $`\widehat{\tau }`$, the fields $`\widehat{A}_M`$ and $`\widehat{F}_{MN}`$ carry a hat in order to be distinguished from their ordinary geometry analogue. This convention notation will also be used in the remainder of this study.
YM gauge theory in noncommutative harmonic space may be constructed in a similar manner as in the standard formulation. This is achieved in practice by extending the usual HS(0,0) classical fields to functionals on NHS($`\eta `$,0) and replacing the ordinary product by the star one given by (3.22). The novelty brought by the harmonic variables is completely controlled by the SU(2) symmetry. Since the $`u_i^\pm `$ variables still obey the same relations as in ordinary HS(0,0), the covariant harmonic derivatives in NHS($`\eta `$,0) defined as
$`𝒟^{++}`$ $`=D^{++}+i\widehat{V}^{++}`$
$`𝒟^{}`$ $`=D^{}+i\widehat{V}^{}`$
$`𝒟^0`$ $`=D^0,`$
where $`\widehat{V}^{++}`$ and $`\widehat{V}^{}`$ are harmonic gauge connections on HS($`\eta `$,0), still obey the SU(2) algebra (2.22) which requires in turns:
$`[D^0,\widehat{V}^{++}]`$ $`=2\widehat{V}^{++};`$
$`[D^0,\widehat{V}^{}]`$ $`=2\widehat{V}^{}`$
$`D^{++}\widehat{V}^{}D^{}\widehat{V}^{++}+i\left(\widehat{V}^{++}\widehat{V}^{}\widehat{V}^{}\widehat{V}^{++}\right)`$ $`=0.`$
The SU(2) symmetry (2.22) and (3.2,3.2) shows also that the self-dual YM constraint eqs in noncommutative NHS($`\eta `$,0) may be obtained from (2.1) by replacing the HS(0,0) objects by their extensions on NHS($`\eta `$,0). The noncommutative self-dual YM constraints read then:
$`[D^0,\widehat{V}^{++}]`$ $`=2\widehat{V}^{++}(a)`$
$`[_\alpha ^+,\widehat{V}^{++}]`$ $`=0(b)`$
$`D^{++}\widehat{V}^{}D^{}\widehat{V}^{++}`$ $`=i\left(\widehat{V}^{++}\widehat{V}^{}\widehat{V}^{}\widehat{V}^{++}\right)(c)`$
$`[_\alpha ^+,\widehat{V}^{}]`$ $`=\widehat{A}_\alpha ^{}(d)`$
$`[_\alpha ^+,\widehat{A}_\beta ^{}]`$ $`=\widehat{F}_{(\alpha \beta )}(e)`$
To solve these constraint eqs, we shall make two analysis by considering first perturbative solutions around the ordinary one. Then we examine the exact solution of these constraints by using noncommutative calculus on NHS($`\eta `$,0).
### 3.3 Perturbative Analysis
Here we shall consider special perturbative solutions preserving manifestly the SU(2) symmetry and corresponding to small values of the deformation parameters $`\eta ^{++}`$, $`\eta ^{}`$ and $`\eta `$. Eqs (3.2) suggest that one may expand the noncommutative connections $`\widehat{V}^{++}`$ and $`\widehat{V}^{}`$ around the ordinary $`V^{++}`$ and $`V^{}`$ ones of (2.29, 2.30) as follows
$`\widehat{V}^{++}`$ $`=`$ $`V^{++}+\eta W^{++}+𝒪(2)`$
$`\widehat{V}^{}`$ $`=`$ $`V^{}+\eta U^{}+\eta ^{}U+\eta ^{++}U^4+𝒪(2)`$ (3.29)
Expanding (3.2.c) to first order in $`\eta ^{++}`$, $`\eta ^{}`$ and $`\eta `$ as in (3.22) and using (3.2,3.3), in particular the analyticity of $`\widehat{V}^{++}`$, we find
$`\eta `$ $`\left(D^{++}U^{}+2U\right)+\eta ^{}D^{++}U`$ (3.30)
$`+\eta ^{++}\left(D^{++}U^4+U^{}\right)\eta D^{}W^{++}\eta ^{}W^{++}`$
$`=i(\eta \{[V^{++},U^{++}]+[W^{++},V^{}]\}`$
$`+\eta ^{}[V^{++},U]+\eta ^{++}[V^{++},U^4])+{\displaystyle \frac{1}{2}}\mathrm{\Omega }^{\alpha \beta }\{\eta ^{++}(_\alpha ^{}V^{++}_\beta ^{}V^{}`$
$`_\alpha ^{}V^{}_\beta ^{}V^{++})\eta (_\alpha ^{}V^{++}_\beta ^+V^{}_\alpha ^+V^{}_\beta ^{}V^{++})\}+𝒪(2).`$
These eqs imply in turn
$`D^{++}UW^{++}+i[V^{++},U]=0`$
$`D^{++}U^{}D^{}W^{++}+2U+i[V^{++},U^{}]+i[W^{++},V^{}]={\displaystyle \frac{1}{2}}\mathrm{\Omega }^{\alpha \beta }\{_\alpha ^{}V^{++}_\beta ^+V^{}`$
$`_\alpha ^+V^{}_\beta ^{}V^{++})`$
$`D^{++}U^4+i[V^{++},U^4]+U^{}={\displaystyle \frac{1}{2}}\mathrm{\Omega }^{\alpha \beta }\{_\alpha ^{}V^{++}_\beta ^{}V^{}`$
$`_\alpha ^{}V^{}_\beta ^{}V^{++})`$ (3.31)
To solve these eqs, we choose the fields $`W^{++}`$, $`U`$, $`U^{}`$ and $`U^4`$ of the form
$`W_{}^{++}{}_{\beta }{}^{\alpha }`$ $`=ax_{}^{+}{}_{}{}^{\alpha }x_{}^{+}{}_{\beta }{}^{}`$
$`U_\beta ^\alpha `$ $`=b_1x_{}^{+}{}_{}{}^{\alpha }x_{}^{}{}_{\beta }{}^{}+b_2x_{}^{}{}_{}{}^{\alpha }x_{}^{+}{}_{\beta }{}^{}`$
$`U_{}^{}{}_{\beta }{}^{\alpha }`$ $`=cx_{}^{}{}_{}{}^{\alpha }x_{}^{}{}_{\beta }{}^{}`$
where $`a`$, $`b_1`$, $`b_2`$ and $`c`$ are parameters to be determined. Lengthy but straightforward calculations lead to:
$`a`$ $`={\displaystyle \frac{1}{\rho ^4}}`$ (3.33)
$`b_1`$ $`={\displaystyle \frac{1}{2(x^2+\rho ^2)^2}}`$
$`b_2`$ $`={\displaystyle \frac{2x^2+\rho ^2}{2\rho ^2(x^2+\rho ^2)^2}}`$
$`c`$ $`={\displaystyle \frac{1}{(x^2+\rho ^2)^2}}.`$
## 4 Exact Solution
Here we give the exact solution of the self-dual YM constraint (3.2) on noncommutative HS($`\eta `$,0). As for ordinary HS(0,0) harmonic space, (3.2.a-b) show that $`\widehat{V}^{++}`$ is a harmonic function on NHS($`\eta `$,0) depending on $`x_{}^{+}{}_{}{}^{\alpha }`$ only. This means that according to (3.2.a-b), $`\widehat{V}^{++}`$ may be written as
$$(V^{++})^{\alpha \beta }=x_{}^{+}{}_{}{}^{\alpha }Ax_{}^{+}{}_{}{}^{\beta }+C^{++}\mathrm{\Omega }^{\alpha \beta }$$
(4.1)
where $`A`$ and $`C^{++}`$ are harmonic functions without any dependence on the $`x^\pm `$’s. The $`A`$ and $`C^{++}`$ scale as the inverse of length squared. Therefore, they behave as the inverse of $`\eta `$ since according to (3)), one can check that we have the following identities
$`x_{}^{+}{}_{\gamma }{}^{}x_{}^{+}{}_{}{}^{\gamma }`$ $`=`$ $`x^+x^+=i\eta ^{++}`$
$`x_{}^{}{}_{\gamma }{}^{}x_{}^{}{}_{}{}^{\gamma }`$ $`=`$ $`x^{}x^{}=i\eta ^{}`$
$`x_{}^{+}{}_{\gamma }{}^{}x_{}^{}{}_{}{}^{\gamma }`$ $`=`$ $`x^+x^{}=i\eta +z^2`$
$`x_{}^{}{}_{\gamma }{}^{}x_{}^{+}{}_{}{}^{\gamma }`$ $`=`$ $`x^{}x^+=i\eta z^2,`$ (4.2)
where we have used $`\mathrm{\Omega }_{\alpha \beta }\mathrm{\Omega }^{\alpha \beta }=2`$ and set $`z^2=\frac{1}{2}(x^+x^{}x^{}x^+)0`$. Note that (4.1,4) satisfy
$`D^{}(V^{++})^{\alpha \beta }`$ $`=`$ $`\left(x_{}^{+}{}_{}{}^{\alpha }Ax_{}^{}{}_{}{}^{\beta }+x_{}^{}{}_{}{}^{\alpha }Ax_{}^{+}{}_{}{}^{\beta }\right)`$ (4.3)
$`+x_{}^{+}{}_{}{}^{\alpha }(D^{}A)x_{}^{+}{}_{}{}^{\beta }+\left(D^{}C^{++}\right)\mathrm{\Omega }^{\alpha \beta }.`$
and
$`D^{++}\eta ^{++}`$ $`=0`$
$`D^{++}\eta `$ $`=\eta ^{++}`$
$`D^{++}\eta ^{}`$ $`=2\eta .`$
The next step is to find $`(V^{})^{\alpha \beta }`$ by solving (3.2.c). Harmonic analysis on $`NHS(\eta ,0)`$ shows that we should look for a solution of the form
$`(V^{})^{\alpha \beta }`$ $`=`$ $`x_{}^{}{}_{}{}^{\alpha }Bx_{}^{}{}_{}{}^{\beta }+x_{}^{+}{}_{}{}^{\alpha }E^{}x_{}^{}{}_{}{}^{\beta }`$ (4.5)
$`+x_{}^{}{}_{}{}^{\alpha }K^{}x_{}^{+}{}_{}{}^{\beta }+x_{}^{+}{}_{}{}^{\alpha }G^{(4)}x_{}^{+}{}_{}{}^{\beta }+H^{}\mathrm{\Omega }^{\alpha \beta }`$
where $`B`$, $`E^{}`$, $`K^{}`$, $`G^{(4)}`$ and $`H^{}`$ are to be determined . Taking the harmonic derivative of (4.5), we get
$`D^{++}(V^{})^{\alpha \beta }`$ $`=`$ $`x_{}^{+}{}_{}{}^{\alpha }\left[B+D^{++}E^{}\right]x_{}^{}{}_{}{}^{\beta }+x_{}^{}{}_{}{}^{\alpha }\left[B+D^{++}K^{}\right]x_{}^{+}{}_{}{}^{\beta }`$
$`+x_{}^{+}{}_{}{}^{\alpha }\left[E^{}+K^{}+D^{++}G^{(4)}\right]x_{}^{+}{}_{}{}^{\beta }`$
$`+x_{}^{}{}_{}{}^{\alpha }\left[D^{++}B\right]x_{}^{}{}_{}{}^{\beta }+D^{++}H^{}\mathrm{\Omega }^{\alpha \beta }.`$
Moreover using (4.1) and (4.5) we have
$`i(V_{}^{++}{}_{\gamma }{}^{\alpha }V_{}^{}{}_{\beta }{}^{\gamma }`$ $``$ $`V_{}^{}{}_{\gamma }{}^{\alpha }V_{}^{++}{}_{\beta }{}^{\gamma })`$
$`=`$ $`x_{}^{+}{}_{}{}^{\alpha }\left[iA(x^+x^{})B+iA(x^+x^+)E^{}\right]x_{}^{}{}_{}{}^{\beta }`$
$`+x_{}^{}{}_{}{}^{\alpha }\left[iB(x^{}x^+)AiK^{}(x^+x^+)A\right]x_{}^{+}{}_{}{}^{\beta }`$
$`+x_{}^{+}{}_{}{}^{\alpha }[iA(x^+x^{})K^{}+iA(x^+x^+)G^{(4)}`$
$`E^{}(x^{}x^+)AiG^{(4)}(x^+x^+)A]x^+^\beta .`$
Putting back (4.1,4),(4.5) and (4) in the constraint (3.2.c) one gets an equation transforming in the $`\mathrm{𝟏}\mathrm{𝟑}`$ representation of the SU(2) symmetry which up on projecting it along the different U(1) cartan subsymmetry directions, we get the following system of four equations
$`\left[D^{++}E^{}+iA(x^+x^+)E^{}+\left[1+iA(x^+x^{})\right]BA\right]`$ $`=0(a)`$
$`[D^{++}K^{}iK^{}(x^+x^+)A+B[1iA(x^{}x^+)]A]A]`$ $`=0(b)`$
$`[D^{++}G^{(4)}+E^{}(1i(x^{}x^+)A)`$
$`+[1+iA(x^+x^{})]K^{}D^{}A]`$ $`=0(c)`$
$`D^{++}B`$ $`=0(d)`$
$`D^{++}H^{}D^{}C^{++}`$ $`=0(e)`$
To solve these, it is interesting to note the following: (1) the $`\eta `$ parameter is a C number independent of the $`x^\pm `$’s and so commute with the star product. it scales as (length)<sup>-2</sup>. $`\eta ^{}`$ and $`\eta ^{++}`$ are proportional to the norm of the sp(1) isospinor $`x_{}^{+}{}_{}{}^{\alpha }`$ and $`x_{}^{}{}_{}{}^{\alpha }`$ which are no longer zero due to noncommutativity. (2) In the limit $`\eta 0`$, $`x^+x^{}0`$ and $`x^{}x^+0`$ as may be check explicitly by help of (2.1) (3) Dimensional arguments show that the unknown functions $`A`$, $`B`$, $`C^{}`$, $`E^{}`$,$`K^{}`$ and $`G^{(4)}`$ in (4) scale as
$`\left[A\right]`$ $`=`$ $`\left[B\right]=\left[E^{}\right]=\left[K^{}\right]=\left[G^4\right]=\left[\eta \right]`$ (4.9)
$`=`$ $`2\left[x^\pm \right]=(\mathrm{length})^2`$
(4) In the limit $`\eta 0`$, eqs (4) should coincide with the ordinary case and so
$`\underset{\eta 0}{lim}A`$ $`=i`$ $`{\displaystyle \frac{1}{\rho ^2}}`$
$`\underset{\eta 0}{lim}B`$ $`=i`$ $`{\displaystyle \frac{1}{\rho ^2+z^2}}`$ (4.10)
$`\underset{\eta 0}{lim}E^{}`$ $`=`$ $`\underset{\eta 0}{lim}K^{}=\underset{\eta 0}{lim}G^{(4)}=0.`$
These limits show that $`E^{}`$, $`K^{}`$ and $`G^4`$ should be proportional to $`\eta `$. Eq (4.d) shows that $`B`$ does not depend on the harmonics. This means that if it carries a dependence on $`\eta `$, then this should be in terms of the invariant $`\eta ^{++}\eta ^{}\eta ^2=\eta ^{(ij)}\eta _{ij}=\stackrel{}{\eta }^2`$. However the dimensional arguments show that the natural solution is just the ordinary one which reads as
$$B=i\frac{2}{(2\rho ^2+x^+x^{}x^{}x^+)}=i\frac{1}{\rho ^2+z^2},$$
(4.11)
where we have used (4).
Now we turn to solve (4.a-b) which we rewrite using (4) as
$`\left[D^{++}E^{}\eta ^{++}AE^{}\right]+\left[(1+iAx^+x^{})BA\right]`$ $`=`$ $`0(a)`$
$`\left[D^{++}K^{}\eta ^{++}AK^{}\right]+\left[(1iAx^{}x^+)BA\right]`$ $`=`$ $`0(b).`$ (4.12)
Taking $`A`$ to be
$$A=i\frac{1}{\rho ^2i\eta }$$
(4.13)
and using (4.11), the last terms of eq (4.a) vanish. It follows that (4) become
$`\left(D^{++}+i{\displaystyle \frac{\eta ^{++}}{\rho ^2i\eta }}\right)E^{}`$ $`=0`$ (4.14)
$`\left(D^{++}i{\displaystyle \frac{\eta ^{++}}{\rho ^2i\eta }}\right)K^{}+2\eta AB`$ $`=0.`$
Integrating these eqs, we find that the most general solution is given by
$`E^{}`$ $`=`$ $`\mathrm{cnst}{\displaystyle \frac{\rho ^2i\eta }{(\stackrel{}{\eta }^2+\eta ^2)}}`$
$`K^{}`$ $`=`$ $`\eta ^{}AB`$ (4.15)
$`=`$ $`{\displaystyle \frac{\eta ^{}}{(\rho ^2i\eta )(2\rho ^2+x^+x^{}x^{}x^+)}}`$
However the constant appearing in the solution of (4) should vanish due to the constraints (4). Therefore
$$E^{}=0$$
(4.16)
Now we turn to solve (4.c). Using the previous solutions (4.11), (4.13), (4) and (4.16), eq (4.c) can be brought to the following form
$$D^{++}G^{(4)}+\frac{A}{B}K^{}=D^{}A.$$
At first sight, this eq seems to be difficult to handle; but if we consider the following remarkable features
$`K^{}`$ $`=`$ $`AB\eta ^{}`$
$`D^{}A`$ $`=`$ $`\eta ^{}A^2,`$ (4.17)
it reduces to
$$D^{++}G^{(4)}=0$$
(4.18)
and so
$$G^{(4)}=0.$$
(4.19)
Finally, concerning (4.e), the $`H^{}`$ and $`C^{++}`$ are solved as
$`C^{++}`$ $`=`$ $`D^{++}\lambda `$
$`H^{}`$ $`=`$ $`D^{}\lambda ,`$ (4.20)
where $`\lambda `$ is an arbitrary function on the $`𝐒^2`$.
Summary
The noncommutative YM SU(2) instanton formulated on NHS($`\eta `$,0) is described by the harmonic connections $`V_{}^{++}{}_{\beta }{}^{\alpha }`$ and $`V_{}^{}{}_{\beta }{}^{\alpha }`$ given
$`V_{}^{++}{}_{\beta }{}^{\alpha }`$ $`=`$ $`x_{}^{+}{}_{}{}^{\alpha }Ax_{}^{+}{}_{\beta }{}^{}+D^{++}\lambda \delta _\beta ^\alpha `$
$`V_{}^{}{}_{\beta }{}^{\alpha }`$ $`=`$ $`x_{}^{}{}_{}{}^{\alpha }Bx_{}^{}{}_{\beta }{}^{}\eta ^{}x_{}^{}{}_{}{}^{\alpha }ABx_{}^{+}{}_{\beta }{}^{}`$
$`+D^{}\lambda \delta _\beta ^\alpha `$
with
$`A`$ $`=`$ $`i{\displaystyle \frac{1}{\rho ^2i\eta }}`$
$`B`$ $`=`$ $`i{\displaystyle \frac{1}{\rho ^2i\eta +x^+x^{}}}.`$ (4.22)
At this level certain interesting comments may be done: (1) The solutions of the noncommutative YM self-dual constraint eqs we have obtained are exact solutions. They extend the perturbative one given in the end of section 3.
(2) our solution admits a U(1) symmetry;
$`V^{++}`$ $``$ $`V^{++}+D^{++}\lambda `$
$`V^{}`$ $``$ $`V^{}+D^{}\lambda `$ (4.23)
for any $`\lambda `$. (3) As long as $`\eta 0`$, our solution is non singular. This result agrees with the absence of small instanton in noncommutative geometry. (4) In the limit $`\eta 0`$, we recover the ordinary solution. Using the associativity of the star product, (3.22) and the algebra of the differential operators
$`_\mu ^{}`$ $`=`$ $`\eta ^{}_\mu ^+\eta _\mu ^{}`$
$`_\mu ^+`$ $`=`$ $`\eta ^{++}_\mu ^{}\eta _\mu ^+,`$ (4.24)
one discovers the perturbative solution obtained in section 3 (3.3-3.33). (5) In the ordinary geometry, the parameter $`\rho ^2`$ is interpreted as the size of the YM instantons. It is also a real Kahler parameter of the resolution of ADE singularities by blowing up two-spheres. In the noncommutative YM theory in NHS($`\eta `$,0), $`\rho ^2`$ is shifted by $`\eta `$ and becomes a quaternionic parameter $`\varrho `$ as shown here below
$`\varrho `$ $`=`$ $`\rho ^2i\eta `$ (4.25)
$`=`$ $`\rho ^2iu_{(k}^+u_{l)}^{}\eta ^{(kl)}.`$
(6) Finally observe that it is possible to work out an other solution of the constraint (3.2) by taking
$$A=i\frac{1}{\rho ^2+i\eta }.$$
(4.26)
This choice affects (4) which becomes
$`\left(D^{++}+i{\displaystyle \frac{\eta ^{++}}{\rho ^2+i\eta }}\right)E^{}2\eta AB`$ $`=0`$ (4.27)
$`\left(D^{++}+i{\displaystyle \frac{\eta ^{++}}{\rho ^2+i\eta }}\right)K^{}`$ $`=0.`$
The second class of solutions of (3.2) are given by (4) and (4) up to performing $`\eta \eta `$ and $`K^{}E^{}`$ .
## 5 Conclusion
In this paper we have studied Yang-Mills Instantons on noncommutative harmonic space NHS($`\eta `$,0). This approach has the advantage of allowing to explicit exact solutions of the noncommutative self dual Yang Mills constraint eqs. It also has the merit of going beyond the perturbative solution described in .
We have first developed harmonic space noncommutative geometry and have shown that NHS($`\eta ,\theta `$) has two remarkable NHS subspaces in addition to the usual ordinary one. This property may be seen by remarking that NHS($`\eta ,\theta `$) has two deformation parameters $`\eta `$ and $`\theta `$ transforming as (1,0) and (0,1) representations of SU(2)$`\times `$SP(1)$``$ SO(4) Lorentz group. According to whether $`\eta `$ and $`\theta `$ are zero or not, we obtain four subspaces, three of them are noncommutative. Second we have reformulated the noncommutative self-duality Yang-Mills constraints in NHS($`\eta ,\theta `$) by extending the idea of harmonic analyticity. In this formulation, the basic objects carrying the Yang-Mills self-dual constraints are given by the harmonic connections $`\widehat{V}^{++}`$ and $`\widehat{V}^{}`$ of a gauged SU(2) symmetry. The latter is a generalisation of a well known trick allowing to go from the standard analysis to the harmonic one. The noncommutative gauge fields $`\widehat{A}_M`$ and curvatures $`\widehat{F}_{\mu \nu }`$ are related to $`\widehat{V}^{++}`$ and $`\widehat{V}^{}`$ as shown in eqs (3.2). Third we have studied the solutions of (3.2) by considering in a first step perturbative solutions around the ordinary one. In a second step, we have given an explicit derivation of an exact solution of the self-dual Yang-Mills constraints. It should be noted here that besides the power of harmonic space analysis, this exact solution has been made possible due to the choice of $`\theta ^{(\alpha \beta )}=0`$ and useful properties described in section 4.
Acknowledgements
Belhaj would like to thank the organizers of the Spring Workshop on Superstrings and related Matters ( March 2000), The Abdus Salam International Centre for theoretical Physics Trieste, Italy for hospitality ; where a part of this work is done. He is grateful to N. Nekrasov and S. Kachru for valuable discussions.
Sahraoui would like to thank DFG under contract 445 Mar 113/5/0 to have supported his stay at the university of Muenchen, Pro. J. Wess to have invited him and Dr. R. Dick for his kindness and hospitality. Many thanks for Pro S. Theisen for his helpful discussions and suggestions.
This research is supported by the program SARS 99/2000. contract CNCPRST-Université Mohammed V, Rabat.
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# Screening effect due to heavy lower tails in one-dimensional parabolic Anderson model
## 1. Introduction
### 1.1. Model and main aim
In a recent paper \[BK00\], we have studied the asymptotic behavior of the solution $`u(t,z)`$ to the so-called parabolic Anderson model for non-positive i.i.d. potentials. Here we answer an open question concerning the almost-sure asymptotics of $`u(t,0)`$ as $`t\mathrm{}`$ in dimension one for potentials lacking the first logarithmic moment. Interestingly, a new phenomenon arises: too heavy tails of the potential at $`\mathrm{}`$ hamper the mass flow to remote areas, thus rendering the more favorable regions inaccessible. This effect is unique to $`d=1`$, since only in one-dimensional topology particles are not able to bypass deep broad valleys in the potential landscape.
The general model is defined as follows. Let $`u:[0,\mathrm{})\times ^d[0,\mathrm{})`$ be the solution to the parabolic problem
$$\begin{array}{cccc}\hfill _tu(t,z)& =& \kappa \mathrm{\Delta }^\mathrm{d}u(t,z)+\xi (z)u(t,z),\hfill & (t,z)(0,\mathrm{})\times ^d,\hfill \\ \hfill u(0,z)& =& 1,\hfill & z^d,\hfill \end{array}$$
(1.1)
where $`_t`$ is the time derivative, $`\kappa >0`$ is the diffusion constant, $`\mathrm{\Delta }^\mathrm{d}`$ is the discrete Laplacian on $`^d`$, $`[\mathrm{\Delta }^\mathrm{d}f](z)=_{xz}[f(x)f(z)]`$, and $`\xi =(\xi (z))_{z^d}`$ is an i.i.d. field. We use $``$ to denote the expectation with respect to $`\xi `$ and $`\mathrm{Prob}()`$ to denote the underlying probability measure. One interpretation of the quantity $`u(t,z)`$ is the total expected mass accumulated at time $`t`$ by a particle starting at $`z`$ at time $`0`$ and diffusing through a random field of sources (sites $`x`$ with $`\xi (x)>0`$) and sinks (sites $`x`$ with $`\xi (x)<0`$). The references \[GM90\], \[CM94\] and \[K00\] provide more explanation and other interpretations.
Besides \[BK00\], the large-$`t`$ behavior of the solution to (1.1) has extensively been studied (in general dimension) for various other classes of distributions: see \[GM90, GM98, GH99\] for $`\xi `$ having the so-called double-exponential upper tail, and \[GK98, GKM99\] for a continuous variant of (1.1) with $`\xi `$ either Gaussian or (smeared) Poissonian field. The techniques used in these studies go back to the pioneering work of Donsker and Varadhan \[DV75, DV79\]; however, there is also an intimate relation to Sznitman’s method of enlargement of obstacles \[S98\]. We refer to \[K00\] for a comprehensive discussion of these relations and a unified presentation of the above results. Henceforth, we shall focus on the almost-sure behavior of $`u(t,0)`$ in the non-positive case, i.e., $`\xi [\mathrm{},0]^^d`$.
In dimensions $`d2`$, the analysis in \[BK00\] produced a fairly complete picture. Indeed, interesting behavior occurs only when $`p=\mathrm{Prob}(\xi (0)>\mathrm{})>p_\text{c}(d)`$, the threshold for site percolation on $`^d`$, and when conditioned on the event that the origin lies in the infinite cluster of sites $`x`$ with $`\xi (x)>\mathrm{}`$. Below and, provided there is no critical percolation (which is rigorously known for $`d=2`$ \[R78\] and $`d19`$ \[HS90\]), also at $`p_\text{c}(d)`$, and also when the origin lies in a finite cluster for $`p>p_\text{c}(d)`$, the quantity $`u(t,0)`$ decays exponentially in $`t`$ with a $`\xi `$-dependent rate.
In dimension $`d=1`$, we have $`p_\text{c}(d)=1`$, which necessitated setting $`\mathrm{Prob}(\xi (0)=\mathrm{})=0`$ in \[BK00\]. However, the latter condition was not sufficient because the existence of the first logarithmic moment, i.e., $`\mathrm{log}(\xi (0)1)<\mathrm{}`$, also had to be assumed in order to establish an asymptotics analogous to the supercritical case in $`d2`$. In particular, two intriguing questions remained unanswered:
1. Is $`\mathrm{log}(\xi (0)1)<\mathrm{}`$ optimal in the sense that $`\mathrm{log}(\xi (0)1)=\mathrm{}`$ implies a strictly different asymptotic behavior of $`u(t,0)`$?
2. What is the precise decay rate when the finiteness of $`\mathrm{log}(\xi (0)1)`$ is robustly violated (keeping however the restriction to “no atom at $`\mathrm{}`$”)?
In this paper we give answers to these questions under mild regularity conditions on the lower tail of the distribution of $`\xi `$. In particular, we show that $`\mathrm{log}(\xi (0)1)<\mathrm{}`$ is only marginally non-optimal for the behavior described in \[BK00\], see Remark 3 after Theorem 1.1. As it turns out, the decay of $`u(t,0)`$ is determined solely by upper and lower tails of $`\mathrm{Prob}(\xi (0))`$. The reason why the intermediate part of the distribution does not play any role is that these tails give rise to two dominant and mutually competing mechanisms (field-shape optimization in the upper tail versus screening effect in the lower tail) whose balancing determines the decay rate. See Subsection 2.2 for more precise heuristic explanation.
### 1.2. Our assumptions
We proceed by stating precisely the needed assumptions, both on upper and lower tails of $`\xi (0)`$. First we restrict ourselves to dimension $`d=1`$ for the sequel of this paper. In accord with \[BK00\], we consider the distributions with the upper tail of the form
$$\mathrm{Prob}\left(\xi (0)x\right)=\mathrm{exp}\left\{x^{\frac{\gamma }{1\gamma }+o(1)}\right\},x0,$$
(1.2)
for some $`\gamma [0,1)`$. However, instead of the distribution function, it is more convenient to work with the cumulant generating function
$$H(\mathrm{})=\mathrm{log}e^{\mathrm{}\xi (0)},\mathrm{}0.$$
(1.3)
The regime in (1.2) corresponds to the behavior $`H(\mathrm{})=\mathrm{}^{\gamma +o(1)}`$ as $`\mathrm{}\mathrm{}`$.
Assumption (H). Let $`\mathrm{esssup}\xi (0)=0`$ and suppose there are constants $`A>0`$ and $`\gamma [0,1)`$, and a positive increasing function $`t\alpha _t`$ such that
$$\underset{t\mathrm{}}{lim}\frac{\alpha _t^3}{t}H\left(\frac{t}{\alpha _t}y\right)=Ay^\gamma ,y>0.$$
(1.4)
The limit in (1.4) is necessarily uniform on compact sets in $`(0,\mathrm{})`$, the pair $`(A,\alpha _t)`$ is unique up to a scaling transformation. Moreover, $`t\alpha _t`$ is regularly varying and $`\alpha _t=t^{\nu +o(1)}`$ as $`t\mathrm{}`$ where $`\nu =(1\gamma )/(3\gamma )(0,1/3]`$. In particular, $`t/\alpha _t\mathrm{}`$, i.e., Assumption (H) indeed controls the upper tails of $`\xi (0)`$. We say that $`H`$ is in the $`\gamma `$-class if (1.4) holds.
Next we formulate our assumption on the lower tails of $`\xi (0)`$ at $`\mathrm{essinf}\xi (0)=\mathrm{}`$. As the opposite case has already been handled in \[BK00\], we shall focus on the case where $`\mathrm{log}(\xi (0)1)`$ is not integrable. Central to our attention are lower tails of the form
$$\mathrm{Prob}\left(\mathrm{log}(\xi (0)1)>x\right)=x^{\zeta +o(1)},x\mathrm{},$$
(1.5)
with some $`\zeta [0,1]`$. In terms of the modified cumulant generating function
$$G(\mathrm{})=\mathrm{log}(\xi (0)1)^{1/\mathrm{}},\mathrm{}>0,$$
(1.6)
the behavior (1.5) roughly corresponds to $`G(\mathrm{})=\mathrm{}^{\zeta +o(1)}`$ as $`\mathrm{}\mathrm{}`$. Note that $`G`$ is positive and decreasing since $`\mathrm{essinf}\xi (0)<1`$. The following is a weak regularity condition for $`G`$ at infinity.
Assumption (G). Let $`\mathrm{log}(\xi (0)1)=\mathrm{}`$ but $`\mathrm{Prob}(\xi (0)=\mathrm{})=0`$. Suppose that for each $`\eta (0,1)`$ there is a function $`\stackrel{~}{G}_\eta :(0,\mathrm{})(0,\mathrm{})`$ with the following properties:
1. $`\stackrel{~}{G}_\eta (\mathrm{})G(\mathrm{})^{\eta +o(1)}`$ as $`\mathrm{}\mathrm{}`$.
2. $`\mathrm{}1/\stackrel{~}{G}_\eta (\mathrm{})`$ is increasing and concave for $`\mathrm{}`$ large enough.
3. The random variable $`1/\stackrel{~}{G}_\eta (\mathrm{log}(\xi (0)1))`$ has the first moment.
###### Remark 1.
As it turns out, Assumption (G) is needed only for the proof of the lower bound in our main result (see Theorem 1.1 below); the upper bound requires no assumptions at all. The role of Assumption (G) and particularly of its part (i) is the following: Abbreviate $`Y=\mathrm{log}(\xi (0)1)`$ and note that, for any $`\delta (0,1]`$, $`G(\mathrm{})Y^\delta \mathrm{}^\delta `$. Therefore, $`G(\mathrm{})\mathrm{}^{\zeta _{}+o(1)}`$ where $`\zeta _{}=sup\{\delta 0:Y^\delta <\mathrm{}\}`$. However, a lower bound of the same (even asymptotic) form requires some regularity of $`\mathrm{}G(\mathrm{})`$ as $`\mathrm{}\mathrm{}`$, which is the essence of (i–iii).
###### Remark 2.
In the view of Remark 1, it is immediate that Assumption (G) holds for $`G(\mathrm{})=\mathrm{}^{\zeta +o(1)}`$ with some $`\zeta (0,1]`$. The reason why we prefer the above (little cumbersome) setting as opposed to simple regularity of $`G`$ is that many cases with $`G(\mathrm{})=\mathrm{}^{o(1)}`$ are automatically included. Indeed, consider the following example: Let $`\theta >0`$ and $`\mathrm{Prob}(\mathrm{log}(\xi (0)1)dx)C/[x\mathrm{log}^{1+\theta }(x)]dx`$ as $`x\mathrm{}`$, where $`C`$ is the normalizing constant. Then $`G(\mathrm{})C^{}(\mathrm{log}\mathrm{})^\theta `$ and Assumption (G) holds with $`\stackrel{~}{G}_\eta (\mathrm{})=G(\mathrm{})[\mathrm{log}\mathrm{log}(\mathrm{}e)]^{1+\theta ^{}}`$ for any $`\eta <1`$ and any $`\theta ^{}>0`$.
### 1.3. Main result
We begin by defining the scale function of the almost-sure asymptotics:
$$\frac{b_t}{\alpha _{b_t}^2}=\mathrm{log}G(t),t>0.$$
(1.7)
In other words, $`tb_t`$ is the inverse of the function $`tt\alpha _t^2`$ (which we may and shall assume to be strictly increasing), evaluated at $`\mathrm{log}G(t)`$. Note that, since $`lim_{\mathrm{}\mathrm{}}G(\mathrm{})=0`$, we have $`b_t\mathrm{}`$ as $`t\mathrm{}`$. If $`G(\mathrm{})=\mathrm{}^{\zeta +o(1)}`$ as $`\mathrm{}\mathrm{}`$ for some $`\zeta (0,1]`$, then $`\alpha _{b_t}^2=\zeta ^\beta (\mathrm{log}t)^{\beta +o(1)}`$, where $`\beta =2\nu /(12\nu )=2(1\gamma )/(13\gamma )(0,2]`$. In the case $`\zeta =0`$, $`\alpha _{b_t}=o(\mathrm{log}^\beta t)`$ as $`t\mathrm{}`$.
Here is our main result. The constant $`\stackrel{~}{\chi }`$ appearing in (1.8) depends only on $`A`$, $`\gamma `$ and $`\kappa `$ and will be defined in Subsection 2.1.
###### Theorem 1.1
Let $`d=1`$ and suppose that Assumption (H) and Assumption (G) hold. Define $`tb_t`$ as in (1.7), and let $`\stackrel{~}{\chi }`$ be the constant in Theorem 1.5 of \[BK00\]. Then
$$\underset{t\mathrm{}}{lim}\frac{\alpha _{b_t}^2}{t}\mathrm{log}u(t,0)=\stackrel{~}{\chi },\mathrm{Prob}\text{-almost surely.}$$
(1.8)
Interestingly, if $`\mathrm{}G(\mathrm{})`$ has a power-law decay as $`\mathrm{}\mathrm{}`$, the lower-tail dependence of the rate can explicitly be computed. This allows for an easy comparison with the assertion in Theorem 1.5 of \[BK00\]. Let $`tb_t^{}`$ be the scale function introduced in \[BK00\]:
$$\frac{b_t^{}}{\alpha _{b_t^{}}^2}=\mathrm{log}t.$$
(1.9)
Recall $`\zeta _{}=sup\{\delta 0:[\mathrm{log}(\xi (0)1)]^\delta <\mathrm{}\}`$. For $`G`$ decaying with a power law, necessarily, $`G(\mathrm{})=\mathrm{}^{\zeta _{}+o(1)}`$. The following is an immediate consequence of Theorem 1.1 and the regularity of $`t\alpha _t`$:
###### Corollary 1.2
Let $`d=1`$, suppose $`\mathrm{Prob}(\xi (0)=\mathrm{})=0`$ and suppose that Assumption (H) holds. Assume that either $`\zeta _{}=1`$ or $`\zeta _{}(0,1)`$ and $`G(\mathrm{})=\mathrm{}^{\zeta _{}+o(1)}`$ as $`\mathrm{}\mathrm{}`$. Then
$$\underset{t\mathrm{}}{lim}\frac{\alpha _{b_t^{}}^2}{t}\mathrm{log}u(t,0)=\zeta _{}^\beta \stackrel{~}{\chi },\mathrm{Prob}\text{-almost surely.}$$
(1.10)
where $`\beta =2(1\gamma )/(13\gamma )`$.
###### Remark 3.
By comparison of Corollary 1.2 and Theorem 1.5 of \[BK00\], $`\zeta _{}=1`$ is necessary and sufficient for the assertion of the latter to hold, at least in the class of distribution with $`G`$ decaying as a positive power. In particular, the condition that $`\mathrm{log}(\xi (0)1)<\mathrm{}`$ in \[BK00\] is only marginally non-optimal because Theorem 1.5 also literally holds if we just assume that $`[\mathrm{log}(\xi (0)1)]^\delta `$ be integrable for any $`\delta <1`$. This answers the first of the questions above.
###### Remark 4.
The cases with $`\zeta _{}>0`$ have a different absolute size of the rate while the time dependence remains as for $`\zeta _{}=1`$. However, when $`\zeta _{}=0`$, also the time dependence changes. For instance, in the aforementioned example $`\mathrm{Prob}(\mathrm{log}(\xi (0)1)dx)C/[x\mathrm{log}^{1+\theta }(x)]dx`$ as $`x\mathrm{}`$ (see Remark 2), $`\alpha _{b_t}^2=[\mathrm{log}\mathrm{log}t]^{\beta +o(1)}`$, which grows much slower than in the case $`\zeta _{}>0`$. For yet thicker lower tails, even slower growths are possible. We conclude that the result of Theorem 1.5 of \[BK00\] qualitatively changes only when $`\xi (0)`$ lacks all positive logarithmic moments.
The remainder of this paper is organized as follows: In the next section we define some important objects and use them to give a heuristic outline of the proof. The actual proof comes in Section 3. Since many steps can almost literally be taken over from \[BK00\], we stay as terse as possible. The essentially novel part are Lemmas 3.2, 3.3, and 3.4.
## 2. Definitions and heuristics
### 2.1. Auxiliary objects
For the sake of both completeness and later reference, we will now introduce the objects needed to define the quantity $`\stackrel{~}{\chi }`$ in Theorem 1.1. Then we proceed by recalling the Feynman-Kac representation and some formulas for Dirichlet eigenvalues.
#### 2.1.1. Definition of $`\stackrel{~}{\chi }`$
Let $`_R`$ be the set of continuous functions $`f:[0,\mathrm{})`$ satisfying $`\mathrm{supp}f[R,R]`$ and having total integral equal to one. Let $`C^+(R)`$ (resp., $`C^{}(R)`$) be the set of continuous functions $`[R,R][0,\mathrm{})`$ (resp. $`[R,R](\mathrm{},0]`$). For $`H`$ in the $`\gamma `$-class, let $`_R:C^+(R)(\mathrm{},0]`$ be the functional defined by
$$_R(f)=A_{[R,R]}f^\gamma \mathsf{\hspace{0.17em}1}_{\{f>0\}}𝑑x,$$
(2.1)
where $`A`$ is as in (1.4).
Let $`_R:C^{}(R)[0,\mathrm{}]`$ be the Legendre transform of $`_R`$:
$$_R(\psi )=sup\{(f,\psi )_R(f):fC^+(R),\mathrm{supp}f\mathrm{supp}\psi \},$$
(2.2)
where $`(f,\psi )=f(x)\psi (x)𝑑x`$. Conventionally, $`_R(0)=\mathrm{}`$. If $`H`$ is in the $`\gamma `$-class with a $`\gamma [0,1)`$, $`_R(\psi )`$ can explicitly be computed: for any $`\psi C^{}(R)`$, $`\psi 0`$,
$$_R(\psi )=\{\begin{array}{cc}(1\gamma ^1)(A\gamma )^{\frac{1}{1\gamma }}|\psi (x)|^{\frac{\gamma }{1\gamma }}𝑑x,\hfill & \text{if }\gamma (0,1),\hfill \\ A|\mathrm{supp}\psi |,\hfill & \text{if }\gamma =0,\hfill \end{array}$$
(2.3)
where $`|\mathrm{supp}\psi |`$ is the Lebesgue measure of $`\mathrm{supp}\psi `$. (Here $`_R(\psi )=\mathrm{}`$ whenever $`\gamma (0,1)`$ and the integral diverges.)
The last object we need is the principal eigenvalue of the operator $`\kappa \mathrm{\Delta }+\psi `$ on $`L^2([R,R])`$ with Dirichlet boundary conditions:
$$\lambda _R(\psi )=sup\{(\psi ,g^2)\kappa g_2^2:gC_\mathrm{c}^{\mathrm{}}(\mathrm{supp}\psi ,),g_2=1\},$$
(2.4)
with the interpretation $`\lambda _R(0)=\mathrm{}`$. Then
$$\stackrel{~}{\chi }=\underset{R>0}{sup}sup\{\lambda _R(\psi ):\psi C^{}(R),_R(\psi )1\}.$$
(2.5)
As was proved in \[BK00\], $`\stackrel{~}{\chi }(0,\mathrm{})`$.
###### Remark 5.
In $`d=1`$, the minimizer of an associated variational problem (namely, that for the annealed or moment asymptotics) can explicitly be computed, see \[BK98\]. Proposition 1.4 of \[BK00\] then allows $`\stackrel{~}{\chi }`$ to be evaluated in a closed form. Except for $`\gamma =0`$, no such expression is known in higher dimensions.
#### 2.1.2. Feynman-Kac formula, Dirichlet eigenvalues
Let $`(X(s))_{s[0,\mathrm{})}`$ be the continuous-time simple random walk on $``$ with generator $`\kappa \mathrm{\Delta }^\mathrm{d}`$. We use $`𝔼_x`$ to denote the expectation with respect to the walk starting at $`x`$. The Feynman-Kac representation for $`u(t,)`$ then reads
$$u(t,x)=𝔼_x\left[\mathrm{exp}\left\{_0^t\xi \left(X(s)\right)𝑑s\right\}\right].$$
(2.6)
Given $`R>0`$, let $`Q_R=[R,R]`$ and let $`u_R(t,x)`$ be the solution to the system (1.1) in $`Q_R`$ and Dirichlet boundary condition $`u_R(,x)=0`$ for $`xQ_R`$. Let $`\tau _R`$ be the first exit time from $`Q_R`$, i.e., $`\tau _R=inf\{s>0:X(s)Q_R\}`$. Then
$$u_R(t,x)=𝔼_x\left[\mathrm{exp}\left\{_0^t\xi \left(X(s)\right)𝑑s\right\}\mathrm{𝟣}\{\tau _R>t\}\right]$$
(2.7)
Note that $`Ru_R(t,x)`$ is increasing.
In the forthcoming developments we will also need the principal Dirichlet eigenvalue of the operator $`\kappa \mathrm{\Delta }^\mathrm{d}+\xi `$ in the box $`z+Q_R`$ centered at $`z`$:
$$\lambda _{z;R}^\mathrm{d}(\xi )=sup\{\underset{xQ_R}{}\xi (x+z)g(x)^2+\kappa \underset{xQ_R}{}g(x)[\mathrm{\Delta }^\mathrm{d}g](x):g\mathrm{}^2(Q_R),g_2=1\}.$$
(2.8)
Note that, by the standard eigenvalue expansion (see \[BK00\]),
$$\mathrm{e}_R(z)^2e^{t\lambda _{z;R}^\mathrm{d}(\xi )}u_R(t,z)\mathrm{\#}Q_Re^{t\lambda _{z;R}^\mathrm{d}(\xi )},$$
(2.9)
where $`\mathrm{e}_R()`$ is the $`\mathrm{}^2`$-normalized principal eigenvector in $`Q_R`$. In particular, the logarithmic asymptotics of $`u_R(t,z)`$ and the asymptotics of $`t\lambda _{z;R}^\mathrm{d}(\xi )`$ coincide provided $`R=R(t)`$ does not grow too fast with $`t`$ (which ensures that $`t\mathrm{e}_{R(t)}(z)^2`$ does not decay too fast).
### 2.2. Heuristic explanation
As alluded to in the introduction, (1.8) results from the competition of two mechanisms: (1) searching for optimal shapes of the potential by the walk in (2.6) and (2) screening off far away sites by regions of strongly negative potential. Let us describe this interplay in detail. To avoid cluttering of indices we often use $`\alpha (b_t)`$ in the place of $`\alpha _{b_t}`$.
Consider a “macrobox” $`Q_{r(t)}=[r(t),r(t)]`$ with $`r(t)\mathrm{exp}[b_t\alpha (b_t)^2]`$, where we think of $`b_t`$ as of a yet undetermined scale function. Fix $`R>0`$ and a shape function $`\psi C^{}(R)`$ satisfying $`_R(\psi )<1`$. A Borel-Cantelli argument shows that there exists a randomly located microbox in $`Q_{r(t)}`$, with diameter $`2R\alpha (b_t)`$, where $`\xi `$ is shaped like $`\psi _t()\psi (/\alpha (b_t))/\alpha (b_t)^2`$. Let us assume that $`R`$ and $`\psi `$ approximately maximize (2.5), i.e., $`\lambda _R(\psi )\stackrel{~}{\chi }`$. Then the dominating strategy for the walk is to move in a short time to that favorable microbox and spend the rest of the time until $`t`$ in it. The contribution coming from the long stay in the microbox is roughly $`\mathrm{exp}[t\lambda _{R\alpha (b_t)}(\psi _t)]`$, which can be approximated by $`\mathrm{exp}[t\alpha (b_t)^2\lambda _R(\psi )]\mathrm{exp}[t\alpha (b_t)^2\stackrel{~}{\chi }]`$, using the scaling properties of the Laplace operator.
The size of the macrobox is determined by the amount of mass the walk loses on the way from the origin to the favorable microbox, while traveling through long stretches of large negative potential. A calculation shows that the penalty it pays is roughly of order $`\mathrm{exp}[_{x=1}^{r(t)}\mathrm{log}(\xi (x)1)]`$. (An optimal strategy is not to spend more than $`(\xi (x)1)^1`$ time units at each site $`x`$ on the way.) Under our assumptions on the lower tails of $`\xi (0)`$, a Borel-Cantelli argument shows that this penalty is roughly $`\mathrm{exp}[G^1(1/r(t))]`$ where $`G^1`$ denotes the inverse function of $`G`$.
As it turns out, the two mechanisms run at optimal “speed” when the two exponents are roughly of the same order, i.e., $`G^1(1/r(t))t\alpha (b_t)^2t`$, because $`\alpha _{b_t}t`$. Recalling that $`r(t)\mathrm{exp}[b_t\alpha (b_t)^2]`$, this reasoning leads to (1.7). A fine tuning of $`r(t)`$ makes the contribution from the travel to the microbox negligible compared to the contribution from the stay in it, i.e., we shall in fact have $`G^1(1/r(t))=o(t\alpha (b_t)^2)`$. Hence, we obtain (1.8) with $`\stackrel{~}{\chi }`$ as in (2.5).
## 3. Proof of Theorem 1.1
As in \[BK00\], the main result will be proved by separately proving upper and lower bounds in (1.8). The proof of Corollary 1.2 comes at the very end of this section.
### 3.1. The upper bound
Recall the notation of Subsection 2.1, in particular that $`Q_R=[R,R]`$. Let
$$r(t)=\frac{3}{G(t)}\mathrm{log}G(t).$$
(3.1)
Note that $`r(t)=t^{\zeta +o(1)}`$ as $`t\mathrm{}`$ if $`\stackrel{~}{G}(\mathrm{})=\mathrm{}^{\zeta +o(1)}`$ as $`\mathrm{}\mathrm{}`$. Abbreviate $`B_R(t)=Q_{r(t)+2R}`$.
The crux of the proof of the upper bound in Theorem 1.1 is the following generalization of Proposition 4.4 of \[BK00\] adapted to the new definition of $`r(t)`$.
###### Proposition 3.1
There exists a constant $`C=C(\kappa )>0`$ and a random variable $`C_\xi (0,\mathrm{})`$ such that, $`\mathrm{Prob}`$-almost surely, for all $`R,t>C`$,
$$u(t,0)C_\xi e^t+e^{Ct/R^2}3r(t)\mathrm{exp}\left\{t\underset{zB_R(t)}{\mathrm{max}}\lambda _{z;2R}^\mathrm{d}(\xi )\right\}.$$
(3.2)
Proof of Theorem 1.1, upper bound. With Proposition 3.1 in the hand, the proof goes along very much the same lines as in \[BK00\]. Indeed, let $`r(t)`$ be as in (3.1) and set $`R`$ in (3.2) to be $`R\alpha (Kb_t)`$, where $`K>0`$ will be chosen later and $`R`$ will tend to $`\mathrm{}`$. Let $`H`$ be in the $`\gamma `$-class and recall that $`\alpha _t=t^{\nu +o(1)}`$ with $`\nu =(1\gamma )/(3\gamma )`$.
Abbreviate $`B(t)=B_{R\alpha (Kb_t)}(t)`$ and $`\lambda (z)=\lambda _{z;2R\alpha (Kb_t)}^\mathrm{d}(\xi )`$, and note that $`r(t)e^{o(t\alpha (b_t)^2)}`$. Then, using also that $`lim_t\mathrm{}\alpha (Kb_t)/\alpha (b_t)=K^\nu `$, we have from (3.2) that
$$\underset{t\mathrm{}}{lim\; sup}\frac{\alpha _{b_t}^2}{t}\mathrm{log}u(t,0)\frac{C}{K^{2\nu }R^2}+\underset{t\mathrm{}}{lim\; sup}\left[\alpha _{b_t}^2\underset{zB(t)}{\mathrm{max}}\lambda (z)\right],$$
(3.3)
$`\mathrm{Prob}`$-almost surely. Abbreviating $`M(t)=\mathrm{max}_{zB(t)}\lambda (z)`$, we have to prove that, for any $`\epsilon >0`$,
$$\underset{t\mathrm{}}{lim\; sup}\alpha _{b_t}^2M(t)\stackrel{~}{\chi }+\epsilon ,\mathrm{Prob}\text{-almost surely},$$
(3.4)
for some appropriate $`K(0,\mathrm{})`$ and sufficiently large $`R`$.
Note that the eigenvalues $`\lambda (z)`$ have identical distribution. Furthermore, their exponential moments can be estimated by
$$\underset{R\mathrm{}}{lim\; sup}\underset{t\mathrm{}}{lim\; sup}\frac{\alpha _{b_t}^2}{b_t}\mathrm{log}e^{Kb_t\lambda (z)}K^{12\nu }\chi ,$$
(3.5)
where $`\chi (0,\mathrm{})`$ is a constant related to $`\stackrel{~}{\chi }`$, see \[BK00\]. Since $`tM(t)`$ is increasing and $`t\alpha _{b_t}`$ slowly varying, it suffices to prove (3.4) for $`t`$ taking only a discrete set of values; the main difference compared to \[BK00\] is that now we take
$$\frac{1}{G(t)}\{e^n:n\}.$$
(3.6)
Let $`G(t)=e^n`$ and note that (1.7) implies that $`b_t\alpha _{b_t}^2=n`$. The proof now proceeds exactly as in \[BK00\]: We let $`p_n(\epsilon )=\mathrm{Prob}(M(t)\alpha _{b_t}^2\stackrel{~}{\chi }+\epsilon )`$ and use the Chebyshev inequality and (3.5) to derive that $`p_n(\epsilon )`$ is summable on $`n`$ for all $`\epsilon >0`$, provided $`K`$ is chosen appropriately and $`R`$ is sufficiently large. The claim is finished using the Borel-Cantelli lemma. ∎
It remains to prove Proposition 3.1. In \[BK00\], the choice $`t\mathrm{log}t`$ for $`r(t)`$ allowed us to use a simple probability estimate for the simple random walk; in particular, the corresponding bound (3.2) held true uniformly in all non-positive potentials. In our present cases, $`r(t)`$ is typically much smaller than $`t\mathrm{log}t`$ and the potential has to cooperate to get the bound (3.2). Unlike in \[BK00\], the role of the potential is actually dominant in the cases of our present interest.
###### Lemma 3.2
For any $`b(2\kappa ,\mathrm{})`$ there is a random variable $`C(\xi )(0,\mathrm{})`$ such that, $`\mathrm{Prob}`$-almost surely,
$$u(t,0)u_R(t,0)C(\xi )\left(\underset{x=0}{\overset{R}{}}\frac{b}{\xi (x)b}+\underset{x=R}{\overset{0}{}}\frac{b}{\xi (x)b}\right),R,t0.$$
(3.7)
Proof. Let $`(X_k)_{k_0}`$ be the embedded discrete-time simple random walk on $``$ and let $`\mathrm{}_n(x)`$ be its local times defined by $`\mathrm{}_n(x)=_{k=1}^n\mathrm{𝟣}\{X_k=x\}`$. Let $`𝔼_y^\mathrm{d}`$ denote the expectation with respect to the discrete-time walk, starting at $`y`$. Abbreviate $`\xi _k=\xi (X_k)`$ and $`\widehat{u}_R(t,0)=u(t,0)u_R(t,0)`$. Then, by (2.6) and (2.7),
$$\widehat{u}_R(t,0)=e^{2\kappa t}\underset{nR}{}(2\kappa )^n𝔼_0^\mathrm{d}\left[_{\mathrm{}_n(t)}𝑑t_1\mathrm{}𝑑t_n\mathrm{exp}\left\{\underset{k=0}{\overset{n}{}}\xi _kt_k\right\}\mathrm{𝟣}\{\mathrm{supp}\mathrm{}_nQ_R\}\right],$$
(3.8)
where $`\mathrm{}_n(t)=\{(t_1,\mathrm{},t_n)(0,\mathrm{})^n:t_1+\mathrm{}+t_nt\}`$, and $`t_0`$ is a shorthand for $`t(t_1+\mathrm{}+t_n)`$.
Fix $`b>2\kappa `$ and define
$$𝒜_n=\{x\mathrm{supp}\mathrm{}_n:\xi (x)b\}.$$
(3.9)
Let
$$_n=\{k\{1,\mathrm{},n\}:X_k𝒜_n\}$$
(3.10)
be the set of all the times at which the walk visits a point $`x`$ with $`\xi (x)>b`$.
By relaxing the constraint $`t_1+\mathrm{}+t_nt`$ in $`\mathrm{}_n(t)`$ to $`t_kt`$ for every $`k_n`$, neglecting the terms with $`k_n\{0\}`$ in the exponential, and integrating out $`t_1,\mathrm{},t_n`$, we get
$$\widehat{u}_R(t,0)e^{2\kappa t}\underset{nR}{}\underset{m=0}{\overset{n}{}}\frac{(2\kappa t)^m}{m!}𝔼_0^\mathrm{d}\left[\mathrm{𝟣}\left\{\mathrm{\#}_n=m\right\}\mathsf{\hspace{0.17em}1}\{\mathrm{supp}\mathrm{}_nQ_R\}\underset{0<kn:k_n}{}\frac{2\kappa }{\xi _k}\right].$$
(3.11)
Neglecting the first indicator and the restriction to $`mn`$, we can carry out the sum over $`m`$ in (3.11) and find that
$$\widehat{u}_R(t,0)\underset{nR}{}𝔼_0^\mathrm{d}\left[\mathrm{𝟣}\{\mathrm{supp}\mathrm{}_nQ_R\}\underset{x𝒜_n}{}\left(\frac{2\kappa }{\xi (x)}\right)^{\mathrm{}_n(x)}\right].$$
(3.12)
On $`\{\mathrm{supp}\mathrm{}_nQ_R\}`$, the walk visits either all sites in $`\{0,\mathrm{},R\}`$ or all sites in $`\{R,\mathrm{},0\}`$. Hence, we can estimate
$$\mathrm{𝟣}\{\mathrm{supp}\mathrm{}_nQ_R\}\underset{x𝒜_n}{}\frac{2\kappa }{\xi (x)}\underset{x=1}{\overset{R}{}}\frac{b}{\xi (x)b}+\underset{x=R}{\overset{1}{}}\frac{b}{\xi (x)b}.$$
(3.13)
The claim (3.7) then follows from the assertion
$$\underset{n=1}{\overset{\mathrm{}}{}}𝔼_0^\mathrm{d}\left[\underset{x𝒜_n}{}\left(\frac{2\kappa }{b}\right)^{\mathrm{}_n(x)1}\right]<\mathrm{}\mathrm{Prob}\text{-almost surely,}$$
(3.14)
where we used that $`\xi (x)b`$ whenever $`x𝒜_n`$. (The term with $`x=0`$ in (3.7) can be added or removed at the cost of changing $`C(\xi )`$ by a finite amount.)
Let us prove that (3.14) holds. First we note that $`𝒜_n`$ contains in every sufficiently large interval in $``$ a positive fraction of sites. Indeed, put $`p=\mathrm{Prob}(\xi (0)>b)(0,1]`$ and note that by Cramér’s theorem we have $`\mathrm{Prob}(\mathrm{\#}(𝒜_nI)\frac{p}{2}\mathrm{\#}I)e^{c\mathrm{\#}I}`$ for every bounded interval $`I`$ and some $`c>0`$ independent of $`I`$. A routine application of the Borel-Cantelli lemma implies that
$$\text{interval }I[n,n]:\mathrm{\#}In^{1/4}\mathrm{\#}(𝒜_nI)>\frac{p}{2}\mathrm{\#}I,$$
(3.15)
for $`n`$ large enough, $`\mathrm{Prob}`$-almost surely.
Now we prove that with high probability, there are sufficiently large intervals which are traversed from one end to the other at least twice by the random walk $`(X_k)_{k=0,\mathrm{},n}`$. Fix $`K_n=3\mathrm{log}n`$ and abbreviate $`k_n=n/K_n`$. We divide the walk into $`K_n`$ pieces $`(X_k^{(i)})_{k=0,\mathrm{},k_n}`$ (neglecting a small overshoot) with $`X_k^{(i)}=X_{(i1)k_n+k}X_{(i1)k_n}`$ for $`i=1,\mathrm{},K_n`$. Note that these $`K_n`$ walks are independent copies of each other. Let us introduce the events
$$B_n=\underset{i=1}{\overset{K_n1}{}}\left\{\mathrm{sgn}X_{k_n}^{(i)}=\mathrm{sgn}X_{k_n}^{(i+1)}\right\}\text{and}C_n=\underset{i=1}{\overset{K_n}{}}\left\{\underset{1kk_n}{\mathrm{max}}\left|X_k^{(i)}\right|L_n\right\},$$
(3.16)
where $`L_n=\sqrt{k_n/(\eta \mathrm{log}n)}`$. It is elementary that $`_0^\mathrm{d}(B_n)2^{K_n+1}n^{2+o(1)}`$ as $`n\mathrm{}`$. Furthermore, with the help of a concatenation argument and convergence of simple random walk to Brownian motion we derive that $`_0^\mathrm{d}(C_n)n^{2+o(1)}`$, whenever $`\eta >0`$ is large enough. Now we estimate
$$𝔼_0^\mathrm{d}\left[\underset{x𝒜_n}{}\left(\frac{2\kappa }{b}\right)^{\mathrm{}_n(x)1}\right]_0^\mathrm{d}(B_n)+_0^\mathrm{d}(C_n)+𝔼_0^\mathrm{d}\left[\mathsf{\hspace{0.17em}1}_{B_n^\mathrm{c}C_n^\mathrm{c}}\underset{x𝒜_n}{}\left(\frac{2\kappa }{b}\right)^{\mathrm{}_n(x)1}\right].$$
(3.17)
Note that, on $`B_n^\mathrm{c}C_n^\mathrm{c}`$, there is an interval $`I[n,n]`$ with $`\mathrm{\#}IL_n`$ such that every point of $`I`$ is visited by at least two of the subwalks, i.e., we have $`\mathrm{}_n(x)2`$ for any $`xI`$. If $`n`$ is sufficiently large, we deduce from (3.15) that there are at least $`pL_n/2`$ points $`x`$ with $`\mathrm{}_n(x)2`$. By using this in (3.17), we have
$$𝔼_0^\mathrm{d}\left[\underset{x𝒜_n}{}\left(\frac{2\kappa }{b}\right)^{\mathrm{}_n(x)1}\right]n^{2+o(1)}+\left(\frac{2\kappa }{b}\right)^{L_np/2},n\mathrm{}.$$
(3.18)
The right hand side is clearly summable on $`n`$ since $`2\kappa /b<1`$. This finishes the proof. ∎
Our next task is to get a good estimate on the size of the products in (3.7).
###### Lemma 3.3
Suppose that $`\mathrm{log}(\xi (0)1)=\mathrm{}`$. Then, for all $`b1`$,
$$\underset{n\mathrm{}}{lim}\frac{1}{G^1(1/n)}\underset{x=1}{\overset{2n\mathrm{log}n}{}}\mathrm{log}\left(\frac{\xi (x)b}{b}\right)=\mathrm{}\mathrm{Prob}\text{-almost surely.}$$
(3.19)
Proof. Abbreviate $`N_n=2n\mathrm{log}n`$ and let $`b1`$. Then
$$\mathrm{log}\left(\frac{\xi (x)b}{b}\right)\mathrm{log}(\xi (x)1)\mathrm{log}b.$$
(3.20)
Using this estimate and the Chebyshev inequality, we have for any $`\theta >0`$ that
$$\mathrm{Prob}\left(\underset{x=1}{\overset{N_n}{}}\mathrm{log}\left(\frac{\xi (x)b}{b}\right)\theta G^1(1/n)\right)\mathrm{exp}\left\{N_nG(1/\lambda )+N_n\lambda \mathrm{log}b+\lambda \theta G^1(1/n)\right\},$$
(3.21)
for any $`\lambda >0`$. Set $`\lambda =1/G^1(1/n)`$ and note that we have $`G(1/\lambda )/\lambda \mathrm{}`$ as $`\lambda 0`$, due to $`\mathrm{log}(\xi (0)1)=\mathrm{}`$. Consequently, the term with $`\mathrm{log}b`$ is negligible and the right-hand side of (3.21) is bounded by $`n^{2+o(1)}`$. The claim is finished by the Borel-Cantelli lemma. ∎
Proof of Proposition 3.1. Pick any $`b(2\kappa ,\mathrm{})`$. Let $`t_0`$ be so large such that the sum in (3.19) with $`n=1/G(t)`$ for all $`tt_0`$ exceeds $`G^1(1/n)`$. Note that $`r(t)2n\mathrm{log}n`$. Combining the results of Lemma 3.2 for $`R=r(t)`$ and Lemma 3.3, we derive, for sufficiently large $`n`$ resp. $`t`$, the bound
$$u(t,0)u_{r(t)}(t,0)2C(\xi )\mathrm{exp}\left(G^1(1/n)\right),$$
(3.22)
where $`C(\xi )`$ is the constant from (3.7). But $`G^1(1/n)t`$ by our choice of $`n`$, which means that $`u(t,0)u_{r(t)}(t,0)C_\xi e^t`$, where $`C_\xi =2C(\xi )e^{t_0}`$. The rest of the argument does not involve the particular form of $`r(t)`$ and can directly be taken over from \[BK00\]. ∎
### 3.2. The lower bound
Unlike the upper bound, the lower bound was basically proved already in \[BK00\], up to a change of the spatial scale and Lemma 3.4 below. For this reason, we shall only indicate the necessary changes.
First we prove the following converse of Lemma 3.3:
###### Lemma 3.4
Fix $`\eta (0,1)`$ and let $`\stackrel{~}{G}_\eta `$ satisfy (ii) and (iii) in Assumption (G). Then there exists a $`\varrho (0,\mathrm{})`$ such that
$$\underset{n\mathrm{}}{lim\; sup}\frac{1}{\stackrel{~}{G}_\eta ^1(\varrho /n)}\underset{x=1}{\overset{n}{}}\mathrm{log}\left(\xi (x)1\right)1\mathrm{Prob}\text{-almost surely.}$$
(3.23)
Proof. The argument is based on the asymptotic sublinearity of $`1/\stackrel{~}{G}_\eta `$ at infinity. However, in order to have sublinearity on the whole interval $`(0,\mathrm{})`$, we first construct an auxiliary modification of $`\stackrel{~}{G}_\eta `$.
Let $`x_0>0`$ be such that $`1/\stackrel{~}{G}_\eta `$ is positive, increasing, and concave on $`[x_0,\mathrm{})`$. Let $`D_0`$ to be the right derivative of $`1/\stackrel{~}{G}_\eta `$ at $`x_0`$. Define $`\widehat{G}_\eta :(0,\mathrm{})(0,\mathrm{})`$ by the formula
$$1/\widehat{G}_\eta (x)=\{\begin{array}{cc}D_0x\hfill & \text{if }xx_0,\hfill \\ 1/\stackrel{~}{G}_\eta (x)+D_0x_01/\stackrel{~}{G}_\eta (x_0)\hfill & \text{if }x>x_0.\hfill \end{array}$$
(3.24)
Note that $`1/\widehat{G}_\eta `$ is positive, increasing, concave and hence sublinear on $`(0,\mathrm{})`$. Moreover, Assumption (G)(iii) holds true for $`\stackrel{~}{G}_\eta `$ replaced by $`\widehat{G}_\eta `$.
For $`a1`$, abbreviate $`Y_a(x)=\mathrm{log}(\xi (x)a)`$. Choose $`a=e^{x_0}`$ and estimate, for $`n\mathrm{}`$,
$$\frac{1}{\stackrel{~}{G}_\eta \left(_{x=1}^nY_a(x)\right)}\frac{1+o(1)}{\widehat{G}_\eta \left(_{x=1}^nY_a(x)\right)}(1+o(1))\underset{x=1}{\overset{n}{}}\frac{1}{\widehat{G}_\eta (Y_a(x))},$$
(3.25)
where we used the fact that $`_{x=1}^nY_a(x)\mathrm{}`$ almost surely, and sublinearity of $`1/\widehat{G}_\eta `$. Since we have that $`1/\widehat{G}_\eta (Y_a(x))<\mathrm{}`$, the Strong Law of Large Numbers tells us that the right-hand side of (3.25) is almost surely no more than $`\varrho n`$, where for $`\varrho `$ we can take, for instance,
$$\varrho =21/\widehat{G}_\eta (Y_a(0)).$$
(3.26)
Hence, we derive
$$\underset{x=1}{\overset{n}{}}Y_1(x)\underset{x=1}{\overset{n}{}}Y_a(x)\stackrel{~}{G}_\eta ^1(\varrho /n),$$
(3.27)
which directly yields the desired claim. ∎
Another important ingredient is the following adaptation of the crucial Proposition 5.1 of \[BK00\] to the present situation. For $`\eta (0,1)`$, choose $`\varrho `$ as in Lemma 3.4 and let this time
$$\gamma _t=\frac{\varrho }{\stackrel{~}{G}_\eta (t\alpha _{b_t}^3)}$$
(3.28)
be the size of the macrobox $`Q_{\gamma _t}`$ (see Subsection 2.2). Note that $`t^{\eta \zeta +o(1)}\gamma _tt^{\zeta +o(1)}`$ as $`t\mathrm{}`$ if $`G(\mathrm{})=\mathrm{}^{\zeta +o(1)}`$ as $`\mathrm{}\mathrm{}`$. Suppose without loss of generality that $`t\gamma _t`$ is increasing.
Define for each $`\psi C^{}(R)`$ a “microbox”
$$Q^{(t)}=\{\begin{array}{cc}Q_{R\alpha (b_t)}\hfill & \text{if }\gamma 0,\hfill \\ Q_{R\alpha (b_t)}\mathrm{supp}\psi _t\hfill & \text{if }\gamma =0,\hfill \end{array}$$
(3.29)
where $`\psi _t:(\mathrm{},0]`$ is the function $`\psi _t()=\psi (/\alpha _{b_t})/\alpha _{b_t}^2`$. The crucial input for the lower bound is the following claim, which says that, with probability one provided $`_R(\psi )<1`$ and $`t`$ is large, there is at least one microbox $`Q^{(t)}`$ in $`Q_{\gamma _t}`$, where $`\xi `$ is no less than (the accordingly shifted) $`\psi _t`$.
###### Proposition 3.5
Let $`R>0`$ and fix $`\psi C^{}(R)`$ satisfying $`_R(\psi )<1`$. Let $`\epsilon >0`$ and suppose Assumptions (G) and (H) hold. Then the following holds almost surely: For each $`\eta (_R(\psi ),1)`$, there is a $`t_0=t_0(\xi ,\psi ,\epsilon ,R,\eta )<\mathrm{}`$ such that for each $`tt_0`$, there is a $`y_tQ_{\gamma _t}`$ with
$$\xi (z+y_t)\psi _t(z)\epsilon \alpha _{b_t}^2zQ^{(t)}.$$
(3.30)
Proof. We begin by formalizing the event in (3.30); in order to later approximate continuous $`t`$ by a discrete variable, we write $`\epsilon /2`$ instead of $`\epsilon `$:
$$A_y^{(t)}=\underset{zQ^{(t)}}{}\left\{\xi (y+z)\psi _t(z)\frac{\epsilon }{2\alpha (b_t)^2}\right\}.$$
(3.31)
Note that the probability of $`A_y^{(t)}`$ does not depend on $`y`$ and note that different $`A_y^{(t)}`$’s are independent if the $`y`$’s have distance larger than $`3R\alpha (b_t)`$ from each other. The proof of Lemma 5.5 in \[BK00\] shows that $`\mathrm{Prob}(A_0^{(t)})G(t)^{_R(\psi )+o(1)}`$ as $`t\mathrm{}`$ (the only modification required is to replace every occurrence of $`t`$ in the meaning $`\mathrm{exp}\{b_t\alpha (b_t)^2\}`$ by $`1/G(t)`$).
In order to prove our claim, it is sufficient to show the summability of
$$p_t=\mathrm{Prob}\left(\underset{yQ_{\gamma _t}}{}\left(A_y^{(t)}\right)^\mathrm{c}\right)$$
(3.32)
over all $`t>0`$ such that $`1/G(t)\{e^n:n\}`$. (The sufficiency follows from the facts that $`\alpha (b_t)/\alpha (b_{et})1`$ as $`t\mathrm{}`$ and that $`tb_t`$ and $`t\gamma _t`$ are increasing. The error terms are absorbed into an extra $`\epsilon /2`$ in (3.30) compared to (3.31), see \[BK00\].)
Using the independence of $`A_y^{(t)}`$ for $`yB(t)=Q_{\gamma _t}3R\alpha (b_t)`$ and the bound $`\mathrm{Prob}(A_0^{(t)})G(t)^{_R(\psi )+o(1)}`$, we easily derive
$$p_t\left(1G(t)^{_R(\psi )+o(1)}\right)^{\mathrm{\#}B(t)}\mathrm{exp}\left\{\frac{G(t)^{_R(\psi )+o(1)}}{\alpha (b_t)\stackrel{~}{G}_\eta (t\alpha (b_t)^3)}\right\},$$
(3.33)
where we used that $`\mathrm{\#}B(t)2\gamma _t/(3R\alpha (b_t))`$ and then applied the definition of $`\gamma _t`$.
Use concavity of $`1/\stackrel{~}{G}_\eta `$ to estimate $`1/\stackrel{~}{G}_\eta (t\alpha (b_t)^3)\alpha (b_t)^3/\stackrel{~}{G}_\eta (t)`$ and use Assumption (G)(i) to bound $`\stackrel{~}{G}_\eta (t)`$ by $`G(t)^{\eta +o(1)}`$. Furthermore, since $`\alpha (b_t)`$ is bounded from above by a positive power of $`b_t\alpha (b_t)^2`$, we see from (1.7) that $`\alpha (b_t)=G(t)^{o(1)}`$. Applying all this reasoning on the right-hand side of (3.33), we see that $`p_t\mathrm{exp}(G(t)^{_R(\psi )\eta +o(1)})`$ as $`t\mathrm{}`$, which is summable on the sequence of $`t`$ such that $`1/G(t)\{e^n:n\}`$. This finishes the proof. ∎
Now we finish the proof of our main result.
Proof of Theorem 1.1, lower bound. Let $`\epsilon >0`$ and fix $`R>0`$ and $`\psi C^{}(R)`$ such that $`_R(\psi )<1`$. Let $`\eta (_R(\psi ),1)`$, define $`\gamma _t`$ as in (3.28) and let $`y_t`$ be as in Proposition 3.5; suppose $`y_t0`$ without loss of generality. Let $`r_x=[\xi (x)1]^1`$. As in \[BK00\], the lower bound will be obtained by restricting the walk in (2.6) to perform the following: The walk keeps jumping toward $`y_t`$, spending at most time $`r_x`$ at each site $`x`$ such that it reaches $`y_t`$ before time $`\gamma _t`$. Then it stays at $`y_t`$ until time $`\gamma _t`$ and then within $`y_t+Q^{(t)}`$ for the remaining time $`t\gamma _t`$.
Inserting this event into (2.6) and invoking Markov property at time $`\gamma _t`$ we get
$$u(t,0)\text{II}\times \text{III},$$
(3.34)
where the same argument as in \[BK00\] shows that $`\text{III}e^{t\alpha (b_t)^2[\lambda _R(\psi )\epsilon ]}`$ for large $`t`$, while for II we have
$`\text{II}={\displaystyle _{\mathrm{}_{y_t}(\gamma _t)}}`$ $`dt_0\mathrm{}dt_{y_t1}e^{2\kappa \gamma _t}\mathrm{exp}\left\{{\displaystyle \underset{k=0}{\overset{n}{}}}\xi _kt_k\right\}{\displaystyle \underset{x=0}{\overset{y_t1}{}}}\mathrm{𝟣}\{t_xr_{x1}\}`$ (3.35)
$`e^{2\kappa \gamma _t}{\displaystyle \underset{x=0}{\overset{y_t1}{}}}\left[r_xe^{r_x\xi (x)}\right]e^{(2\kappa +1)\gamma _t}\mathrm{exp}\left\{{\displaystyle \underset{x=0}{\overset{y_t1}{}}}\mathrm{log}\left(\xi (x)1\right)\right\},`$
where we recalled the notation of (3.8). Now $`y_t\gamma _t`$, so using Lemma 3.4 we have that
$$\text{II}e^{(2\kappa +1)\gamma _t}\mathrm{exp}\left\{\stackrel{~}{G}_\eta ^1(\varrho /\gamma _t)(1+o(1))\right\}=e^{(2\kappa +1)\gamma _tt\alpha (b_t)^3(1+o(1))},$$
(3.36)
where we used the definition of $`\gamma _t`$. Since $`1/\stackrel{~}{G}_\eta `$ is asymptotically concave, $`\gamma _t=\varrho /\stackrel{~}{G}_\eta (t\alpha _{b_t}^3)O(t\alpha _{b_t}^3)`$ and the exponent is $`o(t\alpha _{b_t}^2)`$. Consequently,
$$u(t,0)e^{t\alpha (b_t)^2[\lambda _R(\psi )\epsilon +o(1)]},$$
(3.37)
where $`o(1)`$ still depends on $`\eta `$. The proof is finished by letting $`t\mathrm{}`$ (which eliminates the dependence on $`\eta `$), optimizing over $`\psi `$ and $`R`$ with $`_R(\psi )<1`$ and letting $`\epsilon 0`$. ∎
## Acknowledgments
M.B. would like to thank Yimin Xiao and Oded Schramm for discussions about the behavior of sums of i.i.d. random variables with infinite mean.
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# Dynamics of Fluctuating Bose-Einstein Condensates
## Abstract
We present a generalized Gross-Pitaevskii equation that describes also the dissipative dynamics of a trapped partially Bose condensed gas. It takes the form of a complex nonlinear Schrödinger equation with noise. We consider an approximation to this Langevin field equation that preserves the correct equilibrium for both the condensed and the noncondensed parts of the gas. We then use this formalism to describe the reversible formation of a one-dimensional Bose condensate, and compare with recent experiments. In addition, we determine the frequencies and the damping of collective modes in this case.
The observation of Bose-Einstein condensation in ultracold trapped atomic vapors has offered the possibility to study the equilibrium and nonequilibrium properties of these degenerate gases experimentally and to compare the results with ab initio calculations. The latter have as an input only the mass and scattering length of the particular atom of interest, and the parameters involved in the experimental setup. To describe the various equilibrium and nonequilibrium properties like for example topological excitations, relaxation rates, mode frequencies, damping rates and density profiles, theories have been developed at different levels of sophistication.
At the most elementary level, the Gross-Pitaevskii equation already captures many of the experimentally observed phenomena . It describes in Hartree approximation the zero-temperature dynamics of the condensate, and has been used to explain and predict many features of these Bose-condensed systems. At the next level, the static and dynamic properties of the noncondensed or thermal part of the gas are to be included. This can in first instance be done by including into the Gross-Pitaevskii equation the Hartree-Fock interaction with the thermal cloud and coupling it to an equation for the dynamics of the thermal part of the gas. The latter is in good approximation given by a Boltzmann equation for the single-particle distribution function .
All these theories describe in essence only the average dynamics of the gas. Near the critical region however, fluctuations in the order parameter are generally much larger than the average value of the order parameter itself. Therefore, in this region, it is necessary to include fluctuations into a description of the trapped gas. Also far below the critical temperature, it can be essential in some cases to include fluctuations into a description of the dynamics of the gas. For example, to understand the phenomenon of phase diffusion , one needs to consider fluctuations that disturb the phase of the condensate. From a fundamental point of view, the desired dissipative generalization of the Gross-Pitaevskii equation should obey the so-called fluctuation-dissipation theorem . This guarantees that both the condensed as well as the noncondensed components of the gas will relax to thermal equilibrium. As a result, the order parameter fluctuates around its mean value, and the central quantity describing the dynamics of the condensate is not this mean value, but the actual probability distribution of the order parameter. With all this in mind, a unified theory describing the coherent dynamics of the condensate wave function, the incoherent scattering between the various components of the gas, as well as the fluctuations around the average value of the order parameter has been developed by one of us .
The purpose of this Letter is two-fold. First, we show how the coupled dynamics of the thermal cloud and the condensate can be solved in a selfconsistent way, that on the one hand leads to the correct equilibrium distribution of the trapped gas, and at the same time takes into account both mean-field effects as well as fluctuations. Second, we show that fluctuations can be of crucial importance when trying to understand recent experimental results.
In general, a theoretical description of a trapped interacting Bose gas is possible in terms of the Langevin equation
$`i\mathrm{}{\displaystyle \frac{\mathrm{\Phi }(𝐱,t)}{t}}`$ $`=`$ $`[{\displaystyle \frac{\mathrm{}^2^2}{2m}}+V_{\mathrm{ext}}(𝐱)\mu iR(𝐱,t)`$ (2)
$`+T^{2B}|\mathrm{\Phi }(𝐱,t)|^2]\mathrm{\Phi }(𝐱,t)+\eta (𝐱,t).`$
Here $`\mathrm{}`$ is Planck’s constant, $`m`$ is the mass of the atom, $`V_{\mathrm{ext}}(𝐱)`$ is the external trapping potential, and $`T^{2B}=4\pi \mathrm{}^2a/m`$ is the s-wave approximation to the two-body scattering matrix, with $`a`$ the scattering length. This Langevin equation can be derived using a field-theoretic formulation of the Keldysh formalism , and describes the fluctuations as well as the mean-field effects of both the condensed and the noncondensed parts of the gas. The derivation requires that the high-energy part of the system is sufficiently close to equilibrium that it can be described as having a temperature $`T`$ and a chemical potential $`\mu `$, and can play the role of a ‘heat bath’. This requirement is usually well satisfied and enters through the explicit expression for the imaginary part in our generalized Gross-Pitaevskii equation, which, if the gas is sufficiently close to or below the critical temperature, can be approximated by
$`iR(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{\beta }{4}}\mathrm{}\mathrm{\Sigma }^K(𝐱)[{\displaystyle \frac{\mathrm{}^2^2}{2m}}+V_{\mathrm{ext}}(𝐱)`$ (4)
$`\mu +T^{2B}|\mathrm{\Phi }(𝐱,t)|^2].`$
Here $`\beta `$ is $`1/k_BT`$, with $`k_B`$ Boltzmann’s constant.
If $`R(𝐱,t)`$ has this particular form, the trapped gas will relax to equilibrium, because it enforces the fluctuation-dissipation theorem. Note that Eq. (4) is always valid for the energy levels below the chemical potential, i.e., for the condensate, but causes the energy distribution function for the noncondensed cloud to relax to its ‘classical’ value $`N(ϵ)=[\beta (ϵ\mu )]^1`$. Therefore, it cannot describe the exponential decay of the density of noncondensed atoms at the edges of the thermal cloud. Instead, the density of thermal atoms decays algebraically. In addition, the classical approximation overestimates the average number of atoms for the eigenstates above the chemical potential. Both defects can be cured , but are unimportant for the condensate and the low-energy part of the thermal cloud, where most of the atoms reside, and for which our theory is intended to be valid.
Finally, the correlations of the Gaussian noise $`\eta (𝐱,t)`$ in Eq. (2) are given by
$`\eta ^{}(𝐱,t)\eta (𝐱^{},t^{})`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}^2}{2}}\mathrm{\Sigma }^K(𝐱)\delta (tt^{})\delta (𝐱𝐱^{}),`$ (5)
where the average denotes an average over the different realizations of the noise. The general expressions for $`R(𝐱,t)`$ and $`\mathrm{\Sigma }^K(𝐱)`$ can be found in Ref., and lead to a Keldysh selfenergy $`\mathrm{\Sigma }^K(𝐱)`$ equal to
$`\mathrm{}\mathrm{\Sigma }^K(𝐱)`$ $`=`$ $`{\displaystyle \frac{4i[T^{2B}]^2}{(2\pi )^5\mathrm{}^6}}{\displaystyle 𝑑𝐩_1𝑑𝐩_2𝑑𝐩_3\delta (𝐩_1𝐩_2𝐩_3)}`$ (7)
$`\delta (ϵ_1ϵ_2ϵ_3)(1+N_1)N_2N_3,`$
with $`N_iN(ϵ_i)`$ the Bose distribution for the eliminated part of the gas, and $`ϵ_i=𝐩_i^2/2m+V_{\mathrm{ext}}(𝐱)\mu `$. It determines the strength of the fluctuations through Eq. (5) and is related to the damping by means of the fluctuation-dissipation theorem.
At this point, we explain briefly the experimental setup we are considering in the rest of this Letter. It is inspired by a recent experiment by Stamper-Kurn et. al , and the ideal realization of the conditions mentioned above. In the experiment we are considering, a Bose gas is trapped and cooled to a temperature above the transition temperature. Subsequently, a dimple is created in the external trapping potential, say along the z-axis, by means of optical techniques. This dimple is steep enough such that there is only one energy level in the potential perpendicular to the z-axis. We assume the dimple to be well approximated by a harmonic potential with trapping frequency $`\omega _{}`$. Factorizing the wave function $`\mathrm{\Phi }_0(𝐱)=\mathrm{\Phi }(x,y)\mathrm{\Phi }(z)`$, the atoms trapped inside this dimple form effectively a one-dimensional gas, with an interaction strength $`g=T^{2B}/2\pi l_{}^2`$. Here, $`l_{}=\sqrt{\mathrm{}/m\omega _{}}`$ is the harmonic oscillator length, and $`\mathrm{\Phi }_0`$ the harmonic groundstate in the dimple. By changing the depth of the dimple, its lowest energy level can become lower than the chemical potential of the noncondensed three-dimensional gas. If this situation occurs, the atoms will condense into this ground state. Notice that during this process, the noncondensed gas in the three dimensional trapping potential will remain close to equilibrium, and represents the ‘heat bath’.
We now turn to a numerical solution of Eq. (2) under these conditions, using a combination of well-known techniques . To ensure particle number conservation in the absence of the term $`iR(𝐱,t)`$, we use an implicit method that represents the time evolution operator $`\mathrm{exp}(iH\delta t/\mathrm{})`$ as $`(1+iH\delta t/2\mathrm{})^1\times (1iH\delta t/2\mathrm{})`$. Here, $`H=\mathrm{}^2^2/2m+V_{\mathrm{ext}}(𝐱)\mu iR(𝐱,t)+g|\overline{\mathrm{\Phi }(𝐱,t)}|^2`$, with $`\overline{\mathrm{\Phi }(𝐱,t_i)}`$ the selfconsistent average $`[\mathrm{\Phi }(𝐱,t_i)+\mathrm{\Phi }(𝐱,t_i+\delta t)]/2`$. The numerical method for solving Eq. (2) can now be found from its solution $`\mathrm{\Phi }(𝐱,t_i+\delta t)=\mathrm{exp}(iH\delta t/\mathrm{})[\mathrm{\Phi }(𝐱,t_i)(i/\mathrm{})\mathrm{exp}(iHt_i/\mathrm{})_{t_i}^{t_i+\delta t}𝑑t^{}\mathrm{exp}(iHt^{}/\mathrm{})\eta (𝐱,t^{})]`$ by introducing a new noisy variable $`\xi _i(𝐱)=\mathrm{exp}(iHt_i/\mathrm{})_{t_i}^{t_i+\delta t}dt^{}\mathrm{exp}(iHt^{}/\mathrm{})\eta (𝐱,t^{})]`$. The correlations of $`\xi _i(𝐱)`$ are $`\xi _i^{}(𝐱)\xi _j(𝐱^{})=\frac{i\mathrm{}^2}{2}\mathrm{\Sigma }^K(𝐱)\delta _{ij}\delta t\delta (𝐱𝐱^{})+𝒪(\delta t^2)`$. The spatial discretization is straightforward.
As a first application, we show in Fig. 1 that in the case of a noninteracting vapor, the density of the harmonically trapped one-dimensional gas relaxes to the correct equilibrium given by $`n(z)=_\alpha |\varphi _\alpha (z)|^2/\beta (ϵ_\alpha \mu )`$, where $`ϵ_\alpha `$ are the energy eigenvalues and $`\varphi _\alpha (z)`$ the corresponding eigenfunctions of the Hamiltonian. The equilibrium is shown for several sizes of the spatial mesh, and we see that the density distribution converges towards the continuum limit given by $`[2l_z^2\beta \mathrm{}\omega _z]^{1/2}[\beta (V_{\mathrm{ext}}\mu )]^{1/2}`$ where $`l_z=\sqrt{\mathrm{}/m\omega _z}`$ is the harmonic oscillator length along the $`z`$-direction, only for rather small mesh sizes. This is caused by the relatively large contribution of high-energy states due to the classical behavior of the thermal cloud. We emphasize, however, that given a certain mesh size, the correct equilibrium corresponding to that particular discretization is reproduced numerically. Moreover, it is important to realize that if we do not include the noise, the density of the gas would be zero. Hence, it is clear that the fluctuation-dissipation theorem ensures that the noise and the imaginary term in Eq. (2) cooperate in order to occupy the energy levels thermally. Thus, including fluctuations is crucial in treating the effects of the thermal cloud.
Second, we calculate the damping of the breathing mode in a trapped one-dimensional Bose gas below the critical temperature. Shown in Fig. 2 is $`g`$ times the value of the density $`n(\mathrm{𝟎})`$ in the center of the one-dimensional trap. First, the gas relaxes towards equilibrium at an effective chemical potential $`\mu ϵ_0=30\mathrm{}\omega _z`$, where $`ϵ_0`$ is the eigenvalue of $`\mathrm{\Phi }_0(x,y)`$ in the dimple. The fact that the value of $`gn(\mathrm{𝟎})`$ does not relax to $`30\mathrm{}\omega _z`$ can be explained by realizing that mean-field effects are included in our formalism, because $`\mathrm{\Phi }(z,t)`$ describes both the condensed and the noncondensed parts of the one-dimensional gas. This implies that the chemical potential for the condensed part of the gas is effectively lowered by the interaction with the noncondensed part of the gas. After allowing the gas to equilibrate, at $`t=0.61`$ s the trapping frequency is instantaneously changed to $`\omega _z^{}=0.95\omega _z`$, and the gas starts to oscillate. The frequency $`\omega `$ and the damping rate $`\gamma `$ of the oscillation are found by fitting to $`n(\mathrm{𝟎},t)=n_{\mathrm{eq}}(\mathrm{𝟎})+\delta n(\mathrm{𝟎})\mathrm{exp}(\gamma t)[1\mathrm{cos}(\omega t)]`$, which describes both the relaxation of the density to its new equilibrium, as well as the excitation and damping of the oscillation. The results are $`\omega =136.84`$ s<sup>-1</sup>, which is close to the result of $`\sqrt{5/2}\omega _z=129.1`$ s<sup>-1</sup> expected from the Gross-Pitaevskii equation , $`\gamma =12.3`$ s<sup>-1</sup>, $`g\delta n(\mathrm{𝟎})=0.84\mathrm{}\omega _z`$, and $`gn_{\mathrm{eq}}(\mathrm{𝟎})=28\mathrm{}\omega _z`$. Note that the damping is in principle caused both by the collisions with the reservoir of thermal atoms in the three-dimensional trapping potential, which is described by $`\mathrm{\Sigma }^K(𝐱)`$, as well as by the nonlinear term in the Langevin equation which induces damping upon averaging over the different realizations of the noise. The latter includes collisional, as well as Landau damping.
In Fig. 3, a snapshot of the density profile is shown at several times during the initial growth of the condensate. These are taken from the simulation presented in Fig. 2. As expected, the condensate shows up as a peak in the density profile. Also shown are the Thomas-Fermi solution $`[\mu V_{\mathrm{ext}}(z)]/g`$ for the condensate density and the noninteracting result for the noncondensed part of the gas which diverges at the critical point where the effective chemical potential becomes equal to zero. This is due to the fact that in one dimension, mean-field theory completely fails near the critical temperature. Fig. 3 shows that far enough from the condensate, the noninteracting result is obtained. Again, in the center of the trap one obtains the Thomas-Fermi result only approximately due to interactions.
Finally, we consider the reversible formation of a one-dimensional Bose condensate, similar in spirit to the experiments by Stamper-Kurn et. al . The dimple perturbing the external trapping potential is now oscillating in such a way, that the lowest energy level in the dimple crosses the chemical potential of the three dimensionally trapped gas several times according to $`ϵ_0=\mu +\mu \mathrm{sin}(\omega t)`$, with $`\omega =2\pi `$ s<sup>-1</sup> and $`\mu =30\mathrm{}\omega _z`$. The calculations have been done with and without noise. The results show that without noise, the condensate evaporation and growth cannot be described properly. Indeed, in that case, our findings depend strongly on the initial conditions. Moreover, the periodic growth occurs with decreasing amplitude.
To describe the growth cycles correctly, one needs to include fluctuations into the generalized Gross-Pitaevskii equation, which is to be expected since we are at times in the critical region. Notice that we quantitatively reproduce the lagging behind of the condensate as observed in an experiment performed with an essentially three-dimensional dimple. We observed a lagging behind of roughly $`0.1`$s, whereas in experiment $`0.07`$s was measured .
In conclusion, we have shown that our generalized stochastic Gross-Pitaevskii equation consistently takes into account the dissipative dynamics of a trapped partially Bose-condensed gas, and describes the equilibrium density profile, condensate growth, coherent dynamics, and damping in a unified way. In our opinion, the one-dimensional experiment considered here would be ideal for a detailed comparison between theory and experiment in the problem of condensate growth, because the three-dimensional cloud remains in equilibrium. In addition, a numerical study of the two-dimensional case would be of interest and might be used to investigate for example the dissipative dynamics and the formation of vortices in a rotating Bose gas.
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# Two pion electroproduction Presented at MESON’2000 Workshop, Krakow, May 19 - May 23, 2000
## 1 Introduction
The $`\gamma N\pi \pi N`$ reaction in nuclei has captured some attention recently and has proved to be a source of information on several aspects of resonance formation and decay as well as a test for chiral perturbation theory at low energies. A model for the $`\gamma p\pi ^+\pi ^{}p`$ reaction was developed in containing 67 Feynman diagrams by means of which a good reproduction of the cross section was found up to about $`E_\gamma 1`$ GeV.
A more reduced set of diagrams, with 20 terms , was found sufficient to describe the reaction up to $`E_\gamma 800`$ MeV where the Mainz experiments are done .
The extension of this kind of work to virtual photons should complement the knowledge obtained through the ($`\gamma `$,$`2\pi `$) and the related reactions. The coupling of the photons to the resonances depends on $`q^2`$ and the dependence can be different for different resonances. Hence, the interference of different mechanisms pointed above will depend on $`q^2`$ and with a sufficiently large range of $`q^2`$, one can pin down the mechanism of ($`\gamma `$,$`2\pi `$) with real or virtual photons with more precision than just with real photons, which would help settle the differences between present theoretical models.
However, there are already interesting two pion electroproduction experiments selecting $`\mathrm{\Delta }`$ in the final state. The reactions are, $`epe^{}\pi ^{}\mathrm{\Delta }^{++}`$ and $`epe^{}\pi ^+\mathrm{\Delta }^0`$ . It is thus quite interesting to extend present models of $`(\gamma ,2\pi )`$ to the realm of virtual photons and compare with existing data. In our paper we do so, extending the model of ref. to deal with the electroproduction process. This model is flexible enough and one can select the diagrams which contain $`\mathrm{\Delta }\pi `$ in the final state in order to compare directly with the measured cross sections.
The extension of the model requires three new ingredients: the introduction of the zeroth component of the photon coupling to resonances (calculations where done in in the Coulomb gauge, $`ϵ^0`$, where the zeroth component is not needed), the implementation of the $`q^2`$ dependence of the amplitudes, which will be discussed in forthcoming sessions, and the addition of the explicit terms linked to the $`S_{1/2}`$ helicity amplitudes which vanish for real photons.
Experiments on ($`\gamma _v`$,$`2\pi `$) are presently being done in the Thomas Jefferson Laboratory , both for $`N\mathrm{\Delta }`$ and $`N\pi \pi `$ production.
## 2 Model for $`eNe^{}\mathrm{\Delta }\pi `$
We will evaluate cross sections of virtual photons integrated over all the variables of the pions and the outgoing nucleon. In this case the formalism is identical to the one of inclusive $`eNe^{}X`$ scattering or pion electroproduction after integrating over the pion variables . For the model of the $`\gamma _vN\mathrm{\Delta }\pi `$ reaction we take the same diagrammatic approach as in ref. and select the diagrams which have a $`\mathrm{\Delta }`$ in the final state. The diagrams which contribute to the process are depicted in fig.1
We follow the paper from Devenish et al. in our approach to electromagnetic transitions for Roper and $`N^{}(1520)`$ resonances. As we are working with virtual photons we need to care about these couplings and hence include terms which vanish for real photons.
Gauge invariance is one of the important elements in a model involving photons and implies that
$$T^\mu q_\mu =0$$
(1)
However, as discussed in the study of the $`eNe^{}N\pi `$ reaction in , and as can be easily seen by inspection of the diagrams and the amplitudes, the constraint of eq. (1) still requires the equality of four electromagnetic form factors,
$$F_1^p(q^2)=F_1^\mathrm{\Delta }(q^2)=F_{\gamma \pi \pi }=F_c(q^2)$$
(2)
The form factors of eq. (2) are respectively the $`\gamma NN`$, $`\gamma \mathrm{\Delta }\mathrm{\Delta }`$, $`\gamma \pi \pi `$ and $`\gamma \mathrm{\Delta }N\pi `$ ones. These form factors are usually parametrized in different forms, except for $`F_1^p(q^2)`$ and $`F_1^\mathrm{\Delta }(q^2)`$ which are taken equal, as it would come from ordinary quark models.
Although the model is gauge invariant with the prescription of eq. (2) there is the inconvenience that the results depend upon which one of the three form factors we take for all of them.
We should note however, that the dominant term, by large, is the $`\mathrm{\Delta }`$ Kroll Ruderman and pion pole terms. This is also so in the test of gauge invariance where the two terms involving the $`F_1^p(q^2)`$ form factor in diagrams D4, D6 give only recoil contributions of the order of O($`p_\pi `$/m) in eq. (1). This justifies the use of $`F_c(q^2)`$ or $`F_{\gamma \pi \pi }(q^2)`$ for all the form factors.
There is, however, another way to respect gauge invariance, while at the same time using different form factors which is proposed in and to which we refer in what follows as Berends et al. approach.
## 3 Results and conclusions
We have tested our results with the experimental data of refs. . We show the cross section of $`\gamma _vp\mathrm{\Delta }^{++}\pi ^{}`$ and $`\gamma _vp\mathrm{\Delta }^0\pi ^+`$ ($`\mathrm{\Delta }^0\pi ^{}p`$), as a function of W, the virtual photon-proton ($`\gamma _vp`$) center of mass energy, and for different values of $`Q^2`$. We have made different calculations. One of them corresponds to using all form factors equal (which we set to $`F_{\gamma \pi \pi }`$) with two different values of $`\lambda _\pi ^2`$, 0.5 $`GeV^2`$ and 0.6 $`GeV^2`$. In we see that the cross section increases by about 10 $`\%`$ when going from $`\lambda _\pi ^2`$=0.5 $`GeV^2`$ and $`\lambda _\pi ^2`$=0.6 $`GeV^2`$. We also show the results taking $`F_1^p`$, $`F_1^\mathrm{\Delta }`$ and setting $`F_c=F_{\gamma \pi \pi }`$ with $`\lambda _\pi ^2`$=0.6 $`GeV^2`$. This latter calculation is not gauge invariant. However we see that the deviation with respect to the gauge invariant one assuming all form factors equal is very small . This reflects the fact that the relevant terms in the model are those involving $`F_{\gamma \pi \pi }`$ and $`F_c`$, the pion pole and $`\mathrm{\Delta }`$ Kroll Ruderman terms.
We also evaluate the cross section using Berends gauge invariant approach with different form factors . We show the results in fig. 2. The continuous line in the figure is obtained with this prescription using $`F_1^p`$, $`F_1^\mathrm{\Delta }`$ but setting $`F_c=F_{\gamma \pi \pi }`$ with $`\lambda _\pi ^2`$ = 0.5 $`GeV^2`$.
We see that these results are remarkably similar to those where $`F_c`$ and $`F_{\gamma \pi \pi }`$ had the same values as here but $`F_1^p`$, $`F_1^\mathrm{\Delta }`$ were set equal to $`F_{\gamma \pi \pi }`$ in order to preserve gauge invariance.
The dotted line in fig. 2 corresponds to the same parametrization for $`F_c`$ as for $`F_{\gamma \pi \pi }`$ but parameter $`\lambda _c^2`$= 0.8 $`GeV^2`$. This shows the sensitivity of the results to $`F_c`$ which appears in the dominant Kroll-Ruderman term.
In summary we could remark the following points: We have shown in that the peak in the cross section is due to an interference between the $`\mathrm{\Delta }`$ Kroll Ruderman term and the $`N^{}(1520)`$ excitation process followed by $`\mathrm{\Delta }\pi `$ decay. Different sets of form factors have been used in our model in order to show the sensitivity of the results to these changes. These tests should be useful in view of the coming data and the possibility to extract relevant information from them. We have calculated the separation of the transverse and longitudinal cross sections and found that the transverse one largely dominates the cross sections. Finally, it is also interesting to note that the present model is just part of a more general $`\gamma _vN\pi \pi N`$ model which selects only the terms where a $`\pi N`$ pair of the final state appears forming a $`\mathrm{\Delta }`$ state.
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# Untitled Document
THE MAXIMAL KINEMATICAL INVARIANCE GROUP OF FLUID
DYNAMICS AND EXPLOSION-IMPLOSION DUALITY
L. O’Raifeartaigh<sup>(a)</sup> Deceased and V. V. Sreedhar<sup>(b)</sup> sreedhar@stp.dias.ie
School of Theoretical Physics
Dublin Institute for Advanced Studies
10, Burlington Road
Dublin 4, Ireland
Abstract
It has recently been found that supernova explosions can be simulated in the laboratory by implosions induced in a plasma by intense lasers. A theoretical explanation is that the inversion transformation, ($`\mathrm{\Sigma }:t1/t,𝐱𝐱/t`$), leaves the Euler equations of fluid dynamics, with standard polytropic exponent, invariant. This implies that the kinematical invariance group of the Euler equations is larger than the Galilei group. In this paper we determine, in a systematic manner, the maximal invariance group $`𝒢`$ of general fluid dynamics and show that it is a semi-direct product $`𝒢=SL(2,R)G`$, where the $`SL(2,R)`$ group contains the time-translations, dilations and the inversion $`\mathrm{\Sigma }`$, and $`G`$ is the static (nine-parameter) Galilei group. A subtle aspect of the inclusion of viscosity fields is discussed and it is shown that the Navier-Stokes assumption of constant viscosity breaks the $`SL(2,R)`$ group to a two-parameter group of time translations and dilations in a tensorial way. The 12-parameter group $`𝒢`$ is also known to be the maximal invariance group of the free Schrödinger equation. It originates in the free Hamilton-Jacobi equation which is central to both fluid dynamics and the Schrödinger equation. DIAS-STP-00-16
Introduction
Considerable efforts are being devoted at present in centres such as the National Ignition Facility at the Lawrence Livermore Laboratory, U.S.A., and the Laser MegaJoule in Bordeaux, France, to simulate astrophysical systems in the laboratory. The motivation for this programme comes from observational evidence that the dynamics of the mixing of gases during a supernova explosion is very similar to the turbulent splashing and mixing of a plasma as a fusion capsule is bombarded by high intensity laser beams . Since the former physical situation deals with an explosion while the latter concerns an implosion, and because the time and length scales involved in the two cases are drastically different, it was a challenging task to produce a theoretical explanation for the observed similarity.
In a recent paper , Drury and Mendonça produced such an explanation by showing that the Euler equations of fluid dynamics which govern both the systems are left invariant by an inversion transformation $`\mathrm{\Sigma }:t1/t,𝐱𝐱/t`$. The discovery of this transformation was inspired by similar transformations of the form $`𝐱f(t)𝐱`$ that are used by cosmologists to factor out an arbitrary uniform expansion or contraction of the system . As is well-known, the Euler equations of fluid dynamics are invariant under the Galilei group . The results of show that the maximal kinematical<sup>1</sup> Meaning transformations involving only space and time coordinates as opposed to those involving internal or gauge degrees of freedom. invariance group of fluid dynamics $`𝒢`$ is larger than the Galilei group . In this paper we determine, in a systematic manner, the structure of $`𝒢`$ for general fluid dynamics. We also discuss subtleties associated with the inclusion of viscosity fields and show that the Navier-Stokes assumption of constant viscosity leads to a reduction of $`𝒢`$ in a tensorial manner.
It will be shown that $`𝒢`$ is of the form
$$𝒢=SL(2,R)G$$
$`(1)`$
where $`G`$ is the connected, static Galilei group which induces the transformations
$$𝐱R𝐱+𝐚+𝐯t,tt$$
$`(2)`$
and $`SL(2,R)`$ is the group
$$t\frac{\alpha t+\beta }{\gamma t+\delta },𝐱\frac{𝐱}{\gamma t+\delta };\alpha \delta \beta \gamma =1$$
$`(3)`$
which includes parity, time translations, dilations and the inversion $`\mathrm{\Sigma }`$.
$`𝒢`$ is isomorphic to the Niederer group, which is the maximal invariance group of the free Schrödinger equation . The reason for this isomorphism is that the kinetic part of the fluid-dynamic Lagrangian (see (17)) is of the free Hamilton-Jacobi form, which is the classical limit of the Schrödinger operator.
The Strategy of Proof
The general fluid dynamic equations in $`n`$-dimensional space are
$$D\rho =\rho 𝐮$$
$`(4a)`$
$$\rho D𝐮=p+𝐕$$
$`(4b)`$
$$Dϵ=(ϵ+p)𝐮$$
$`(4c)`$
where the convective derivative $`D`$ and the viscosity terms $`𝐕`$ are defined by
$$D=\frac{}{t}+𝐮\text{and}V_i=_j\left(\eta (_ju_i+_iu_j\frac{2}{n}\delta _{ij}_ku_k)\right)+_i(\zeta _ku_k)$$
$`(5)`$
respectively. In the above equations $`\rho `$, u, $`p`$ and $`ϵ`$ stand for the density, the velocity vector field, the pressure, and the energy density of the fluid respectively and $`\eta ,\zeta `$ are the shear and bulk viscosity fields. The above differential equations are augmented by an algebraic condition called the polytropic equation of state which relates the pressure to the energy density thus:
$$p=(\gamma _o1)ϵp+ϵ=\gamma _oϵ$$
$`(6)`$
$`\gamma _o`$ being a constant called the polytropic exponent. This equation may be used to eliminate $`p`$ from (4). Further, by making the substitution,
$$ϵ=\chi \rho ^{\gamma _o}$$
$`(7)`$
the original set of equations (4) reduces to
$$D\rho =\rho 𝐮$$
$`(8a)`$
$$\rho D𝐮=(\gamma _o1)(\chi \rho ^{\gamma _o})+𝐕$$
$`(8b)`$
$$D\chi =0$$
$`(8c)`$
To find the maximal invariance group of these equations we first note that, for any system (such as the above) for which there is no feedback from the fields to the space-time transformations, the maximal invariance group $`𝒢`$ of the general class of configurations is also an invariance group for any special sub-class of configurations; though not necessarily a maximal one, since a restricted set of configurations could have a larger symmetry group. Formally,
$$𝒞_s𝒞𝒢𝒢_s$$
$`(9)`$
where $`𝒞`$ denotes general configurations and $`𝒞_s`$ denotes a special class. Thus a general strategy for finding maximal invariance groups is to first find a (tractable) sub-class of configurations $`𝒞_s`$ for which the maximal invariance group $`𝒢_s`$ can be found, and then see what conditions are imposed on $`𝒢_s`$ by the general configurations. This is the strategy we shall adopt; except that in our case we choose a sub-class for which the generalisation to arbitrary configurations imposes no further conditions i.e. $`𝒢_s𝒢`$. The tractable sub-class we choose is obtained by making three simplifications, namely, ignoring the effects of viscosity by tuning the viscosity fields $`\eta `$ and $`\zeta `$ to zero, by setting $`\chi =1`$, and by letting the velocity vector $`𝐮`$ be curl-free. The resulting sub-class is the configuration of inviscid, isentropic and irrotational flows which are well-known in fluid dynamics.
The Sub-class of Inviscid, Isentropic, and Irrotational Fluids
The first simplification we consider is to ignore viscosity effects in which case (8) reduces to
$$D\rho =\rho 𝐮$$
$`(10a)`$
$$\rho D𝐮=(\gamma _o1)(\chi \rho ^{\gamma _o})$$
$`(10b)`$
$$D\chi =0$$
$`(10c)`$
Second, we note that it is consistent with the field equations to set $`\chi =1`$, with $`𝐮`$ not necessarily curl-free, in which case the equations in (10) reduce to,
$$D\rho =\rho 𝐮\text{and}D𝐮=\gamma _o(\gamma _o1)\rho ^{\gamma _o2}\rho $$
$`(11)`$
The reason for making this simplification is that the system becomes a local Lagrangian one. This can be seen by using the standard Clebsch parametrization
$$𝐮=\varphi \nu \theta $$
$`(12)`$
for the vector-field $`𝐮`$, in which case the local Lagrangian density is,
$$=\rho \left[\dot{\varphi }+\nu \dot{\theta }\frac{1}{2}(\varphi +\nu \theta )^2\rho ^{\gamma _o1}\right]$$
$`(13)`$
It will be convenient to write this Lagrangian density in the form
$$=\rho \left[\dot{\varphi }\frac{1}{2}(\varphi )^2+\nu 𝒟\theta \frac{\nu ^2}{2}(\theta )^2\rho ^{\gamma _o1}\right]\text{where}𝒟=\frac{}{t}\varphi $$
$`(14)`$
is the convective derivative $`D`$ restricted to the curl-free part of $`𝐮`$. Clearly $`\nu `$ is a Lagrange multiplier, and, by varying it, we obtain
$$\nu =\frac{𝒟\theta }{(\theta )^2}\text{and}=\rho \left[\dot{\varphi }\frac{1}{2}(\varphi )^2\frac{1}{2}\frac{(𝒟\theta )^2}{(\theta )^2}\rho ^{\gamma _o1}\right]$$
$`(15)`$
Note that the first equation here is equivalent to $`D\theta =0`$ since
$$D\theta =𝒟\theta (\nu \theta )\theta =𝒟\theta \nu (\theta )^2=(\theta )^2\left(\frac{𝒟\theta }{(\theta )^2}\nu \right)$$
$`(16)`$
We now proceed to the third and last simplification mentioned.
The third simplification we make is to consider the curl-free case $`𝐮=\varphi `$ or $`\nu =\theta =0`$. Then the Action corresponding to the Lagrangian density (13) reduces to
$$S=d^nx𝑑t\rho \left[\dot{\varphi }\frac{1}{2}(\varphi )^2\right]\rho ^{\gamma _o}$$
$`(17)`$
Note that the quantity in square brackets is just the Hamilton-Jacobi function for a free particle. Following our strategy we now seek the maximal invariance group of the Action (17).
The most general transformation involving the fields, as shown in the Appendix, takes the linear, inhomogeneous form
$$\xi _i=\xi _i(x,t),\tau =\tau (x,t),\stackrel{~}{\varphi }=s(x,t)\varphi +\lambda (x,t),\stackrel{~}{\rho }=\mu (x,t)\rho +\mu ^{}(x,t)$$
$`(18)`$
where $`s`$ is a constant. However, one sees at once that if $`\mu ^{}0`$ the form invariance of the Action is violated by inhomogeneous terms. Furthermore, since
$$\frac{}{t}=\frac{\tau }{t}\frac{}{\tau }+\frac{\xi _i}{t}\frac{}{\xi _i}\text{and }\frac{}{x_i}=\frac{\tau }{x_i}\frac{}{\tau }+\frac{\xi _j}{x_i}\frac{}{\xi _j}$$
$`(19)`$
we see that we obtain a term $`(\stackrel{~}{\varphi }/\tau )^2`$ unless
$$\frac{\tau }{x_i}=0$$
$`(20)`$
This means that the $`\varphi /t`$ term does not pick up an $`x`$-dependent factor relative to $`\varphi /\tau `$. Finally we note that under the transformations (18), the second term in the Hamilton-Jacobi part of the Action (17) gives
$$\frac{1}{2}h_{ij}\frac{\varphi }{x_i}\frac{\varphi }{x_j}\frac{1}{2}g_{ij}\frac{\varphi }{\xi _i}\frac{\varphi }{\xi _j}\text{where}g_{ij}=\frac{\xi _i}{x_k}\frac{\xi _j}{x_l}h_{kl}\mathrm{\Omega }(x)g_{ij}^{}(\xi )$$
$`(21)`$
where $`h_{ij}`$ is the flat Euclidean metric. As is well-known, the most general transformations that have the above property are the special conformal transformations with parameter $`b`$, for which $`\mathrm{\Omega }=(1+2bx+b^2x^2)^2`$, and the Euclidean transformations for which $`\mathrm{\Omega }=1`$. It is easy to see however that, for the Action in (17) to remain invariant, the $`\varphi /t`$ part of the Hamilton-Jacobi function must transform in the same way as the $`(\varphi /x)^2`$ part. But since (20) forbids $`\varphi /t`$ to pick up an $`x`$-dependent factor, the conformal transformations in (21) are not allowed. Thus the possible transformations reduce to
$$\xi _i=M_{ij}(t)x_j+w_i(t),\tau =\tau (t),\stackrel{~}{\varphi }=s\varphi +\lambda (x,t)\stackrel{~}{\rho }=\mu (x,t)\rho $$
$`(22)`$
where $`s`$ is a constant and $`M`$ is Euclidean i.e. $`M=f(t)R(t)`$, where $`f`$ is a scale factor and $`R(t)`$ is a rotation matrix. In that case
$$\left(\begin{array}{cc}\frac{\tau }{t}& \frac{\tau }{x_j}\\ \frac{\xi _i}{t}& \frac{\xi _i}{x_j}\end{array}\right)=\left(\begin{array}{cc}\dot{\tau }& 0\\ \dot{\xi }_i& M_{ij}\end{array}\right)\left(\begin{array}{c}\dot{\tau }0\\ \text{det}M0\end{array}\right)$$
$`(23)`$
and the Action becomes
$$\begin{array}{cc}\hfill S& =\frac{d^n\xi d\tau }{\text{det}M^n\dot{\tau }}\mu ^1\stackrel{~}{\rho }\left[(\dot{\tau }\frac{}{\tau }+\dot{\xi }_i\frac{}{\xi _i})\frac{(\stackrel{~}{\varphi }\lambda )}{s}\frac{\text{det}M^2}{2s^2}(\frac{\stackrel{~}{\varphi }}{\xi _i}\frac{\lambda }{\xi _i})^2(\mu ^1\stackrel{~}{\rho })^{\gamma _o1}\right]\hfill \\ & =\frac{d^n\xi d\tau }{s\mu \text{det}M^n}\stackrel{~}{\rho }\left[(\frac{}{\tau }+\frac{\dot{\xi }_i\frac{}{\xi _i}}{\dot{\tau }})(\stackrel{~}{\varphi }\lambda )\frac{\text{det}M^2}{2s\dot{\tau }}(\frac{\stackrel{~}{\varphi }}{\xi _i}\frac{\lambda }{\xi _i})^2\frac{s(\mu ^1)^{\gamma _o1}}{\dot{\tau }}\stackrel{~}{\rho }^{\gamma _o1}\right]\hfill \end{array}$$
$`(24)`$
Invariance then requires that
$$s\mu (\text{det}M)^n=1,(\text{det}M)^2=s\dot{\tau },\text{and}(\mu ^1)^{\gamma _o1}=\frac{\dot{\tau }}{s}$$
$`(25)`$
and
$$(\frac{}{\tau }+\frac{\xi _i}{\tau }\frac{}{\xi _i})(\stackrel{~}{\varphi }\lambda )\frac{1}{2}\left(\frac{\stackrel{~}{\varphi }}{\xi }\frac{\lambda }{\xi }\right)^2=\frac{\stackrel{~}{\varphi }}{\tau }\frac{1}{2}(\frac{\stackrel{~}{\varphi }}{\xi })^2$$
$`(26)`$
Equations (25) can be readily solved to get
$$s=1(\text{det}M)^2=\dot{\tau },\mu =(\text{det}M)^n,\text{and}\gamma _o=1+\frac{2}{n}$$
$`(27)`$
The result for $`\gamma _o`$ is in agreement with ; in particular in three-dimensional space, $`n=3`$, and $`\gamma _o=\frac{5}{3}`$, as expected. Substituting these results in (26), and requiring that it be satisfied for all $`\stackrel{~}{\varphi }/\xi _i`$, splits it into
$$\frac{\xi _i}{\tau }=\frac{\lambda }{\xi }_i\text{and}\frac{\lambda }{\tau }\frac{1}{2}\left(\frac{\lambda }{\xi }\right)^2=0,\text{where}\dot{\tau }=\text{det}M^2$$
$`(28)`$
Notice that the second equation is just the free Hamilton-Jacobi equation with Action $`\lambda `$. The maximal invariance group can now be obtained by solving the above equations for the functions $`\xi _i`$ and $`\lambda `$. Thus the invariance group $`𝒢`$ is determined by the structure of the Hamilton-Jacobi equation for a free particle. However, we will now show, even before solving these equations, that the group of invariance remains the same when the simplifications made in this section are relaxed and we return to the case of the general fluid configurations.
Proof of Maximality for General Fluid Configurations
We first note that the first and last equations in (28) are exactly the conditions for the restricted derivative $`𝒟`$ to transform covariantly since
$$\begin{array}{cc}\hfill 𝒟(x,t,\varphi )& =\frac{}{t}\frac{\varphi }{x_i}\frac{}{x_i}=\left(\dot{\tau }\frac{}{\tau }+\frac{\xi _i}{t}\frac{}{\xi _i}\right)\left(\frac{\xi _j}{x_i}\frac{\xi _k}{x_i}\right)\frac{\varphi }{\xi _j}\frac{}{\xi _k}\hfill \\ & =\left(\dot{\tau }\frac{}{\tau }\text{det}M^2\frac{\stackrel{~}{\varphi }}{\xi _i}\frac{}{\xi _i}\right)+\left(\frac{\xi _i}{t}+\text{det}M^2\frac{\lambda }{\xi _i}\right)\frac{}{\xi _i}=\text{det}M^2𝒟(\xi ,\tau ,\stackrel{~}{\varphi })\hfill \end{array}$$
$`(29)`$
Let us now relax the assumption that $`𝐮`$ is curl-free, with $`\chi =1`$ still being true. From the Lagrangian density for this case,
$$=\rho \left[\dot{\varphi }\frac{1}{2}(\varphi )^2\frac{1}{2}\frac{(𝒟\theta )^2}{(\theta )^2}\rho ^{\gamma _o1}\right]$$
$`(30)`$
and the covariant transformations,
$$𝒟(x,t,\varphi )=(\text{det}M)^2𝒟(\xi ,\tau ,\stackrel{~}{\varphi })\text{and}(x)=(\text{det}M)(\xi )$$
$`(31)`$
we see that (30) will be invariant under the given space-time transformations provided only that $`\nu `$ and $`\theta `$ are scalars. Thus the case of general $`𝐮`$ puts no restrictions on the space-time transformations but determines the transformation properties of the $`\nu `$ and $`\theta `$ components of $`𝐮`$. In particular we see that the full convective derivative $`D`$ and $`𝐮`$ have the transformation properties
$$D(x,t,\varphi )=(\text{det}M)^2D(\xi ,\tau ,\stackrel{~}{\varphi }),\text{and}\stackrel{~}{𝐮}=\frac{1}{\text{det}M}\left(𝐮\lambda \right)$$
$`(32)`$
Note that the inhomogeneous part of the $`𝐮`$-transformation comes from the curl-free part of $`𝐮`$. Furthermore we see that the $`𝒢`$ transformations do not preserve the condition of incompressibility ($`𝐮=0`$) and hence we need to consider general, compressible fluid configurations.
All this is within the Lagrangian framework. In order to relax the $`\chi =1`$ simplification, we need to allow for general fields $`\chi `$. For this we must go outside the Lagrangian framework and consider the field equations. But that is simple since the field equations with $`\chi =1`$ are invariant and we see by inspection that the equations for general $`\chi `$ remain invariant provided that $`\chi `$ transforms as a scalar. Thus the set of equations in (10) is invariant with respect to the maximal invariance group provided $`\nu ,\theta `$ and $`\chi `$ transform as scalars. Thus the maximal invariance group for $`\nu =\theta =0`$, $`\chi =1`$ (corresponding to our sub-class without viscosity) is the maximal invariance group for all inviscid configurations.
Finally we include viscosity effects. It is straightforward to verify that the viscosity terms in (4b) are covariant only if $`𝐮=0`$ or $`_i\zeta =0`$. However, as already mentioned, the former condition is not invariant under $`𝒢`$ transformations and hence covariance requires that $`\zeta `$ is a constant in space (although not in time). In this case, the fluid equations (4) would remain invariant if the viscosity fields transform as tensors of rank $`n`$, where $`n`$ is the space dimension i.e. $`\eta \text{det}M^n\eta `$ and $`\zeta \text{det}M^n\zeta `$. This in turn implies that the equations would remain invariant under the Navier-Stokes assumption of constant viscosity only if det$`M`$ is a constant which breaks the $`𝒢`$ symmetry in a tensorial way. Therefore the maximal invariance group of the Navier-Stokes equations contains just the static Galilei group plus dilations and time translations corresponding to det$`M`$ being a constant, as will be clear when we solve for the $`𝒢`$ transformation functions in the next section.
It is worth mentioning that the breakdown of $`𝒢`$ by the Navier-Stokes restriction does not contradict (9) because, in contrast to the other restrictions $`\eta =\zeta =0`$, $`\chi =1`$, $`\times 𝐮=0`$ that we have considered, this restriction has a feed-back effect on the non-linear part of the space-time transformations.
The Solutions for the Transformation Functions
Returning to the equations in (28), it is first useful to note from the definition of $`\xi `$ in (22) that
$$\frac{\xi _i}{\tau }=A_{ij}\xi _j+W_i\text{ where}AM_\tau M^T=f_\tau f^1+R_\tau R^T\text{and}W=M\frac{}{\tau }(M^1w)$$
$`(33)`$
where the subscript $`\tau `$ denotes differentiation with respect to $`\tau `$. Hence the first equation in (28) becomes
$$\frac{\lambda }{\xi _i}=A_{ij}\xi _jW_i\left(\frac{^2\lambda }{\xi _i\xi _j}=A_{ij}\right)$$
$`(34)`$
Here the equation in the brackets shows that $`A`$ is a symmetric matrix. But since $`R`$ is a rotation matrix, $`R_\tau R^T`$ is in the Lie algebra of the rotation group and hence is antisymmetric. Thus we have the result $`R_\tau =0`$ and hence $`R`$ is a constant (rigid) rotation matrix.<sup>2</sup> With this in mind, we suppress the matrix for convenience, and restore it later.
It then follows that det$`M=f`$ and from (28) and (32) we have
$$\dot{\tau }=f^2\text{and}\stackrel{~}{𝐮}=\frac{1}{f}\left(𝐮\lambda \right)$$
$`(35)`$
The definitions of $`A`$ and $`W`$ then simplify to
$$A=f_\tau f^1\text{and}W_i=f_\tau \left(\frac{w_i}{f}\right)$$
$`(36)`$
Integrating (34) with respect to $`\xi `$ we have
$$\lambda =A\frac{\xi ^2}{2}W_i\xi _ih(\tau )$$
$`(37)`$
Substituting this into the second equation in (28) we get
$$\frac{(\xi )^2}{2}\frac{A}{\tau }\xi _i\frac{W_i}{\tau }\frac{h}{\tau }=\frac{1}{2}(A\xi +W)^2$$
$`(38)`$
Comparing the coefficients of the powers of $`\xi `$ breaks this into
$$\frac{A}{\tau }=A^2,\frac{W_i}{\tau }=AW_i=\frac{_\tau f}{f}W_i,\frac{h}{\tau }=\frac{W^2}{2}$$
$`(39)`$
We can solve the second equation in (39) explicitly to get
$$W_i=\frac{v_i}{f}_t\left(\frac{w_i}{f}\right)=v_iw_i=f(t)(v_it+a_i)$$
$`(40)`$
where the $`v_i`$ and $`a_i`$ are constants. It follows that
$$\xi _ifx_i+w_i=f(t)(x_i+a_i+v_it)$$
$`(41)`$
The third equation in (39) can then be written as
$$\frac{h}{\tau }=\frac{v^2}{2f^2}\frac{h}{t}=\frac{v^2}{2}h=h_0\frac{v^2}{2}t$$
$`(42)`$
Since the transformations of the fields are given by $`\stackrel{~}{\rho }=f^n\rho `$ and $`\stackrel{~}{\varphi }=\varphi +\lambda `$, we see that everything is determined by the function $`f`$ which satisfies
$$\frac{}{\tau }\left(\frac{\dot{f}}{f^3}\right)=\left(\frac{\dot{f}}{f^3}\right)^2$$
$`(43)`$
Using $`\dot{\tau }=f^2`$, this is easily seen to be equivalent to
$$_\tau ^2f=0\text{and}\frac{\stackrel{\text{}}{𝜏}}{\dot{\tau }}\frac{3}{2}\left(\frac{\ddot{\tau }}{\dot{\tau }}\right)^2=0$$
$`(44)`$
The left hand side of the second equation is the Schwarzian derivative of $`\tau `$ and it follows that the general solution of the above equation is
$$\tau =\frac{\alpha t+\beta }{\gamma t+\delta }\text{and}f=\frac{1}{\gamma t+\delta }\text{where}\alpha \delta \beta \gamma =1$$
$`(45)`$
Thus the most general transformations of the time coordinate form an $`SL(2,R)`$ group. The second equation in (35) then reduces to
$$\stackrel{~}{𝐮}=(\gamma t+\delta )𝐮\gamma (𝐱+𝐚)+\delta 𝐯$$
$`(46)`$
From (41) it also follows that
$$\xi _i=\frac{1}{\gamma t+\delta }(x_i+a_i+v_it)$$
$`(47)`$
Returning to the general case, and using the inversions
$$\gamma \tau \alpha =\frac{1}{\gamma t+\delta }\text{and}\gamma \xi v=\frac{\gamma (x+a)\delta v}{\gamma t+\delta }$$
$`(48)`$
we have
$$\lambda =\lambda _0\frac{(\gamma \xi v)^2}{2\gamma (\gamma \tau \alpha )}=\lambda _0+\frac{[\gamma (x+a)\delta v]^2}{2\gamma (\gamma t+\delta )}\text{where}\lambda _0=\left(h_0+\frac{\delta }{\gamma }\frac{v^2}{2}\right)$$
$`(49)`$
The case $`\text{det}M=f=`$ a constant corresponds to letting $`\gamma =0`$ and from (45) and (47) it is clear that the group of transformations in this case includes the static Galilei group $`G`$, the dilations, the time translations and parity, but excludes time-reversal. As already mentioned, this is the maximal invariance group of the Navier-Stokes equations. It is also useful to consider the following two special cases.
I. $`\beta =\gamma =0,\alpha =1`$: The Connected, Static Galilei Transformations: In this case, we have
$$g:\tau =t,\xi =R𝐱+𝐚+𝐯t$$
$`(50)`$
where we have restored the rotations, and from (49) it follows
$$\lambda =h_0𝐯(R𝐱+𝐚)\frac{v^2}{2}t$$
$`(51)`$
These equations describe connected, static Galilei transformations which exclude parity and time-reversal.
II. $`𝐚=𝐯=\mathrm{𝟎},R=1`$: The Inversion Transformations: In this case, we have
$$\sigma :\tau =\frac{\alpha t+\beta }{\gamma t+\delta },\xi =\frac{R𝐱}{\gamma t+\delta };\alpha \delta \beta \gamma =1$$
$`(52)`$
These are the $`SL(2,R)`$ generalisations of the inversion transformations presented in . For the transformations of the fields we have from (49)
$$\lambda =\frac{\gamma x^2}{2(\gamma t+\delta )}h_0$$
$`(53)`$
The Maximal Invariance Group $`𝒢`$
To understand the structure of the group, we study the relationship between the $`SL(2,R)`$ group and the connected static Galilei group $`G`$. Let us first consider a conjugation of a $`gG`$ by a $`\sigma SL(2,R)`$. By making three successive transformations of $`x`$ and $`t`$ we find that
$$\sigma ^1(\alpha ,\beta ,\gamma )g(R,𝐚,𝐯)\sigma (\alpha ,\beta ,\gamma )=g(R,𝐚_\sigma ,𝐯_\sigma )$$
$`(54)`$
where
$$\left(\begin{array}{c}𝐚_\sigma \\ 𝐯_\sigma \end{array}\right)=\left(\begin{array}{cc}\delta & \beta \\ \gamma & \alpha \end{array}\right)\left(\begin{array}{c}𝐚\\ 𝐯\end{array}\right)$$
$`(55)`$
This shows that $`G`$ is an invariant sub-group and so the group structure is
$$𝒢=SL(2,R)G$$
$`(56)`$
where $``$ denotes semi-direct product with $`G`$ as invariant subgroup. More precisely, if we recall that $`G`$ itself takes the form
$$G=R\left(T(𝐚)B(𝐯)\right)$$
$`(57)`$
where $`T`$ and $`B`$ are the translation and boost groups with parameters $`𝐚`$ and $`𝐯`$ respectively, then we see that $`SL(2,R)`$ commutes with $`R`$ and mixes $`T`$ and $`B`$ in the manner shown in (55).
The original inversion $`\mathrm{\Sigma }`$ is the special element of $`SL(2,R)`$ for which $`(\alpha ,\beta ,\gamma ,\delta )=(0,1,1,0)`$. Note that $`\mathrm{\Sigma }^2=P`$ where $`P`$ is the parity. Furthermore if we consider the coset elements $`g_\genfrac{}{}{0pt}{}{}{\mathrm{\Sigma }}(R,𝐚,𝐯)\mathrm{\Sigma }g(R,𝐚,𝐯)`$, where $`gG`$, we have using (54) and the standard property of Galilei transformations viz. $`g(R,𝐚,𝐯)g(R^{},𝐚^{},𝐯^{})=g(RR^{},R𝐚^{}+𝐚,R𝐯^{}+𝐯)`$,
$$\begin{array}{cc}\hfill g_\genfrac{}{}{0pt}{}{}{\mathrm{\Sigma }}(R^{},𝐚^{},𝐯^{})g_\genfrac{}{}{0pt}{}{}{\mathrm{\Sigma }}(R,𝐚,𝐯)& =g_\genfrac{}{}{0pt}{}{}{P}(R^{}R,R^{}𝐚𝐯^{},R^{}𝐯+𝐚^{})\hfill \\ \hfill g_\genfrac{}{}{0pt}{}{}{\mathrm{\Sigma }}^2(R,𝐚,𝐯)& =g_\genfrac{}{}{0pt}{}{}{P}(R^2,R𝐚𝐯,R𝐯+𝐚)\hfill \end{array}$$
$`(58)`$
where we have used the obvious notation $`g_\genfrac{}{}{0pt}{}{}{P}(R,𝐚,𝐯)=Pg(R,𝐚,𝐯)`$. It follows from the above equation that
$$g_\genfrac{}{}{0pt}{}{}{\mathrm{\Sigma }}^4(R,𝐚,𝐯)=g(R^4,(R^21)(R𝐚𝐯),(R^21)(R𝐯+𝐚))$$
$`(59)`$
Since $`R𝐚𝐯`$ and $`R𝐯+𝐚`$ are linearly independent, this shows that every connected Galilei transformation is the fourth power of a coset transformation.
Extension to the Quantum Theory
If a classical system is described by the Galilei group $`G`$ in (57) then, in the field theoretic representation, there is a central extension of $`G`$; namely, the well-known one-parameter mass group whose generator commutes with all the generators of $`G`$. Since a quantum wavefunction can be thought of as a nonrelativistic field, the corresponding quantum system is described by a central extension of the group of the form
$$G=R(T(𝐚),B(𝐯))$$
$`(60)`$
where the group $`(T,B)`$ is no longer abelian, but is a Heisenberg-Weyl group of the form
$$T(𝐚)B(𝐯)=B(𝐯)T(𝐚)e^{i\mathrm{}M𝐚𝐯}$$
$`(61)`$
and $`M`$ is a central constant. It is obvious that this relation is invariant with respect to the transformation (55) of the parameters $`𝐚`$ and $`𝐯`$ induced by $`SL(2,R)`$. Thus $`𝒢`$ remains a symmetry group of the quantised system.
Connection to the Free Schrödinger Equation
The fact that $`𝒢`$ is a covariance group of the Hamilton-Jacobi function $`\varphi /t+1/2(\varphi )^2`$ shows that it is an invariance group of a free particle. The above discussion of the quantum extension implies that it is also an invariance group of a quantised free particle. Indeed it was shown in 1972 that $`𝒢`$ is the maximal invariance group of the free particle Schrödinger equation . Since the invariance under the Galilei group is well-known, it suffices to verify that the Schrödinger equation
$$i\mathrm{}\frac{\psi }{t}+\frac{\mathrm{}^2}{2m}\frac{^2\psi }{x_i^2}=0$$
$`(62)`$
remains form-invariant under the $`𝒢`$. It is easily checked that this is accomplished by the following transformation of the wavefunction
$$\psi (𝐱,t)(\gamma \tau \alpha )^{\frac{n}{2}}e^{i\lambda }\psi (\xi ,\tau )(\gamma t+\delta )^{\frac{n}{2}}e^{i\lambda }\psi $$
$`(63)`$
where $`\lambda `$ is defined in (49). Note that at time $`\tau =\mathrm{}`$, i.e. $`t=\frac{\delta }{\gamma }`$, the wavefunction becomes infinite, but this is precisely the singularity associated with the explosion.
Conclusions
We have investigated the maximal kinematical invariance group of general fluid mechanics and have found that it is a semi-direct product of the form $`SL(2,R)G`$ where $`G`$ is the static Galilei group. The incompressibility condition is not preserved by the above transformations and hence the viscosity terms transform covariantly only if the viscosity field $`\zeta `$ is constant in space (although not in time). In this case the fluid dynamic equations remain invariant if the viscosity fields transform like tensors of rank $`n`$, where $`n`$ is the space dimension. The Navier-Stokes assumption of constant viscosity breaks the $`SL(2,R)`$ part of $`𝒢`$ to a two-parameter group of dilations and time translations in a tensorial manner.
The inversion transformation $`\mathrm{\Sigma }`$ found in is a special element of the $`SL(2,R)`$ and acts like the square root of parity. The transformations generated by the coset elements $`g_\genfrac{}{}{0pt}{}{}{\mathrm{\Sigma }}=\mathrm{\Sigma }g`$ act like the fourth roots of connected Galilei transformations. It is also pointed out that $`𝒢`$ is the Niederer group which is the maximal invariance group of the free Schrödinger equation. The reason for this is that $`𝒢`$ is actually the invariance group of the free Hamilton-Jacobi function, which is common to both fluid dynamics and the classical limit of the Schrödinger equation. It is surprising that such a basic result does not seem to be more widely known!
Appendix
In this appendix we show that the linear inhomogeneous transformations of the field variables considered in (18) are indeed the most general field transformations allowed. In order to show this we begin by noting that the most general transformations can be written as follows:
$$\stackrel{~}{\varphi }\stackrel{~}{\varphi }(\xi ,\tau ,\varphi )\text{and}\stackrel{~}{\rho }\stackrel{~}{\rho }(\xi ,\tau ,\rho )$$
$`(A1)`$
where $`\stackrel{~}{\varphi }`$ and $`\stackrel{~}{\rho }`$ are a priori arbitrary functions of the old fields $`\varphi `$ and $`\rho `$. Note that we are not allowing the fields to mix, but this is reasonable since $`\varphi `$ and $`\rho `$ are really not on the same footing; whereas $`\varphi `$ appears in the Action (17) only through its derivatives, $`\rho `$ appears without any derivatives. Using the fact that $`x,t\xi ,\tau `$, the Hamilton-Jacobi function can be rewritten as follows:
$$\frac{\varphi }{t}\frac{1}{2}\left(\frac{\varphi }{x}\right)^2=\left(\frac{\xi }{t}\frac{\varphi }{\xi }+\frac{\tau }{t}\frac{\varphi }{\tau }\right)\frac{1}{2}\left(\frac{\xi }{x}\frac{\varphi }{\xi }+\frac{\tau }{x}\frac{\varphi }{\tau }\right)^2$$
$`(A2)`$
Now it follows from
$$\stackrel{~}{\varphi }\stackrel{~}{\varphi }(\xi ,\tau ,\varphi )\varphi F(\xi ,\tau ,\stackrel{~}{\varphi })$$
$`(A3)`$
that
$$\frac{\varphi }{\xi }=\frac{F}{\xi }+\frac{\stackrel{~}{\varphi }}{\xi }\frac{F}{\stackrel{~}{\varphi }}\text{and}\frac{\varphi }{\tau }=\frac{F}{\tau }+\frac{\stackrel{~}{\varphi }}{\tau }\frac{F}{\stackrel{~}{\varphi }}$$
$`(A4)`$
We would now like to fix the form of $`F`$ by requiring the covariance of the Hamilton-Jacobi function i.e.
$$\frac{\varphi }{t}\frac{1}{2}\left(\frac{\varphi }{x}\right)^2\frac{\stackrel{~}{\varphi }}{\tau }\frac{1}{2}\left(\frac{\stackrel{~}{\varphi }}{\xi }\right)^2$$
Substituting $`(A4)`$ on the right hand side of $`(A2)`$, we notice that we get terms of the form $`(\stackrel{~}{\varphi }/\tau )^2`$ unless $`\tau /x=0\tau \tau (t)`$. Equating the coefficients of $`\stackrel{~}{\varphi }/\tau `$ and $`(1/2)(\stackrel{~}{\varphi }/\xi )^2`$ because of covariance, we get
$$\frac{F}{\stackrel{~}{\varphi }}=\frac{\frac{\tau }{t}}{\left(\frac{\xi }{x}\right)^2}$$
$`(A5)`$
Requiring the coefficient of the $`\stackrel{~}{\varphi }/\xi `$ term to vanish produces
$$\frac{F}{\xi }=\frac{\frac{\xi }{t}}{\left(\frac{\xi }{x}\right)^2}$$
$`(A6)`$
Similarly requiring the coefficient of the $`\stackrel{~}{\varphi }`$–independent terms to vanish we get
$$\frac{\xi }{t}\frac{F}{\xi }+\frac{\tau }{t}\frac{F}{\tau }\frac{1}{2}\left(\frac{\xi }{x}\right)^2\left(\frac{F}{\xi }\right)^2$$
Substituting this equation in $`(A6)`$ gives, after some algebra,
$$\frac{F}{\tau }=\frac{1}{2}\frac{\left(\frac{\xi }{t}\right)^2}{\left(\frac{\xi }{x}\right)^2\frac{\tau }{t}}$$
$`(A7)`$
From $`(A5)`$ it follows that since $`\stackrel{~}{\varphi }`$ is arbitrary, the left hand side changes whereas the right hand side, for a given $`\xi ,\tau ,x,t`$, is a constant in the field. Hence
$$\frac{F}{\stackrel{~}{\varphi }}s^1(\xi ,\tau )$$
$`(A8)`$
where $`s^1`$ is a constant in the field variable. Solving the above differential equation we get
$$F=s^1\stackrel{~}{\varphi }+\stackrel{~}{\lambda }(\xi ,\tau )\stackrel{~}{\varphi }=s(\xi ,\tau )\varphi +\lambda (\xi ,\tau )$$
$`(A9)`$
where $`\lambda s\stackrel{~}{\lambda }`$. Further by taking a derivative of $`(A9)`$ with respect to $`\xi `$ we see that
$$\frac{F}{\xi }=\frac{s^1}{\xi }\stackrel{~}{\varphi }+s^1\frac{\stackrel{~}{\varphi }}{\xi }+\frac{\stackrel{~}{\lambda }}{\xi }$$
$`(A10)`$
Again from $`(A6)`$ it is clear that the left hand side of the above equation is independent of $`\stackrel{~}{\varphi }`$ and hence, $`s^1`$ is independent of $`\xi `$. Similarly, by appealing to $`(A7)`$ we can show that it is independent of $`\tau `$. Hence $`s`$ is a constant. A similar analysis may be carried out to establish the generality of the $`\rho `$ transformations in (18).
Acknowledgements
We thank Luke Drury for introducing us to this topic and for many useful discussions, David Saakian for raising the question of viscosity, O. Jahn for his interest, and R. Jackiw for valuable correspondence.
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4. L. O’C Drury and J. T. Mendonça, Physics of Plasmas 7, (2000) 5148.
5. S. Weinberg, Gravitation and Cosmology; John Wiley and Sons Inc. (1972).
6. E. C. G. Sudarshan and N. Mukunda, Classical Dynamics: A Modern Perspective, John Wiley & Sons (1974).
7. For related work in this matter see, M. Hassaine and P. A. Horvathy, Ann. of Phys. 282 (2000) 218; Phys. Lett A279 (2001) 215; A. M. Grundland and L. Lalague, Can. J. Phys. 72, (1994) 362, Can. J. Phys. 73, (1995) 463; R. Jackiw, physics/0010042.
8. U. Niederer, Helvetica Physica Acta, 45 (1972) 802; C. R. Hagen Phys. Rev. D5 (1972) 377; R. Jackiw, Phys. Today 25 (1972) 23.
9. L. D. Landau and E. M. Lifshitz, Fluid Mechanics, Pergamon Press (1959).
10. A. Clebsch, Journal für die reine und angewandte Mathematik, 56 (1859), 1; H. Lamb, Hydrodynamics, Cambridge University Press, 1942. For related work, see S. Deser, R. Jackiw and A. P. Polychronakos, physics/0006056; R. Jackiw, V. P. Nair and So-Young Pi, hep-th/0004084.
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# On 𝑘-abelian 𝑝-filiform Lie algebras
## 1 Generalities
###### Definition 1
Let $`𝔤`$ be a finite dimensional vectorial space over $``$. A Lie algebra law over $`^n`$ is a bilinear alternated mapping $`\mu Hom(^n\times ^n,^n)`$ which satisfies the conditions
1. $`\mu (X,X)=0,X^n`$
2. $`\mu (X,\mu (Y,Z))+\mu (Z,\mu (X,Y))+\mu (Y,\mu (Z,X))=0,X,Y,Z^n`$,
( Jacobi identity )
If $`\mu `$ is a Lie algebra law, the pair $`𝔤=(^n,\mu )`$ is called Lie algebra. From now on we identify the Lie algebra with its law $`\mu `$.
###### Remark 2
We say that $`\mu `$ is the law of $`𝔤`$, and where necessary we use the bracket notation to describe the law :
$$[X,Y]=\mu (X,Y),X,Y𝔤$$
The nondefined brackets are zero or obtained by antisymmetry.
Let $`𝔤_n=(^n,\mu )`$ be a nilpotent Lie algebra. For any nonzero vector $`X𝔤_nC^1𝔤_n`$ let $`c\left(X\right)`$ be the ordered sequence of a similitude invariants for the nilpotent operator $`ad_\mu \left(X\right)\dot{,}`$ i.e., the ordered sequence of dimensions of Jordan blocks of this operator. The set of these sequences is ordered lexicographically.
###### Definition 3
The characteristic sequence of $`𝔤_n`$ is an isomorphism invariant $`c\left(𝔤_n\right)`$ defined by
$$c\left(𝔤_n\right)=\underset{X𝔤_nC^1𝔤_n}{\mathrm{max}}\left\{c\left(X\right)\right\}$$
A nonzero vector $`X𝔤_nC^1𝔤_n`$ for which $`c\left(X\right)=c\left(𝔤_n\right)`$ is called characteristic vector.
###### Remark 4
In particular, the algebras with maximal characteristic sequence $`(n1,1)`$ correspond to the filiform algebras introduced by Vergne \[Ve1\]. Thus it is natural to generalize this concept to lower sequences.
###### Definition 5
A nilpotent Lie algebra $`𝔤_n`$ is called $`p`$-filiform if its characteristic sequence is $`(np,1,..^p..,1).`$
###### Remark 6
This definition was first given in \[CGoJ\]. The $`\left(n1\right)`$-filiform Lie algebras are abelian, while the $`\left(n2\right)`$-filiform are the direct sum of an Heisenberg algebra $`𝔥_{2p+1}`$ and abelian algebras. A classification of the $`\left(n3\right)`$-filiform can also be found in \[CGoJ\]. These and any $`\left(n4\right)`$-filiform Lie algebra have non trivial diagonalizable derivations \[AC1\]. This fact is important, for it is telling us that their structure is relatively simple. To search for nilpotent algebras with rank zero ( i.e, with no nonzero diagonalizable derivations ) we must start with the $`\left(n5\right)`$-filiform Lie algebras. As the difficulty of distinguishing isomorphism classes increases considerably for bigger indexes, it seems reasonable to consider additional assumptions made on the algebras to be classified. For example, the filtration given by the central descending sequence can be used to impose additional conditions on the $`p`$-filiformness.
###### Remark 7
For indexes $`p\left(n4\right)`$ the number of isomorphism classes is finite. The index $`p=\left(n5\right)`$ is the first for which an infinity of isomorphism classes exists.
###### Definition 8
Let $`𝔤_n`$ be a nilpotent Lie algebra. The smallest integer $`k`$ such that the ideal $`C^k𝔤_n`$ is abelian is called commutativity index of $`𝔤_n.`$
###### Definition 9
A Lie nilpotent algebra $`𝔤_n`$ is called $`k`$-abelian if $`k`$ is the smallest positive integer such that
$$C_𝔤\left(C^k𝔤\right)C^k𝔤\text{ and }C^{k1}𝔤C^{k1}𝔤$$
###### Remark 10
The preceding definition is equivalent to impose that the commutativity index of $`𝔤_n`$ is exactly $`k`$. In \[GGoK\] a less restrictive definition of $`k`$-abelianity is considered. The purpose there is to study certain topological properties of the variety of filiform laws $`𝔉^m.`$ Our definition is more restrictive: the $`k`$-abelian Lie algebras do not contain the $`\left(k1\right)`$-abelian algebras; the reason is justified by the important structural difference between algebras having its ideal $`C^k𝔤_n`$ abelian and those having it not. On the other side we avoid unnecessary repetitions.
As we are considering here the $`\left(n5\right)`$-filiform Lie algebras, we have to determine which abelianity indexes are admissible. Only the nonsplit algebras are of interest for us, thus from now on we will understand nonsplit Lie algebra when we say Lie algebra.
###### Lemma 11
Let $`𝔤_n`$ be an $`\left(n5\right)`$-filiform Lie algebra. Then there exists a basis $`\{\omega _1,..,\omega _6,\theta _1,..,\theta _{n6}\}`$ of $`\left(^n\right)^{}`$ such that the law is expressible as
$`d\omega _1`$ $`=d\omega _2=0`$
$`d\omega _3`$ $`=\omega _1\omega _2`$
$`d\omega _4`$ $`=\omega _1\omega _3+{\displaystyle \underset{i=1}{\overset{n6}{}}}\alpha _i^2\theta _i\omega _2`$
$`d\omega _5`$ $`=\omega _1\omega _4+{\displaystyle \underset{i=1}{\overset{n6}{}}}\left(\alpha _i^2\theta _i\omega _3+\alpha _i^3\theta _i\omega _2\right)+\beta _2\omega _3\omega _2`$
$`d\omega _6`$ $`=\omega _1\omega _5+{\displaystyle \underset{i=1}{\overset{n6}{}}}\left(\alpha _i^2\theta _i\omega _4+\alpha _i^3\theta _i\omega _3+\alpha _i^4\theta _i\omega _2\right)`$
$`+{\displaystyle \underset{1i,jn6}{}}a_{ij}^1\theta _i\theta _j+\beta _1\left(\omega _5\omega _2\omega _3\omega _4\right)+\beta _2\omega _4\omega _2+\beta _3\omega _3\omega _2`$
$`d\theta _j`$ $`=\epsilon _i^j\theta _i\omega _2\beta _{1,j}\left(\omega _5\omega _2\omega _3\omega _4\right)+\beta _{3,j}\omega _3\omega _2,\mathrm{\hspace{0.33em}1}jn6`$
The proof is trivial.
###### Lemma 12
Any $`\left(n5\right)`$-filiform Lie algebra $`𝔤_n`$ is either $`1`$ or $`2`$-abelian.
Proof. If the algebra is $`1`$-abelian, it is simply an algebra whose derived algebra is abelian \[Bra\]. If it is $`2`$-abelian, then there exist $`X,YC^1𝔤_n`$ such that $`0[X,Y]`$. From the above equations it is immediate to derive the possibilities:
1. $`dimC^1𝔤_n=6`$
2. $`dimC^1𝔤_n=5`$ and $`X,YC^1𝔤_n`$ such that $`[X,Y]0`$
3. $`dimC^1𝔤_n=4`$ and $`X,YC^1𝔤_n`$ such that $`[X,Y]0.`$
###### Remark 13
From this lemma we see how important is to consider our stronger version of the $`k`$-abelianity. In particular we will see its connection with the characteristic nilpotence.
###### Notation 14
For $`n7`$ let $`𝔤_0^n`$ be the Lie algebra whose Cartan-Maurer equations are
$`d\omega _1`$ $`=d\omega _2=0`$
$`d\omega _j`$ $`=\omega _1\omega _{j1},\mathrm{\hspace{0.33em}3}j6`$
$`d\theta _j`$ $`=0,\mathrm{\hspace{0.33em}1}jn6`$
Let $`\{X_1,..,X_6,Y_1,..,Y_{n6}\}`$ be a dual basis of $`\{\omega _1,..,\omega _6,\theta _1,..,\theta _6\}`$. Let $`V_1=X_1,..,X_6_{}`$ and $`V_2=Y_1,..,Y_{n6}_{}`$. We write $`B(V_i,V_j)`$ to denote the space of bilinear alternated mappings from $`V_i`$ to $`V_j`$.
Let us consider the following applications :
1. $`\psi _{i,j}^1B(V_2,V_1),\mathrm{\hspace{0.33em}1}i,jn6`$ :
$$\psi _{i,j}^1(Y_i,Y_l)=\{\begin{array}{cc}X_6& \text{if }i=k,j=l\\ 0& \text{otherwise}\end{array}$$
2. $`\psi _i^jHom(V_2\times V_1,V_1),j=2,3,4`$
$$\psi _i^j(Y_k,X_l)=\{\begin{array}{cc}X_{l+j}& \text{if }i=k,\mathrm{\hspace{0.33em}2}l6j\\ 0& \text{otherwise}\end{array}$$
3. $`\phi _{1,k}B(V_1,V_2),\mathrm{\hspace{0.33em}1}kn6`$ :
$$\phi _{1,k}(X_5,X_2)=\phi _{1,k}(X_3,X_4)=Y_k$$
4. $`\phi _{3,k}B(V_1,V_2),\mathrm{\hspace{0.33em}1}kn6`$ :
$$\phi _{3,k}(X_3,X_2)=Y_k$$
5. $`\phi _1B(V_1,V_1)`$ :
$$\phi _1(X_5,X_2)=\phi _1(X_3,X_4)=X_6$$
6. $`\phi _2B(V_1,V_1)`$ :
$$\phi _{}(X_k,X_2)=X_{k+2},k=3,4$$
7. $`\phi _3B(V_1,V_1)`$ :
$$\phi _3(X_3,X_2)=X_6$$
where the undefined brackets are zero or obtained by antisymmetry.
###### Lemma 15
For $`n7`$$`1kn6`$ and $`l=2,3,4`$ the mappings $`\psi _{1,k},\psi _{3,k},\psi _i^l,`$
$`\phi _{1,k},\phi _{3,k},\phi _1,\phi _2,\phi _3`$ are $`2`$-cocycles of the subspace $`Z^2(𝔤_0^n,𝔤_0^n)`$.
###### Notation 16
For convenience in the exposition, we introduce the following notation
$$\underset{\begin{array}{c}t=i\\ m>k\end{array}}{\overset{j}{}}\psi _{f\left(t\right),f\left(t\right)+1}$$
where the sum is only defined whenever $`mk+1`$.
###### Proposition 17
Any nonsplit $`\left(n5\right)`$-filiform Lie algebra with $`dimC^1𝔤_n=6`$ is isomorphic to one of the following laws :
1. $`𝔤_{2m}^1\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{1,1}+\phi _{3,2}+\psi _2^3+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
2. $`𝔤_{2m}^2\left(m5\right):`$
$$𝔤_0^{2m}+\phi _{1,1}+\phi _{3,2}+\psi _3^3+\underset{t=2}{\overset{m3}{}}\psi _{2t1,2t}^1$$
3. $`𝔤_{2m+1}^3\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _{1,1}+\phi _{3,2}+\psi _3^3+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
4. $`𝔤_{2m}^4\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{1,1}+\phi _{3,2}+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
5. $`𝔤_{2m}^5\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{1,1}+\phi _2+\phi _{3,2}+\psi _2^3+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
Proof. Suppose $`\beta _{1,k}0`$ for $`k1`$. We can take $`\beta _{1,1}=1`$ and $`\beta _{1,i}=0`$ for $`i2`$, as well as $`\beta _1=0`$. The Jacobi conditions imply
$$\alpha _1^2=\alpha _1^3=\alpha _1^1=0,\alpha _1^4=\underset{k2}{}\beta _{3,k}\alpha _{k,j}^4,\beta _{3,k}\alpha _{i,k}^1=0$$
Let $`\beta _{3,k}0`$ for $`k1`$, so that we can choose $`\beta _{3,2}=1,\beta _{3,k}=0`$ for $`k2`$. A change of basis allows to take $`\beta _3=0`$. From the conditions above we deduce $`\alpha _1^4=\alpha _i^1=0`$. Consider tha change $`\omega _2^{}=\alpha \omega _1+\omega _2`$ with $`\alpha 0`$. Then we have
$$\{\begin{array}{c}\alpha _{j,i}^4=0,j2\\ \alpha _j^3\beta _2=0,j2\end{array}$$
There are two cases :
1. 1. If $`\alpha _2^30`$ we suppose $`\alpha _2^3=1`$ and $`\alpha _i^3=0,i2`$ through a linear change. Reordering the forms $`\theta _i`$ we can suppose $`\alpha _{2t1,2t}^1=1`$ for $`2t\frac{n6}{2};\alpha _{i,j}^1`$ for the remaining. We obtain a unique class of nonsplit Lie algebras in even dimension and isomorphic to $`𝔤_{2m}^1.`$
2. If $`\alpha _2^3=0`$
1. If $`\alpha _i^30`$ with $`i3`$ we can suppose $`\alpha ^3{}_{3}{}^{}=1`$ and $`\alpha _j^3=0`$ for $`j3.`$ Reordering the $`\theta _i`$ we obtain one algebra in even and one algebra in odd dimension, which are respectively isomorphic to $`𝔤_{2m}^2`$ and $`𝔤_{2m+1}^3.`$
2. If $`\alpha _i^3=0`$ for $`i3`$ we obtain in an analogous way two even dimensional algebras, respectively isomorphic to $`𝔤_{2m}^4`$ and $`𝔤_{2m}^5.`$
###### Remark 18
It is very easy to see that the obtained algebras are pairwise non isomorphic, as their infinitesimal deformations are not cohomologous cocycles in the cohomology space $`H^2(𝔤_0^{2m},𝔤_0^{2m})`$. This calculations are routine and will be ommited in future.
###### Remark 19
From the linear system associated ( for the elementary properties of these systems see \[AG1\] ) to the algebras above it follows the existence of nonzero eigenvectors for diagonalizable derivations, so that the rank is at least one. Then the algebra of derivations has nonzero semi-simple derivations.
###### Notation 20
We define the set
$$𝔥_2=\left\{𝔤\right|𝔤\text{ is nonsplit, }2\text{-abelian and }\left(n5\right)\text{-filiform }\}$$
We now express conditions making reference to the reduced system of forms given before :
We say that $`𝔤_2`$ satisfies property $`\left(P1\right)`$ if
$`dimC^1𝔤`$ $`=5`$
$`\beta _{1,1}`$ $`=1`$
###### Remark 21
The general condition would be $`\beta _{1,k}0`$ for a $`k1`$ and $`\beta _{3,k}=0`$ for any $`k`$. Now a elementary change of basis allows to reduce it to the preceding form.
###### Proposition 22
Let $`𝔤_n`$ be an $`\left(n5\right)`$-filiform Lie algebra satisfying the property $`\left(P1\right).`$ The $`𝔤_n`$ is isomorphic to one of the following laws:
1. $`𝔤_{2m+1}^6\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _{1,1}+\psi _2^3+\underset{t=2}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
2. $`𝔤_{2m}^7\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{1,1}+\psi _2^3+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
3. $`𝔤_{2m+1}^8\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _{1,1}+\phi _2+\underset{\begin{array}{c}t=2\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
4. $`𝔤_{2m+1}^9\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _{1,1}+\phi _3+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}$$
5. $`𝔤_{2m+1}^{10}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _{1,1}+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m1}{}}\psi _{2t,2t+1}^1$$
Proof. $`\beta _{3,k}=0`$ for all $`k.`$ The characteristic sequence implies $`\alpha _i^4=\alpha _i^3\beta _3=0`$ for all $`i.`$ Moreover, a change of basis of the type $`\omega _2^{}=\omega _4\frac{j^1}{2}\omega _2`$ allows to suppose $`\alpha _{1,j}^4=0.`$
1. If $`\alpha _i^30`$ we suppose $`\alpha _2^3=1,\alpha _i^3=0,i2.`$ A change of basis allows $`\beta _3=0.`$ There are two possibilities: an even dimensional algebra isomorphic to $`𝔤_{2m}^7`$ and an odd dimensional one isomorphic to $`𝔤_{2m+1}^6.`$
2. $`\alpha _i^3=0,i.`$
1. If $`\beta _20`$ we put $`\beta _2=1`$ and $`\beta _3=0`$ with a linear change of basis. We obtain a unique algebra in odd dimension isomorphic to $`𝔤_{2m+1}^8.`$
2. If $`\beta _2=0`$ there are two possibilities, depending on $`\beta _3`$ : we obtain two odd dimensional algebras which are respectively isomorphic to $`𝔤_{2m+1}^9`$ and $`𝔤_{2m+1}^{10}.`$
A Lie algebra $`𝔤𝔥_2`$ satisfies property $`\left(P2\right)`$ if
$`dimC^1𝔤_n`$ $`=5`$
$`\beta _{3,t}`$ $`0\text{ for t}1`$
###### Proposition 23
Let $`𝔤_n`$ be an algebra with property $`\left(P2\right)`$. Then $`𝔤_n`$ is isomorphic to one of the following laws
1. $`𝔤_{2m+1}^{11}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\psi _1^3+\psi _1^4+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
2. $`𝔤_{2m}^{12}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\psi _1^3+\psi _1^4+\psi _2^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
3. $`𝔤_{2m}^{13}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\psi _1^3+\psi _2^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
4. $`𝔤_{2m+1}^{14}\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\psi _1^3+\underset{t=1}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
5. $`𝔤_{2m+1}^{15}\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\phi _1+\psi _1^4+\psi _2^3+\underset{t=1}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
6. $`𝔤_{2m}^{16}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\phi _1+\psi _1^4+\psi _2^3+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
7. $`𝔤_{2m+1}^{17}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\phi _1+\psi _1^4+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
8. $`𝔤_{2m+1}^{18}\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\phi _1+\psi _3^4+\psi _2^3+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
9. $`𝔤_{2m}^{19}\left(m5\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\phi _1+\psi _{24}^1+\psi _2^3+\psi _3^4+\underset{\begin{array}{c}t=2\\ m>5\end{array}}{\overset{m4}{}}\psi _{2t+1,2t+2}^1$$
10. $`𝔤_{2m}^{20}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\phi _1+\psi _2^3+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
11. $`𝔤_{2m+1}^{21}\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\phi _1+\psi _2^3+\underset{t=1}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
12. $`𝔤_{2m}^{22}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\phi _1+\psi _2^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
13. $`𝔤_{2m+1}^{23}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\phi _1+\underset{t=1}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
14. $`𝔤_{2m+1}^{24,\alpha }\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\alpha \phi _2+\psi _1^3+\psi _1^4+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
15. $`𝔤_{2m}^{25,\alpha }\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\alpha \phi _2+\psi _1^3+\psi _1^4+\psi _2^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
16. $`𝔤_{2m}^{26}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\phi _2+\psi _1^3+\psi _2^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
17. $`𝔤_{2m+1}^{27}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\phi _2+\psi _1^3+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
18. $`𝔤_{2m+1}^{28}\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\phi _1+\phi _2+\psi _2^3+\psi _3^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
19. $`𝔤_{2m}^{29}\left(m5\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\phi _1+\phi _2+\psi _2^3+\psi _3^4+\psi _{24}^1+\underset{\begin{array}{c}t=2\\ m>5\end{array}}{\overset{m4}{}}\psi _{2t+1,2t+2}^1$$
20. $`𝔤_{2m}^{30}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\phi _1+\phi _2+\psi _2^3+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
21. $`𝔤_{2m+1}^{31}\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _{3,1}+\phi _1+\phi _2+\psi _2^3+\underset{t=1}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
22. $`𝔤_{2m}^{32}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _{3,1}+\phi _1+\phi _2+\psi _2^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
Proof. The starting assumptions are $`\beta _{1,k}=0`$ for all $`k`$ and $`\beta _{3,k}0`$ for some $`k1.`$ We can suppose $`\beta _{3,1}=1`$, and from the Jacobi conditions we obtain
$$\alpha _1^2=0,\alpha _{i1}^1=2\alpha _i^2\beta _1,i2$$
From the characteristic sequence we deduce the nullity of $`\alpha _i^3`$ for all $`i,`$ so that $`\alpha _{i1}^1=0`$ $`i`$ as well.
1. $`\alpha _1^3=1`$ : a combination of the linear changes of basis allow to take $`\beta _1=0.`$
1. $`\alpha _1^4=1`$ :.
1. If $`\alpha _i^4=0`$, $`i2`$ we reorder the $`\{\theta _2,..,\theta _{n6}\}`$ such that $`\alpha _{2t1,2t}^1=1`$ for $`1t\frac{n6}{2}`$ and the remaining brackets zero. The decisive structure constant is $`\beta _2.`$ If it is zero we obtain an odd dimensional Lie algebra isomorphic to $`𝔤_{2m+1}^{11}.`$ If not, $`\beta _2=\alpha `$ is an essential parameter. So we obtain an infinite family of odd dimensional Lie algebras isomorphic to the family $`𝔤_{2m+1}^{24,\alpha }.`$
2. $`\alpha _i^40`$, $`i2.`$ Without loss of generality we can choose $`\alpha _2^40`$ and the remaining zero for $`i3.`$ It is easy to deduce $`\alpha _{2j}^1=0,j.`$ Reordering $`\{\theta _3,..,\theta _{n6}\}`$ in the previous manner we obtain an even dimensional Lie algebra and an infinite family of even dimensional algebras, which are respectively isomorphic to $`𝔤_{2m}^{12}`$ and $`𝔤_{2m}^{25,\alpha }.`$
2. Now take $`\alpha _1^4=0`$
1. If there is an index $`i2`$ such that $`\alpha _i^40`$ we can suppose $`\alpha _2^4=1`$ and the remaining zero. Reordering the $`\{\theta _3,..,\theta _{n6}\}`$ as before we obtain two even dimensional Lie algebras, respectively isomorphic to $`𝔤_{2m}^{13}`$ and $`𝔤_{2m}^{26}.`$
2. $`\alpha _i^4=0,i.`$ A similar reordering of the $`\theta _i`$ gives two Lie algebras in odd dimension isomorphic to $`𝔤_{2m+1}^{14}`$ and $`𝔤_{2m+1}^{27}.`$
2. Suppose now $`\alpha _1^3=0.`$
1. $`\alpha _1^4=1,\alpha _i^4=0`$ $`i2.`$ A linear change allows to annihilate $`\beta _2.`$
1. $`\alpha _2^3=1`$ , and the remaining zero.
There are two possible cases, depending on $`\alpha _{23}^1:`$ if it is nonzero we obtain an odd dimensional algebra isomorphic to $`𝔤_{2m+1}^{15}`$ and if it is zero an algebra in even dimension isomorphic to $`𝔤_{2m}^{16}.`$
2. $`\alpha _i^3=0`$ for any $`i2`$ : we obtain an odd dimensional Lie algebra isomorphic to $`𝔤_{2m+1}^{17}.`$
2. $`\alpha _1^4=0`$
1. $`\alpha _2^30`$ and $`\alpha _i^3=0,i3.`$ With a linear change we can suppose $`\alpha _2^4=0.`$
A-1) $`\alpha _3^4=1`$ and $`\alpha _i^4=0`$ for $`i4.`$ A linear change allows to suppose $`\alpha _{3j}^1=0`$ for all $`j.`$ If $`\alpha _{2j}^1=0`$ for all $`j`$ we obtain the algebras $`𝔤_{2m+1}^{18}`$ and $`𝔤_{2m+1}^{28}`$. If not, reorder $`\{\theta _4,..,\theta _{n6}\}`$ such that $`\alpha _{24}^1=1.`$ We obtain the algebras $`𝔤_{2m}^9,𝔤_{2m}^{29}.`$
A-2) $`\alpha _i^4=0`$ $`i.`$ If $`\alpha _{23}^1=0`$ we obtain the Lie algebras $`𝔤_{2m}^{20},𝔤_{2m}^{30},`$ and if $`\alpha _{23}^10`$ we obtain the algebras $`𝔤_{2m+1}^{21}`$ and $`𝔤_{2m+1}^{31}.`$
2. $`\alpha _i^3=0i.`$
B-1) $`\alpha _2^4=1,`$ $`\alpha _{2j}^1=0.`$ We obtain two Lie algebras in even dimension isomorphic to $`𝔤_{2m}^{22}`$ and $`𝔤_{2m}^{32}.`$
B-2) $`\alpha _i^4=0`$ for $`i2`$ : there is only one algebra in odd dimension which is isomorphic to $`𝔤_{2m+1}^{23}.`$
There is only one remaining case, namely the corresponding to the $`\left(n5\right)`$-filiform $`2`$-abelian Lie algebras with minimal dimension of its derived algebra.
###### Proposition 24
Let $`𝔤_n`$ be an $`2`$-abelian algebra with $`dimC^1𝔤_n=4.`$ Then $`𝔤_n`$ is isomorphic to one of the following laws :
1. $`𝔤_{2m}^{33}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _1+\psi _1^3+\psi _2^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
2. $`𝔤_{2m}^{34}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _1+\phi _2+\psi _1^3+\psi _2^4+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
3. $`𝔤_{2m+1}^{35}\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _1+\psi _1^3+\psi _2^4+\psi _{13}^1+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
4. $`𝔤_{2m+1}^{36}\left(m4\right):`$
$$𝔤_0^{2m+1}+\phi _1+\phi _2+\psi _1^3+\psi _2^4+\psi _{13}^1+\underset{\begin{array}{c}t=2\\ m>4\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
5. $`𝔤_{2m+1}^{37}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _1+\psi _1^3+\psi _1^4+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
6. $`𝔤_{2m}^{38}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _1+\psi _1^3+\underset{t=1}{\overset{m3}{}}\psi _{2t1,2t}^1$$
7. $`𝔤_{2m}^{39}\left(m4\right):`$
$$𝔤_0^{2m}+\phi _1+\phi _2+\psi _1^3+\underset{t=1}{\overset{m3}{}}\psi _{2t1,2t}^1$$
8. $`𝔤_{2m+1}^{40}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _1+\psi _1^3+\underset{t=1}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
9. $`𝔤_{2m+1}^{41}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _1+\phi _2+\psi _1^3+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
10. $`𝔤_{2m+1}^{42}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _1+\psi _1^4+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
11. $`𝔤_{2m+1}^{43}\left(m3\right):`$
$$𝔤_0^{2m+1}+\phi _1+\phi _2+\psi _1^4+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t,2t+1}^1$$
12. $`𝔤_{2m}^{44}\left(m3\right):`$
$$𝔤_0^{2m}+\phi _1+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
13. $`𝔤_{2m}^{45}\left(m3\right):`$
$$𝔤_0^{2m}+\phi _1+\phi _2+\underset{\begin{array}{c}t=1\\ m>3\end{array}}{\overset{m3}{}}\psi _{2t1,2t}^1$$
Proof. The starting assumptions for this case are
$$\beta _10\text{ and }\beta _{3,k}=0,k1$$
In particular the Jacobi condition forces $`\alpha _i^2=0`$ $`i.`$
1. Suppose $`\beta _1`$ and $`\alpha _1^30`$ $`\left(\text{so }\alpha _j^3=0\text{ for }i2\text{ }\right).`$
We observe that if there exists an $`\alpha _{ij}^10`$ then a linear change of basis allows to suppose $`\alpha _i^4=\alpha _j^4=0.`$ So we have the conditions
$$d_ia_{ij}=d_ja_{ij}=0,\mathrm{\hspace{0.33em}1}i,j$$
(1)
1. $`\alpha _i^40`$ with $`i2`$ .We can suppose $`\alpha _2^4=1`$ $`\left(\text{so }\alpha _{2j}^1=0\text{ by }\left(1\right)\right)`$ and $`\alpha _i^4=0,i2.`$
1. If $`\alpha _{1j}^1=0`$ for all $`j`$ we obtain two even dimensional algebras isomorphic to $`𝔤_{2m}^{33}`$ and $`𝔤_{2m}^{34}.`$
2. If $`\alpha _{1j}^10`$ for an index $`j`$ we can suppose $`\alpha _{13}^1=1.`$ We obtain two odd dimensional algebras isomorphic respectively to $`𝔤_{2m+1}^{35}`$ and $`𝔤_{2m+1}^{36}.`$
2. $`\alpha _i^4`$ $`=0,i2`$
1. If $`\alpha _1^40,`$ then $`\alpha _{1j}^1=0`$ by $`\left(1\right).`$ Reordering $`\{\theta _2,..,\theta _{n6}\}`$ we obtain an algebra isomorphic to $`𝔤_{2m+1}^{37}.`$
2. If $`\alpha _1^4=0`$ and $`\alpha _{12}^10`$ we obtain two algebras isomorphic to $`𝔤_{2m}^{38}`$ and $`𝔤_{2m}^{39}.`$
3. If $`\alpha _1^4=\alpha _{1j}^1=0,j`$ we obtain two algebras in odd dimension isomorphic to $`𝔤_{2m+1}^{40}`$ and $`𝔤_{2m+1}^{41}`$.
2. Suppose $`\beta _10`$ and $`\alpha _i^3=0,i.`$ Additionally we can suppose $`\beta _3=0.`$
1. If $`\alpha _i^40`$ for an index $`i1`$ let $`\alpha _1^4=1`$ and $`\alpha _i^4=0,i2`$ and $`\alpha _{1j}^1=0`$ by $`\left(1\right).`$ We obtain two algebras isomorphic $`𝔤_{2m+1}^{42}`$ and $`𝔤_{2m+1}^{43}.`$
2. If $`d_i=0`$ $`i`$ we obtain two even dimensional algebras respectively isomorphic to $`𝔤_{2m}^{44}`$ and $`𝔤_{2m}^{45}.`$
###### Remark 25
Observe that the algebras $`𝔤_{2m+1}^{42}`$ and $`𝔤_{2m+1}^{43}`$ are central extensions of the five dimensional filiform algebras $`𝔩_5^1`$ and $`𝔩_5^2.`$
###### Corollary 26
Any nonsplit $`\left(n5\right)`$-filiform $`2`$-abelian Lie algebra is isomorphic to one of the laws $`𝔤^i,i\{1,..,45\}.`$
###### Remark 27
As we have seen that a $`\left(n5\right)`$-filiform Lie algebra is either $`1`$\- or $`2`$-abelian, the global classification follows from determining the isomorphism classes of the $`1`$-abelian ones.
### 1.1 Characteristically nilpotent $`\left(n5\right)`$-filiform Lie algebras
The first example of a nilpotent Lie algebra all whose derivations are nilpotent was given by Dixmier and Lister in 1957 \[DL\], as an answer to a question formulated by Jacobson \[Ja\] two years earlier. This new class of Lie algebras was soon recognized to be very important, and called characteristically nilpotent, as they verify a certain sequence for derivations which is a kind of generalization of the central descending sequence for nilpotent Lie algebras ( \[DL\], \[LT\]).
###### Definition 28
A Lie algebra $`𝔤`$ is called characteristically nilpotent if the Lie algebra of derivations $`Der\left(𝔤\right)`$ is nilpotent.
###### Remark 29
It is easily seen that the original definition given by Dixmier and Lister is equivalent to the given above \[LT\].
###### Remark 30
It is trivial to verify that there do not exist characteristically nilpotent, $`\left(np\right)`$-filiform Lie algebras for indexes $`p=1,2`$. For $`p=3,4`$, it has been shown that these algebras have rank $`r1`$ \[AC1\], and that almost any of these laws is the nilradical of a solvable, rigid law.
###### Lemma 31
Let $`𝔤`$ be a $`\left(n5\right)`$-filiform, $`1`$-abelian Lie algebra. Then $`rank\left(𝔤\right)1`$.
Proof. If $`dimC^1𝔤=4`$, the assertion follows immediately from the linear system $`\left(S\right)`$ associated to the algebra, as this system admits nontrivial solutions. If $`dimC^1𝔤=5`$, the only case for which the system could have zero solution is $`\beta _2=\beta _{3,1}=1`$, and the distinct values of $`(\alpha _i^2,\alpha _i^3,\alpha _i^4)`$. For any of these starting conditions it is routine to prove the existence of a nonzero semisimple derivation.
###### Lemma 32
Let $`𝔤_2`$. If $`\alpha _{ij}^10`$ for $`1i,jn6`$ such that
$$\beta _{1,k}=\beta _{3,k}=\alpha _k^t=0,k=i,j,t=2,3,4$$
then $`rank\left(𝔤\right)1.`$
Proof. Consider the endomorphism defined by
$$d\left(Y_i\right)=Y_i,d\left(Y_j\right)=Y_j$$
and zero over the undefined images, where $`(X_1,..,X_6,Y_1,..,Y_{n6})`$ is the dual basis of $`(\omega _1,..,\omega _6,\theta _1,..,\theta _{n6})`$. Clearly $`d`$ is a nonzero semisimple derivation of $`𝔤`$.
###### Proposition 33
A $`\left(n5\right)`$-filiform Lie algebra $`𝔤_n`$ is characteristically nilpotent if and only if it is isomorphic to one of the following laws:
$`𝔤_7^{11},𝔤_9^{15},𝔤_7^{17},𝔤_7^{24,\alpha }\left(\alpha 0\right),𝔤_7^{27},𝔤_9^{36},𝔤_7^{37},𝔤_7^{41}`$
$`𝔤_8^{12},𝔤_8^{16},𝔤_8^{25,\alpha }\left(\alpha 0\right),𝔤_8^{26},𝔤_8^{34},𝔤_8^{39}`$
###### Corollary 34
There are characteristically nilpotent Lie algebras $`𝔤_n`$ with nilpotence index $`5`$ for the dimensions $`n=7,8,9,14,15,16,17,18`$ and $`n21.`$
Proof. As the sum of characteristically nilpotent algebras is characteristically nilpotent \[T1\], the assertion follows from the previous proposition.
###### Remark 35
In fact, for any $`n7`$ there exist characteristically nilpotent Lie algebras of nilindex $`5`$. However, the algebras to be added are not $`p`$-filiform any more \[AC2\].
### 1.2 Nilradicals of rigid algebras as factors of $`k`$-abelian Lie algebras
The second application of $`k`$-abelian Lie algebras is of interest for the theory of rigid Lie algebras. In this paragraph we prove the existence, by giving a family for dimensions $`2m+2\left(m4\right)`$, of $`\left(m1\right)`$-abelian Lie algebras $`𝔤`$ of characteristic sequence $`(2m1,2,1)`$ all whose factor algebras $`\frac{𝔤}{C^k𝔤}`$ $`\left(km\right)`$ are isomorphic to the nilradical of a solvable rigid law.
For $`m4`$ let $`𝔤_m`$ be the Lie algebra whose Cartan-Maurer equations are
$`d\omega _1`$ $`=d\omega _2=0`$
$`d\omega _j`$ $`=\omega _1\omega _{j1},\mathrm{\hspace{0.33em}3}j2m1`$
$`d\omega _{2m}`$ $`=\omega _1\omega _{2m1}+{\displaystyle \underset{j=2}{\overset{m}{}}}\left(1\right)^j\omega _j\omega _{2m+1j}`$
$`d\omega _{2m+1}`$ $`=\omega _2\omega _3`$
$`d\omega _{2m+2}`$ $`=\omega _1\omega _{2m+1}+\omega _2\omega _4`$
It is elementary to verify that this algebra has characteristic sequence $`(2m1,2,1)`$. Moreover, it is $`\left(m1\right)`$-abelian, for the exterior product $`\omega _m\omega _{m+1}`$ proves that $`[C^{m2}𝔤_m,C^{m2}𝔤_m]0`$ and $`[C^{m1}𝔤_m,C^{m1}𝔤_m]=0.`$
###### Notation 36
The dual basis of $`(\omega _1,..,\omega _{m2+2})`$ will be denoted as $`(X_1,..,X_{2m+2})`$.
###### Lemma 37
For any $`4mk2m2`$ the factor algebra $`\frac{𝔤_m}{C^k𝔤_m}`$ has equations
$`d\stackrel{-}{\omega }_1`$ $`=d\stackrel{-}{\omega }_2=0`$
$`d\stackrel{-}{\omega }_j`$ $`=\stackrel{-}{\omega }_1\stackrel{-}{\omega }_{j1},\mathrm{\hspace{0.33em}3}jk+1`$
$`d\stackrel{-}{\omega }_{2m+1}`$ $`=\stackrel{-}{\omega }_2\stackrel{-}{\omega }_3`$
$`d\stackrel{-}{\omega }_{2m+2}`$ $`=\stackrel{-}{\omega }_1\stackrel{-}{\omega }_{2m+1}+\stackrel{-}{\omega }_2\stackrel{-}{\omega }_4`$
where $`\stackrel{-}{\omega _j}=\omega _j\mathrm{mod}C^k𝔤_m`$. Moreover, this algebra is $`1`$-abelian of characteristic sequence $`(k,2,1)`$.
The proof is trivial.
###### Proposition 38
For any $`4mk`$ the algebra $`\frac{𝔤_m}{C^k𝔤_m}`$ is isomorphic to the nilradical of a solvable, rigid Lie algebra $`𝔯_{m,k}`$.
Proof. Let $`𝔯_{m,k}=`$ $`\frac{𝔤_m}{C^k𝔤_m}𝔱_k`$ be the semidirect product of $`\frac{𝔤_m}{C^k𝔤_m}`$ by the torus $`𝔱_k`$ defined by its weights :
| $`\lambda _1,\lambda _2+\left(k1\right)\lambda _1,\lambda _j=\lambda _2+\left(k3+j\right)\lambda _1\left(3jk+1\right)`$ |
| --- |
| $`\lambda _{2m+1}=2\lambda _2+\left(2k1\right)\lambda _1,\lambda _{2m}=2\lambda _2+2k\lambda _1`$ |
over the basis $`\{\stackrel{-}{X}_1,..,\stackrel{-}{X}_{k+1},\stackrel{-}{X}_{2m+1},\stackrel{-}{X}_{2m+2}\}`$ dual to $`\{\stackrel{-}{\omega }_1,..,\stackrel{-}{\omega }_{k+1},\stackrel{-}{\omega }_{2m+1},\stackrel{-}{\omega }_{2m+2}\}`$. Then the law is given by
| $`[V_1,\stackrel{-}{X}_1]=\stackrel{-}{X}_1,[V_1,\stackrel{-}{X}_j]=\left(k1\right)\stackrel{-}{X}_j\left(3jk+1\right),`$ |
| --- |
| $`[V_1,\stackrel{-}{X}_{2m+1}]=\left(2k1\right)\stackrel{-}{X}_{2m+1},[V_1,\stackrel{-}{X}_{2m+2}]=2k\stackrel{-}{X}_{2m+2}`$ |
| $`[V_2,\stackrel{-}{X}_j]=\stackrel{-}{X}_j\left(2jk+1\right),[V_2,\stackrel{-}{X}_j]=2\stackrel{-}{X}_j\left(j=2m+1,2m+2\right)`$ |
| $`[\stackrel{-}{X}_1,\stackrel{-}{X}_j]=\stackrel{-}{X}_{j+1}\left(2jk\right),[\stackrel{-}{X}_2,\stackrel{-}{X}_3]=a\stackrel{-}{X}_{2m+1},[\stackrel{-}{X}_2,\stackrel{-}{X}_4]=b\stackrel{-}{X}_{2m+1}`$ |
Thus the only nonzero brackets not involving the vector $`\stackrel{-}{X}_1`$ are
$$[\stackrel{-}{X}_2,\stackrel{-}{X}_3]=a\stackrel{-}{X}_{2m+1},[\stackrel{-}{X}_2,\stackrel{-}{X}_4]=b\stackrel{-}{X}_{2m+1}$$
Now Jacobi implies $`a=b`$, and by a change of basis $`a=1`$. Thus the law is rigid, and $`𝔱_k`$ is a maximal torus of derivations of $`\frac{𝔤_m}{C^k𝔤_m}`$.
###### Corollary 39
For any $`4mk`$ the factor algebra
$$\frac{\left(\frac{𝔤_m}{C^k𝔤_m}\right)}{\stackrel{-}{X}_{2m+2}}$$
is $`1`$-abelian of characteristic sequence $`(k,1,1)`$ and isomorphic to the nilradical of a solvable rigid law $`𝔰_{m,k}`$. Moreover
$$𝔰_{m,k}\frac{𝔯_{m,k}}{\stackrel{-}{X}_{2m+2}}$$
## 2 Other 2-abelian nilpotent Lie algebras
Let $`E_6`$ be the simple exceptional Lie algebra of dimension $`78`$. Let $`\mathrm{\Phi }`$ be a root system respect to a Cartan subalgebra $`𝔥`$ and $`\mathrm{\Delta }=\{\alpha _1,..,\alpha _6\}`$ a basis of fundamental roots. Recall that the standard Borel subalgebra of $`E_6`$ is given by
$$B\left(\mathrm{\Delta }\right)=𝔥+\underset{\alpha \mathrm{\Phi }^+}{}L_\alpha $$
where $`L_\alpha `$ is the weight space associated to the root $`\alpha `$. Recall also that any parabolic subalgebra $`𝔭`$ is determined, up to isomorphism, by a subsystem $`\mathrm{\Delta }_1\mathrm{\Delta }`$ such that $`P`$ is conjugated to the subalgebra
$$P\left(\mathrm{\Delta }_1\right)=𝔥+\underset{\alpha \mathrm{\Phi }_1\mathrm{\Phi }_2^+}{}L_\alpha $$
where $`\mathrm{\Phi }_1`$ is the set of roots expressed in terms of $`\mathrm{\Delta }\backslash \mathrm{\Delta }_1`$ and $`\mathrm{\Phi }_2^+=\mathrm{\Phi }^+\left(\mathrm{\Phi }\backslash \mathrm{\Phi }_1\right)`$. It is elementary to see that the nilradical is
$$𝔫\left(\mathrm{\Delta }_1\right)=\underset{\alpha \mathrm{\Phi }_2^+}{}L_\alpha $$
and called $`(E_6,\mathrm{\Delta }_1)`$-nilalgebra.
Let
$$=\left\{\begin{array}{c}\{\alpha _1,\alpha _4\},\{\alpha _4,\alpha _6\},\{\alpha _3,\alpha _5\},\{\alpha _3,\alpha _4\},\{\alpha _4,\alpha _5\},\{\alpha _2,\alpha _3\},\\ \{\alpha _2,\alpha _5\},\{\alpha _2,\alpha _4\},\{\alpha _1,\alpha _2,\alpha _3\},\{\alpha _2,\alpha _5,\alpha _6\},\\ \{\alpha _1,\alpha _2,\alpha _5\},\{\alpha _2,\alpha _3,\alpha _6\},\{\alpha _1,\alpha _4,\alpha _6\},\\ \{\alpha _1,\alpha _2,\alpha _6\},\{\alpha _1,\alpha _3,\alpha _5\},\{\alpha _3,\alpha _5,\alpha _6\}\end{array}\right\}$$
As known, the maximal root of $`E_6`$ is
$$\delta =\alpha _1+2\alpha _2+2\alpha _3+3\alpha _4+2\alpha _5+\alpha _6=\underset{i=1}{\overset{6}{}}k_{a_i}\alpha _i$$
We define the $`\mathrm{\Delta }_1`$-height of $`\delta `$ as
$$h_{\mathrm{\Delta }_1}\left(\delta \right)=\underset{\alpha _i\mathrm{\Delta }_1}{}k_{a_i}$$
and the subsets
$$\mathrm{\Delta }_1\left(k\right)=\left\{\alpha \mathrm{\Phi }_2^+\right|h_{\mathrm{\Delta }_1}\left(\alpha \right)=k\}$$
###### Proposition 40
For any $`\mathrm{\Delta }_1`$ the $`(E_6,\mathrm{\Delta }_1)`$-nilalgebra $`𝔫\left(\mathrm{\Delta }_1\right)`$ is $`2`$-abelian.
Proof. As known, for the ideals $`C^k𝔫`$ of the descending central sequence we have
$$C^k𝔫=\underset{\begin{array}{c}\alpha \mathrm{\Delta }_1\left(j\right)\\ jk+1\end{array}}{}L_\alpha $$
Thus, if the derived subalgebra is not abelian, it suffices to show the existence of two roots $`\alpha ,\beta \mathrm{\Delta }_1\left(2\right)`$ such that $`\alpha +\beta \mathrm{\Phi }_2^+`$ and that for any two roots $`\gamma ,\epsilon \mathrm{\Delta }_1\left(3\right)`$ we have $`\gamma +\epsilon \mathrm{\Phi }_2^+.`$ Moreover, let $`\delta _1=_{i=1}^6\alpha _i\mathrm{\Phi }`$
1. $`\mathrm{\Delta }_1=\{\alpha _1,\alpha _4\}`$ : take $`\alpha =\delta _1\alpha _5\alpha _6,\beta =\delta \alpha _1\alpha _2\alpha _4;\alpha +\beta =\delta `$
2. $`\mathrm{\Delta }_1=\{\alpha _3,\alpha _5\}`$ : $`\alpha =\delta _1,\beta =\delta _1\alpha _1\alpha _6+\alpha _4;\alpha +\beta =\delta `$
3. $`\mathrm{\Delta }_1=\{\alpha _4,\alpha _5\}`$ : $`\alpha =\delta _1\alpha _1\alpha _2\alpha _6,\beta =\delta _1;\alpha +\beta =\delta \alpha _2\alpha _4`$
4. $`\mathrm{\Delta }_1=\{\alpha _2,\alpha _3\}`$ : $`\alpha =\delta _1,\beta =\delta _1\alpha _1\alpha _2\alpha _6;\alpha +\beta =\delta \alpha _2\alpha _4`$
5. $`\mathrm{\Delta }_1=\{\alpha _2,\alpha _4\}`$ : $`\alpha =\delta _1,\beta =\delta _1\alpha _1\alpha _6+\alpha _4;\alpha +\beta =\delta `$
6. $`\mathrm{\Delta }_1=\{\alpha _1,\alpha _2,\alpha _3\}`$ : $`\alpha =\delta \delta _1,\beta =\delta _1\alpha _2;\alpha +\beta =\delta \alpha _2`$
7. $`\mathrm{\Delta }_1=\{\alpha _1,\alpha _2,\alpha _5\}`$ : $`\alpha =\delta \delta _1,\beta =\delta _1\alpha _2;\alpha +\beta =\delta \alpha _2`$
8. $`\mathrm{\Delta }_1=\{\alpha _1,\alpha _3,\alpha _6\}`$ : $`\alpha =\delta _1\alpha _6,\beta =\delta _1\alpha _1\alpha _2;\alpha +\beta =\delta \alpha _2\alpha _4`$
9. $`\mathrm{\Delta }_1=\{\alpha _1,\alpha _4,\alpha _6\}`$ : $`\alpha =\delta _1\alpha _6,\beta =\delta _1\alpha _1\alpha _2;\alpha +\beta =\delta \alpha _2\alpha _4`$
10. $`\mathrm{\Delta }_1=\{\alpha _1,\alpha _2,\alpha _6\}`$ : $`\alpha =\delta _1\alpha _1,\beta =\delta _1+\alpha _4;\alpha +\beta =\delta `$
11. $`\mathrm{\Delta }_1=\{\alpha _1,\alpha _3,\alpha _5\}`$ : $`\alpha =\delta _1\alpha _2\alpha _6,\beta =\delta _1+\alpha _4;\alpha +\beta =\delta \alpha _2`$
Let $`X_\alpha ,X_\beta `$ be generators of the weight spaces $`L_\alpha `$ and $`L_\beta `$ : then the preceding relations show that
$$[X_\alpha ,X_\beta ]=X_{\alpha +\beta }0$$
proving that the derived subalgebra is not abelian.
Finally, it is trivial to see that for any subset $`\mathrm{\Delta }_1`$ listed above we have
$$[C^2𝔫\left(\mathrm{\Delta }_1\right),C^2𝔫\left(\mathrm{\Delta }_1\right)]=0$$
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# HIGGS SEARCHES AT THE TEVATRON Talk given at Les Rencontres de Physique de la Vallée d’Aoste, La Thuile, Italy; February 27 - March 4, 2000
## 1 Introduction
Despite increasingly precise scrutiny, the standard model (SM) of particle physics remains largely unshaken. Only the recent experimental evidence for massive neutrinos challenges any of its many predictions. With the single exception of the Higgs boson, all the particles expected in the SM, and no others, have now been observed. This single missing particle represents a major gap in our knowledge of the microscopic world. Tied up in its nature (or natures, in the case that there is more than one Higgs scalar) is the method by which the original $`SU(2)_L`$$`\times `$$`U(1)`$ symmetry of the theory is spontaneously broken to the distinct electromagnetic and weak forces we observe.
Within the SM, the mass of the single physical Higgs boson left after electroweak symmetry breaking is related to the vacuum expectation value of the neutral Higgs field ($`v`$ = 246 GeV) by $`M_h^2=\lambda v^2`$. The Higgs self-coupling parameter, $`\lambda `$, is not specified by the theory, making the Higgs mass an unknown quantity. However, if the Higgs mechanism is to fulfill its role in the SM, the Higgs mass cannot exceed 1 TeV, otherwise unitarity would be violated in the scattering of longitudinally polarized gauge bosons.
Since it does not incorporate gravity, the SM cannot be a fundamental theory of all interactions despite its many other successes. Even if it is a correct effective theory up to the Planck scale of $``$10<sup>19</sup> GeV, where quantum gravitational effects become significant, the SM is still unsatisfying. This is because radiative corrections to the square of the Higgs mass are quadratically divergent in an energy cutoff parameter, which should be the Planck scale if the SM is to be valid up to that range. In order to get a physical Higgs mass on the order of the electroweak scale, these corrections must be cancelled by tuning the bare Higgs mass at the Planck scale to one part in 10<sup>16</sup> – a rather unnatural condition.
Since there is an implied need for new physics beyond the SM, a natural question to ask is at what energy scale ($`\mathrm{\Lambda }`$) should such effects become apparent. The physics would then look SM-like below $`\mathrm{\Lambda }`$, but new particles and interactions would become apparent beyond that scale. This question is intimately related to the mass of the Higgs boson that is used to break electroweak symmetry. At large values of $`M_h`$, the renormalization group equation for the Higgs self-coupling causes $`\lambda `$ to blow up for $`\mathrm{\Lambda }`$ below the Planck mass. The point at which this happens then sets the scale for new physics. On the other hand, if $`M_h`$ is too small, top quark contributions to $`\lambda `$ can drive it negative. To avoid this problem an energy cutoff must be introduced that can be associated with $`\mathrm{\Lambda }`$. Using these constraints, a measurement of the Higgs mass at the relatively low energies of today’s accelerators could clarify the scale for the breakdown of the SM.
Many theoretical frameworks have been proposed to address some of the weaknesses of the SM without abandoning its low energy successes. These will not be discussed here. We will limit our discussion to results of the minimal supersymmetric standard model (MSSM), a guide indicative of extensions that probe the nature of physics and the Higgs beyond the SM.
The MSSM is the simplest version of supersymmetric models that solve the “naturalness” problem discussed above by relating fermionic and bosonic degrees of freedom. This involves adding supersymmetric partners to all SM particles that differ from their SM counterparts by a half unit of spin. An additional requirement is the expansion of the Higgs sector to include more than the single complex doublet used in the SM. The MSSM proposes only a single extra complex doublet. After electroweak symmetry breaking, five physical Higgs particles remain: two CP-even neutral scalars, $`h`$ and $`H`$ (with $`M_h`$$`<`$$`M_H`$ by convention), one CP-odd neutral scalar $`A`$ and two charged scalars $`H^\pm `$.
The parameters of the MSSM that most directly affect the Higgs bosons are $`\mathrm{tan}\beta `$ (the ratio of the vacuum expectation values of the two Higgs doublets) and $`M_A`$ (the mass of the CP-odd Higgs, conventionally chosen to be the free Higgs mass in the theory). Unlike the case of the SM, in the MSSM, the mass of the lightest of the Higgs bosons, $`h`$, is constrained by supersymmetry. At tree level, this mass must be less than that of the Z<sup>0</sup>. Radiative corrections modify this relationship, but an upper bound still exists with ($`M_h`$)$`{}_{\mathrm{max}}{}^{}`$ 130 GeV.
The mass of an SM-like Higgs is also constrained by experimental measurements. Direct searches at LEP 2 currently give the most stringent lower bound on the mass. This bound changes as LEP accumulates more data. A snapshot of the LEP-wide 95% C.L. limit given at the September 1999 LEPC constrains $`M_h`$ $`>`$ 102.6 GeV. The best experimental upper bound on the Higgs mass comes from global fits to electroweak measurements done by the LEP Electroweak Working Group. Their results indicate that $`M_h`$ $`<`$ 215 GeV at 95% C.L.
## 2 The Higgs at the Tevatron
If a Higgs boson is not discovered at LEP 2 in its last year of running, the next place to look will be the Tevatron at Fermilab. The Tevatron is a $`\overline{\text{p}}`$p collider that operated until 1996 at a center of mass energy of $`\sqrt{s}`$ = 1.8 TeV. Two experiments, CDF and , each collected approximately 100 pb<sup>-1</sup> of data during the period 1992–1996. As we will see, this data set, referred to as “Run I” data, is not sufficiently large to have competitive sensitivity to an SM Higgs although interesting studies have been made in certain non-standard models.
Not content with the success of Run I, the Tevatron accelerator is being upgraded to increase its center-of-mass energy to $`\sqrt{s}`$ = 2.0 TeV and to ultimately achieve an instantaneous luminosity of 2$`\times `$10<sup>32</sup> cm<sup>-2</sup>s<sup>-1</sup>. CDF and are also being upgraded to take advantage of the higher luminosity. Data taking for this new “Run II” will start in March 2001. Projections for integrated luminosity in Run II begin at 2 fb<sup>-1</sup>, but could reach another factor of ten. Obviously, with this large expected data set, CDF and have a bright future.
Before we dive into the details of Run I Higgs searches and projections for Run II sensitivities, it is worthwhile to review how Higgs bosons are produced and decay, and how they would be detected at the Tevatron. The cross section for SM Higgs bosons in $`\overline{\text{p}}`$p collisions, as calculated by Spira, is given in Fig. 2 for various production modes. The main SM Higgs decay modes, calculated using the program HDECAY, are given in Fig. 2 as a function of Higgs mass. Corresponding plots for the neutral MSSM Higgs bosons, $`h`$, $`H`$ and $`A`$, also based on the calculations by Spira, are available on the Fermilab Run II Higgs Working Group web page. These have largely similar characteristics to the plots for SM Higgs, however, at large $`\mathrm{tan}\beta `$ couplings to $`b`$-quarks and $`\tau `$-leptons are enhanced, leading to $`h`$($`H`$)$`b`$$`\overline{b}`$ production being favored over the $`h`$($`H`$)W/Z modes.
The clear advantage of the Tevatron over e<sup>+</sup>e<sup>-</sup> machines is its center of mass energy – a factor of 10 higher than at LEP. Of course, not all of this energy is available to the partons participating in the hard scattering producing the Higgs. Nevertheless, higher mass Higgs bosons can be produced at the upgraded Tevatron than at any previous machine. The problem is finding them.
Table 1 spells out the difficulty in identifying a Higgs at the Tevatron – the immense background. Clearly, the most favorable production mode for the SM Higgs, gg$``$$`h_{SM}`$ (through intermediate top quarks), cannot be used in a search for a low mass Higgs ($`M_h`$$`<`$130 GeV), that decays mainly to $`b`$$`\overline{b}`$. Background from QCD dijet production is a factor of $``$10<sup>6</sup> larger than the signal. This has prompted Tevatron Higgs hunters to concentrate on the next highest cross section production modes of a light Higgs in association with a W or Z. Even here, searches must contend with difficult background from boson-pair, W/Z$`b`$$`\overline{b}`$ and top quark production.
If the Higgs has higher mass, identification is slightly easier although the production cross section is smaller because decays to WW and ZZ begin to dominate above $`M_h`$$``$130 GeV. In fact, these unique final states, make it possible to take advantage of the gg$``$$`h_{SM}`$ production mode in this mass region.
Regardless of the mass of the Higgs, searches at the Tevatron will take advantage of all distinguishing features of Higgs decay to overcome backgrounds with cross-sections many times larger than those of the Higgs. The most promising final states for a low mass SM Higgs (that decays mainly to $`b`$$`\overline{b}`$) are: $`\mathrm{}\nu `$$`b`$$`\overline{b}`$, $`\mathrm{}^+\mathrm{}^{}`$$`b`$$`\overline{b}`$, $`\nu \overline{\nu }`$$`b`$$`\overline{b}`$ and $`jj`$$`b`$$`\overline{b}`$. For a high mass SM Higgs (decaying mainly to WW/ZZ) the best final states contain the following distinguishing particles: $`\mathrm{}^+\mathrm{}^{}`$$`\nu \overline{\nu }`$, $`\mathrm{}^\pm \mathrm{}^\pm jj`$ and $`\mathrm{}^\pm \mathrm{}^{}_{}{}^{}\pm \mathrm{}^{}`$. Discussion of the motivation for choosing these final states can be found in Run II Higgs Working Group Report.
A few themes emerge from consideration of Table 1 and the final states listed above. To be maximally sensitive to the Higgs, an ideal Tevatron detector must have all of the following properties:
1. High lepton (e and $`\mu `$) identification efficiency,
2. Excellent missing energy resolution (for neutrino identification),
3. High $`b`$-quark identification efficiency,
4. Good invariant mass resolution for $`b`$$`\overline{b}`$ pairs (to reject $`b`$$`\overline{b}`$ background outside of the $`h_{SM}`$$``$$`b`$$`\overline{b}`$ peak).
The same signatures mentioned for the SM Higgs will also be important for finding neutral Higgs scalars in the MSSM. Since couplings to $`b`$-quarks and $`\tau `$-leptons tend to grow with increasing $`\mathrm{tan}\beta `$, $`\tau `$ identification takes on greater importance and good sensitivity to $`b`$ jets becomes even more crucial.
An MSSM charged Higgs with mass less than $`m_tm_b`$ is expected to be produced at the Tevatron through the decay of top quarks – a very different mechanism than for neutral Higgs particles. In regions of large and small $`\mathrm{tan}\beta `$ (away from $`\mathrm{tan}\beta `$$``$6) the branching ratio for $`t`$ $``$$`H^\pm `$$`b`$ is predicted to be quite large. The decay of the $`H^\pm `$ is mainly to W$`b`$$`\overline{b}`$ and $`c`$$`s`$ for low $`\mathrm{tan}\beta `$ and to $`\tau `$$`\nu _\tau `$ for high $`\mathrm{tan}\beta `$. These final states are sufficiently different from the SM top-quark decays that standard top analyses would have low efficiency for them. This means that, in addition to looking explicitly for the $`H^\pm `$ decay products (especially $`\tau \nu `$) a “disappearance” search is also possible for $`H^\pm `$. For a substantial branching ratio of top to charged Higgs, the measured p$`\overline{\text{p}}`$$``$$`t`$$`\overline{t}`$$`X`$ cross-section would be smaller than the SM expectation. A discrepancy between measurement and prediction can provide evidence for $`H^\pm `$.
## 3 Higgs Searches at Run I
Searches for Higgs have been performed by CDF and in all the main SM final states as well as in models beyond the SM. Results are summarized in Table 2. The reach of CDF and in production cross-section multiplied by $`h_{SM}`$$``$$`b`$$`\overline{b}`$ branching ratio in standard Higgs modes is far weaker than the predictions of the SM. For certain models beyond the SM, however, sizable regions of parameter space can be excluded. The charged Higgs of the MSSM is an especially interesting search since the Tevatron is the only facility that can take advantage of the $`t`$$``$$`H^\pm `$$`b`$ mode.
### 3.1 The Four-$`b`$ Final State at CDF
A preliminary analysis from CDF is another good example of the possibilities of Run I searches in beyond-the-SM scenarios. Here, $`b`$-quark identification of the CDF detector is used to select events with four $`b`$ jets in the final state. This topology is expected at large $`\mathrm{tan}\beta `$ in several SUSY models where the coupling of some of the neutral Higgs scalars to $`b`$-quarks is enhanced. The final state arises in $`b`$$`\overline{b}`$ production when one of the primary $`b`$-quarks radiates a neutral Higgs which, then decays to $`b`$$`\overline{b}`$. The cross section for this process goes as $`\mathrm{tan}^2\beta `$, and can therefore become sizable at large $`\mathrm{tan}\beta `$. The cross section is also affected by details of the mixing between left- and right-handed stop squarks. Results are presented for two extreme cases: that of minimal and maximal stop mixing.
This CDF analysis, which is an update of that presented in the Run II Higgs Working Group Report, uses a multijet trigger requiring a total cluster energy of $`>`$125 GeV and at least four trigger clusters with energy $`>`$15 GeV for an integrated luminosity of 91 pb<sup>-1</sup>. Offline, at least four jets are required, with $`E_T`$$``$15 GeV and $`\eta `$$`<`$2.4<sup>1</sup><sup>1</sup>1The pseudo-rapidity, $`\eta `$, is defined as $`\eta =\mathrm{ln}(\mathrm{tan}\frac{\theta }{2})`$, where $`\theta `$ is the polar angle.. Three of these jets must be tagged by the CDF secondary vertex algorithm as arising from $`b`$-quarks. In order to further reject the large QCD multi-jet background, additional criteria are imposed to take advantage of the distinct topology of four $`b`$ jets arising from $`\varphi `$$`b`$$`\overline{b}`$ ($`\varphi `$ = $`h`$,$`H`$,$`A`$) production. These criteria involve an $`M_\varphi `$ dependent cut on the $`E_T`$ of the three highest-$`E_T`$ jets, a requirement that the azimuthal angle between the two leading $`b`$ jets be larger than 1.9 radians, and an $`M_\varphi `$ dependent cut on the invariant mass of several different jet pairings.
Distributions in several variables used in the analysis, comparing data and SM prediction, are shown in Fig. 3 for CDF data with three $`b`$-tagged jets, prior to the imposition of the mass requirements. A breakdown of the efficiencies and backgrounds in the analysis, along with the number of observed events in the final selection is given in Table 3. No evidence for a $`\varphi `$$`b`$$`\overline{b}`$ signal is seen, providing a 95% C.L. limit on $`\sigma \times B`$ for this process of 25.7 pb. This limit can be interpreted within the framework of the MSSM as an exclusion region in $`\mathrm{tan}\beta `$ vs. $`M_h`$ or $`M_A`$ space. The excluded regions are presented in Fig. 4. As can be seen, the excluded regions extend significantly those previously probed by LEP.
## 4 Higgs Searches at Run II
Despite the good sensitivity of Run I analyses to several non-SM Higgs bosons, the results for the SM Higgs are (as expected) not significant. As we saw in Table 2, $`\sigma \times B`$ limits for final states expected from the SM Higgs are, at best, a factor of 50 higher than predictions. Given this picture, why are we optimistic about Higgs searches in Run II? There are three reasons: higher $`\sqrt{s}`$, higher luminosity and better detectors.
The increase in center of mass energy of the Tevatron for Run II is modest (1.8 to 2.0 TeV). This translates, however, into a substantial gain in cross section for associated Higgs production of approximately 20% in the SM.
As mentioned previously, the integrated luminosity collected in Run II is expected to be at least 2 fb<sup>-1</sup> by 2002 and should reach $``$15 fb<sup>-1</sup> before the start of the LHC. This large data set will be our main lever on the SM Higgs. However, it should not be forgotten that if we see evidence for a Higgs particle, larger control samples such as Z$``$$`b`$$`\overline{b}`$ that will be available in Run II will give us confidence that what we observe actually corresponds to signal.
In order to take maximum advantage of the glorious new data sets, both CDF and are being upgraded. From experience with Run I analyses and some theoretical guidance, a clear picture has emerged of the most important detector properties required for Higgs searches. Missing energy, leptons and $`b`$-quarks are the experimental pillars of Higgs searches at hadron machines. Efficient identification and accurate reconstruction of these objects requires all features of the detectors to work at full capacity, and consequently all aspects are being overhauled. Some of the improvements that have the most impact on Higgs searches are mentioned below.
Improving lepton identification is mainly a question of increasing the coverage of the calorimeters for electrons (which also determines missing energy resolution) and the muon chambers for muons. The Run I calorimeters of both experiments were excellent. Therefore, no changes are being made, aside from those required to adapt to the new beam conditions. CDF and are, however, both increasing the effective coverage of their muon systems.
Identification and reconstruction of $`b`$-quarks depends critically on tracking (although soft lepton identification also plays a role). To improve prospects, the old CDF silicon detector is being replaced with new 3D readout detectors that provide stand-alone silicon tracking to $`\eta `$$`<`$2.0. This amounts to a 40% increase in acceptance. Using this detector, CDF expects to gain in the efficiency of double $`b`$-tagging for $`t`$$`\overline{t}`$ events by a factor of 3.5. The tracking system of will be even more radically revamped. Central to this is the addition of a magnet providing a solenoidal field of 2.0 T in the tracking volume. The tracking system will consist of cylinders of scintillating fibers and a silicon detector with 3D readout extending to $`\eta `$$`<`$1.7. This will allow to join in the $`b`$-quark game at the same level as CDF. Both detectors are also adding dedicated trigger systems to identify $`b`$-quarks online.
### 4.1 Prophecies for Run II
We now embark into the realm of speculation about what will happen in Run II. This is not purely an exercise in fantasy, because it is extremely important to understand how detector limitations will affect Higgs searches while these parameters can still be modified. It is also crucial to know what luminosity is required to achieve sensitivity to the Higgs at different masses, as this will strongly influence the running strategy. As such, a Fermilab-wide working group, consisting of representatives from CDF, and the Theory Group, was established to study Higgs issues at Run II. The results presented in this section are based mainly on a preliminary version of the Working Group report (version 3) available at the time of the conference. The most up to date version (version 6) can be found on the Working Group’s web page.
Of course, to make predictions that have any chance of correctly fortelling the future we need accurate simulations of key performance parameters. Unfortunately, full simulations of the upgraded CDF and detectors are still evolving. However, even with a relatively simple simulation, using parameterized detector response, we can go a long way towards answering questions that are relevant for the construction and early running phases of Run II concerning detector resolution and efficiencies required for Higgs sensitivity and luminosity limitations on the mass reach.
The simulation used for the bulk of the results presented here, referred to as SHW, parameterizes important detector resolutions and efficiecies using an “average” of the foreseen Run II CDF and detectors. Most of the detector parameters in SHW are tunable, allowing studies to be made of how a specific parameter impacts Higgs sensitivity. As a baseline, most analyses use a track reconstruction efficiency of 97% for tracks with $`\eta `$$`<`$2 and $`P_T`$$`>`$300 MeV, relative energy resolutions for the electromagnetic and hadronic calorimeters of 20%/$`\sqrt{E}`$ and 80%/$`\sqrt{E}`$ respectively, $`b`$-tagging efficiency of $``$60% for $`E_T`$=100 GeV and a $`b`$$`\overline{b}`$ relative mass resolution of 10–14%.
A few warnings about SHW are in order. First, since SHW is a parameterized simulation, details of event-by-event detector response are missing. This means that systematic effects and hardware-related background and misidentifications are largely neglected. Their impact can be estimated, however, from extrapolations of Run I results, and thus are not completely ignored. Another difficult issue concerns the trigger. Excellent trigger performance will be crucial to obtaining good results. The most questionable of these, the hadronic event triggers, are considered in the analyses outlined here, however, leptonic triggers are generally taken to be 100% efficient. This is a reasonable assumption if lepton triggers function as foreseen. In general, SHW predictions should be taken in their context – a means to understand the detector and accelerator requirements so as to achieve competitive sensitivity to Higgs. While more elaborate simulations may yield slightly more accurate predictions in some areas, only data will tell us the real story.
Before turning to channel-by-channel sensitivities as a function of Higgs mass and luminosity it is worthwhile to describe the SHW results concerning the key Higgs-search detector parameters: missing energy resolution, lepton identification, $`b`$-quark identification and $`b`$$`\overline{b}`$ mass resolution.
1. Missing $`E_T`$ resolution was excellent in Run I and no gains are foreseen in this area.
2. Lepton identification efficiency in CDF and is mainly governed by geometrical acceptance. Improvements are being made in the muon systems of both detectors.
3. Tagging of $`b`$-quarks plays an important role in any Higgs search. However, signal significance ($`S/\sqrt{B}`$) grows at a faster rate if $`b`$$`\overline{b}`$ mass resolution is improved than if $`b`$-tagging efficiency is increased. This highlights the importance of a good understanding of the $`b`$-jet energy scale.
Sensitivity to an SM Higgs boson from a combination of CDF and expectations, as measured by signal over the square-root of background ($`S/\sqrt{B}`$), for various decay channels, as a function of Higgs mass, is presented in Fig. 6 for an integrated luminosity of 1 fb<sup>-1</sup> per experiment. Several points are apparent. First, the mass reach of the Tevatron experiments is significantly improved by considering final states produced when the Higgs (at high mass) decays to real or virtual boson pairs. Second, good improvements in sensitivity over purely cut-based analyses can be expected when using multi-variate techniques such as neural net analyses. Finally, it is clear that 1 fb<sup>-1</sup> per experiment will not get us to the Higgs. This is quantified in Fig. 6 where the combined CDF and 95% C.L. limits, 3$`\sigma `$ evidence and 5$`\sigma `$ discovery thresholds for a given integrated luminosity delivered to each experiment are plotted as a function of Higgs mass. With the minimal Run II integrated luminosity of 2 fb<sup>-1</sup>, the Tevatron will barely, if at all, extend the expected LEP2 Higgs mass limit of $``$115 GeV. With more than 10 fb<sup>-1</sup> per experiment, an SM Higgs could be excluded up to around 180 GeV. First hints of a real Higgs signal (at 3$`\sigma `$) would only appear beyond what has already been excluded for an integrated luminosity of at least 20 fb<sup>-1</sup> per experiment.
An important consideration in making the projections in Fig. 6 is our confidence in the predictions for background. Estimates of background levels based purely on Monte Carlo are notoriously unreliable; especially those originating from tails of distributions. An example is the QCD $`b`$$`\overline{b}`$ background in the $`\nu \overline{\nu }`$$`b`$$`\overline{b}`$ channel. To take account of this unreliability, a relative uncertainty on the background ($`B`$) in each channel of a minimum of 10% or 1/$`\sqrt{LB}`$ (where $`L`$ is the integrated luminosity) is used in combining the individual channels to produce Fig. 6. Of course, the increased luminosity of Run II should provide better understanding of the background based on control data samples, which can be used tighten selection criteria, thereby reducing background systematics.
Projected sensitivities for both neutral and charged Higges in the framework of the MSSM indicate that a relatively high integrated luminosity (10–15 fb<sup>-1</sup> for exclusion or 20–30 fb<sup>-1</sup> for discovery) will also be needed to reach decisive conclusions. However, if this is delivered, SUSY could be discovered or constrained over significant regions of the MSSM parameter space.
## 5 Conclusions
As we have seen, the Tevatron has been a very active field for Higgs searches. Several interesting limits have come out of Run I analyses relevant to predictions beyond the SM, but results for the minimal Higgs have not dented the SM. Nevertheless, the valuable experience gained in Run I, is already being applied to the upcoming Run II, slated to start in March of 2001. In order to have the best possible detectors for Higgs searches and to help set optimal parameters for the next run, studies have been initiated by the Fermilab Run II Higgs Working Group to determine the effect of detector choices and luminosity on the Higgs reach at the Tevatron. The main improvement in Run II that makes us optimistic about Higgs prospects is certainly the increased luminosity, a factor of 20 or perhaps as much as 300 over that delivered in Run I. Detector improvements will also play a big role. The main gains here come in $`b`$ identification and $`b`$$`\overline{b}`$ mass resolution – both of which are essential for discovering the Higgs. Not to be overlooked is the fact that will fully enter the arena of $`b`$-tagging with their upgraded tracking system in Run II. This will have a major impact on the overall Tevatron sensitivity. The big lesson that Run II Higgs prophecies teach us though is that luminosity will be crucial. With only the initial 2 fb<sup>-1</sup> the Tevatron will not extend the eventual LEP2 Higgs sensitivity, of $``$115 GeV. With 10 fb<sup>-1</sup> we could exclude at 95% C.L. an SM Higgs up to $``$180 GeV and with 20 fb<sup>-1</sup> we could see evidence at the 3$`\sigma `$ level for Higgs masses up to 180 GeV. Similarly, strong sensitivity to a wide region of MSSM parameter space can be made with more than 10 fb<sup>-1</sup>. These sensitivities are especially interesting given that the lightest Higgs is predicted to lie below $``$130 GeV in the MSSM. Needless to say, anticipation at the Tevatron is high!
## 6 Acknowledgements
This talk would have been entirely free of content without the help of a large number of people. I would especially like to acknowledge the wise advice of Max Chertok, Regina Demina, Mark Kruse, André Turcot, Juan Valls and Weiming Yao. Finally, huge thanks go to the conference organizers for providing us with such a stimulating meeting and such a good snow fall.
## 7 References
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# Formation of Structure in Snowfields: Penitentes, Suncups, and Dirt Cones
## I Light Reflection on Clean Snow
The model for penitente growth we derive here contains simplifying assumptions; we hope to capture the essential features while neglecting some effects We will discuss the assumptions and their limits of validity as the model is developed. Some of the most important simplifications include considering the latent heat to be constant and including only first-order, isotropic reflections. We focus on a one-dimensional model of penitentes, assuming invariance in the transverse direction, although it is straightforward to generalize these equations to two dimensions or to multiple reflections.
We consider the height of the snow surface $`h(x,t)`$, and seek an equation for the time evolution of $`h`$.
### A Snow Ablation
Heat incident on the surface leads to ablation—the height $`h`$ decreases as the snow melts or sublimates. We assume that ablated snow vanishes into the air or drains, and therefore that the flow of water along the surface and re-freezing are not important (and similarly that other changes in the nature of the snow are unimportant). This model can apply to either melting or sublimation. We use the term “ablation” to refer to removal of snow in either way.
Suppose a point on the surface absorbs a power per unit horizontal area $`P(x)`$. The latent heat required to ablate a unit volume of snow is $`L`$. Combining this with an effective diffusive smoothing term (see below) gives the evolution equation for the surface:
$$\frac{h}{t}=\frac{P(x)}{L}+D\frac{^2h}{x^2}.$$
(1)
For clean snow, we assume that $`L`$ is a constant (independent of $`x`$). This is true when the surface temperature and humidity are approximately constant. As discussed in the introduction, fully developed penitentes may have melting in the hollows and sublimation in the tips—a situation which requires $`L`$ to vary along the surface. Indeed, the variation in $`L`$ might be the essential effect for large structures. For small angle structures, that is, amplitude small relative to wavelength, $`L=`$ constant should be a good approximation. Later, we will include spatial variation in the effective $`L`$ due to dirt on the snow surface—see section II.
The second term in equation 1 for the surface height is a simple form of the small-scale cutoff: a diffusive term with diffusion constant $`D`$. As we will see below, in the absence of any smoothing term, the model can produce arbitrarily small structures. This is clearly not realistic, because the physics at small scales will cut off the instability. For the qualitative results here, the exact mechanism of the small-scale cutoff is not essential; the main point is that there is some minimum size structure which can form. A natural small-scale cutoff is the extinction length of sunlight, which defines the thickness of the snow layer in which the light scattering takes place. Points on the snow surface within one extinction length are not optically independent, and therefore such nearby points ablate at the same rate. The extinction length depends on the density and grain structure of the snow. The typical extinction length for old snow (grain radius 1 mm) is of order 1 cm. We will choose the diffusion coefficient so that the characteristic cutoff length is of order the optical extinction length. Again, remember that this term in the height equation is a simplified representation of the small-scale physics, and any conclusions which depend sensitively on the form of this term should be considered suspect. Note that diffusion of heat through the snow might seem another natural form of the small-scale cutoff; however the gradients of temperature in the snow are not large enough for thermal diffusion to stabilize short wavelengths.
We note here that recent work by Nodwell and Tiedje considers the scattering of light in the snowpack in quantitative detail. They find a range of length scales where suncups can form (both a minimum and maximum wavelength), a result which our simplified model cannot produce.
### B Light Reflection
In this section we describe the reflection of sunlight from the snow surface. We assume that the sunlight shines directly down (in the $`z`$ direction) and has a uniform power per unit length $`I`$. The parameter characterizing reflections is the albedo $`\alpha `$, which denotes the fraction of light reflected. Thus the absorbed power per unit length is $`(1\alpha )I`$. For old snow—called firn—a typical value is $`\alpha =0.5`$ .
The reflecting properties of snow are different from those of a mirror. Snow looks white because it scatters light in many directions, as we would expect for a rough surface. Here we treat the light using ray optics, and assume the surface reflects isotropically, thus the power is distributed uniformly into $`\pi `$ of solid angle outside the surface. We approximate that the reflection occurs at the surface of the snow. (As mentioned above, the reflection takes place in a layer of order 1 cm thick. We ignore this in formulating the reflections, and include its effects schematically through the diffusive term.)
Using these properties, the total amount of light scattered from an interval around point $`x_1`$ to the interval between $`x`$ and $`x+dx`$ is
$$\frac{\alpha I}{\pi }d\theta dx_1,$$
(2)
where $`d\theta `$ is the angle subtended by the surface between $`x`$ and $`x+dx`$ (see Figure 5).
We can find $`d\theta `$ in terms of the shape of the surface.
$`d\theta ={\displaystyle \frac{dl}{p}}={\displaystyle \frac{|𝐩\times d𝐬|}{p}}={\displaystyle \frac{\mathrm{\Delta }h\mathrm{\Delta }xh^{}(x)}{\mathrm{\Delta }h^2+\mathrm{\Delta }x^2}}.`$where we have used $`h^{}=h/x`$ and
$`\mathrm{\Delta }x`$ $`=`$ $`x_1x`$ (3)
$`\mathrm{\Delta }h`$ $`=`$ $`h(x_1)h(x)`$ (4)
and $`d𝐬`$ is the vector tangent to the surface
$$d𝐬=dx(1,h^{})$$
(5)
We define the vector $`𝐩`$, which points from the point $`x_1`$ to the point $`x`$. From Figure 5, we can see that
$$p=\sqrt{\mathrm{\Delta }x^2+\mathrm{\Delta }h^2}$$
(6)
To find the total power reflected to point $`x`$, we must add up the intensity scattered from all points $`x_1`$:
$$P_r(x)=\frac{\alpha I}{\pi }\frac{dx_1(\mathrm{\Delta }hh^{}(x)\mathrm{\Delta }x)}{\mathrm{\Delta }x^2+\mathrm{\Delta }h^2}$$
(7)
The integrand in this equation is the propagator for light intensity, it describes how the intensity is carried from one point to another on the surface; $`P_r(x)`$ is the intensity due to a single reflection. To include multiple reflections, we can write the power as an integral equation for $`P`$.
$$P(x)=(1\alpha )I+\frac{\alpha }{\pi }\frac{dx_1P(x_1)(\mathrm{\Delta }hh^{}(x)\mathrm{\Delta }x)}{\mathrm{\Delta }x^2+\mathrm{\Delta }h^2}$$
(8)
This can be written as a power series in $`\alpha `$. We will only consider single reflections here, which does not introduce a large error when $`\alpha `$ is small. For old snow, a typical value of $`\alpha 0.5`$. Including the higher-order correction from multiple reflections may be important in determining the precise details of the largest shapes.
This formula for reflected intensity is not complete, because it neglects the line-of-sight constraint. Light cannot scatter from $`x_1`$ to $`x`$ if the path of the light ray is blocked by another part of the surface. This requirement is a nonlinear constraint which is difficult to handle analytically but is straighforward to implement in numerics. We typically indicate the constraint schematically, by writing “line of sight” under the integral:
$`P_r(x)={\displaystyle \frac{\alpha I}{\pi }}{\displaystyle _{\begin{array}{c}\text{ line of}\\ \text{ sight}\end{array}}}{\displaystyle \frac{dx_1(\mathrm{\Delta }hh^{}(x)\mathrm{\Delta }x)}{\mathrm{\Delta }x^2+\mathrm{\Delta }h^2}}`$We can also write a necessary (but not sufficient) criterion for the line of sight constraint, when applied to local analysis within one “basin.” The two points $`x`$ and $`x_1`$ are within line of sight of each other when the dot product of the vector normal to the surface and the vector $`𝐩`$ is less than 0: $`𝐧𝐩=\mathrm{\Delta }h\mathrm{\Delta }xh^{}(x)>0`$. (See Figure 5.) Note, however, that this simple criterion will miss intermediate bumps in the surface. In other words, the constraint may be satisfied but no reflection occurs between $`x_1`$ and $`x`$ because the line of sight is blocked by an intervening peak.
### C Model
The equations combining reflection and ablation are
$$\frac{h}{t}=\frac{\alpha I}{L}(x)+Dh^{\prime \prime }$$
(9)
where we have defined the integral
$$(x)=\frac{1}{\pi }_{\begin{array}{c}\text{ line of}\\ \text{ sight}\end{array}}\frac{dx_1(\mathrm{\Delta }hh^{}(x)\mathrm{\Delta }x)}{\mathrm{\Delta }x^2+\mathrm{\Delta }h^2}.$$
(10)
The intensity of the sun determines a characteristic ablation rate $`IL^1`$, where $`L`$ is the latent heat per unit volume. Combining this velocity with the diffusion coefficient $`D`$ gives a length
$$\overline{\mathrm{}}=\frac{DL}{I}$$
(11)
and time
$$\overline{t}=\frac{DL^2}{I^2}.$$
(12)
The solar constant gives the intensity of solar radiation at the top of the atmosphere $`I=1.4\times 10^6\text{erg cm}^2\text{sec}^1`$; we therefore choose $`I=10^6\text{erg cm}^2\text{sec}^1`$ as the typical value of $`I`$ under bright sunny conditions. The latent heat depends on density. Freshly fallen snow has a density of between 0.05 and 0.2 $`\text{g cm}^3`$, while older snow that has survived one melt season has a density range of 0.4 to 0.8 $`\text{g cm}^3`$ . Here we pick an intermediate density of 0.3 $`\text{g cm}^3`$ for our estimates. This gives a latent heat per unit volume for melting $`L=10^9\text{erg cm}^3`$ and a melting rate $`I/L=10^3\text{cm sec}^1`$. We pick $`D=2.5\times 10^5\text{cm}^2\text{sec}^1`$, where this choice is made so that the most unstable wavelength is 2 cm (see below). In this case the length scale $`\overline{\mathrm{}}=0.25`$ mm. and the time scale $`\overline{t}=25`$ seconds.
For sublimating snow, the latent heat is seven times larger. this gives the slower melting rate $`I/L=1.4\times 10^3\text{cm sec}^1`$, larger length scale $`\overline{\mathrm{}}=1.75`$ mm and time scale $`\overline{t}=1225`$ seconds.
We will now perform a perturbation analysis of Equations 9 to see how the size structures formed compares to the scale $`\overline{\mathrm{}}`$. We have set up the problem so that structures will initially form on a scale roughly comparable to $`\overline{\mathrm{}}`$, and expect the perturbation analysis to give this result.
### D Quasi-Linear Regime
Here we show how an approximate linear analysis of the equations can be performed. This allows us to derive the dispersion relation, which characterizes when the system is stable or unstable. There is a fastest growing mode determined by the competition between reflection and diffusion. The length scale of this mode is related to the basic scale $`\overline{\mathrm{}}`$ from dimensional analysis above; we determine the prefactor here. The results are significant because they describe how the physical parameters affect the instability. We will argue that reflection favors structures on scales as small as possible. On the other hand, the small-scale cutoff limits the smallest structures possible. Therefore we expect the fastest-growing mode to be of order the cutoff size.
The reflection integral is scale invariant: upon rescaling $`x`$ and $`h`$ by the same amount the integral $`(x)`$ is unchanged. Thus in the absence of diffusion, there is no characteristic scale in the problem. Therefore a shape with aspect ratio 1—a shape with variations in $`h`$ comparable to variations in $`x`$—should have a growth rate of order 1 (in the absence of boundary effects). The integral contributes a shape factor independent of the amplitude of the shape $`\delta `$. Therefore the rate of change of amplitude $`\dot{\delta }`$ is constant.
To examine shapes with aspect ratio far from one, we start with an aspect-ratio 1 shape, then transform $`x\lambda x`$ and $`h\delta h`$ . When $`\delta \lambda `$, we find that the integral scales with the basic angle $`\delta /\lambda `$: $`\delta /\lambda `$. Thus for small perturbations, we expect a growth rate proportional to the amplitude ($`\dot{\delta }\delta `$).
For sufficiently small $`\delta /\lambda `$, we treat the contribution from the reflection integral as a numerical factor of order 1. Note that a sinusoidal perturbation is not an eigenshape for small amplitudes; we do not know what the actual eigenshapes are. The dominant contribution is the scaling with $`\delta `$, and we neglect the other (slower) dependence on position, amplitude, etc. Thus the quasilinear equation for a small-amplitude variation in the surface $`h=\delta \mathrm{sin}qxe^{\omega t}`$ is approximately
$$(x)\frac{q}{\pi }\delta \mathrm{sin}qxe^{\omega t}$$
(13)
which gives a dispersion relation
$$\omega =\frac{\alpha I}{\pi L}qDq^2$$
(14)
This argument selects a fastest-growing mode with wavenumber
$`q_{}`$ $`=`$ $`{\displaystyle \frac{\alpha I}{2\pi LD}}`$ (15)
$`\omega _{}`$ $`=`$ $`{\displaystyle \frac{(\alpha I)^2}{4\pi ^2L^2D}}=q_{}^2D`$ (16)
These equations are the dimensional analysis result, with an estimate of the prefactor from the scaling argument. Plugging in values of typical parameters given above, we find the most unstable wavelength for melting $`\lambda _{}=2\pi /q_{}`$ of 2 cm, and characteristic time 4000 sec. In the case of sublimation the wavelength is 14 cm and the time $`2\times 10^5`$ sec. The choice of the diffusion coefficient is now clear: we chose $`D`$ to give a most unstable wavelength of 2 cm. We have put in diffusion as a simplified representation of the small-scale physics, and chosen its value so that the numbers make sense. It is important to remember that because of this choice of $`D`$, the numbers calculated here cannot be considered a prediction of the initial size structures that form. The calculation of real interest is how this instability is changed by dirt, as discussed in the following section.
Although it agrees well with simulations of initial growh of perturbations which compute the reflected intensity at each point , we must remember that this analysis is only quasi-linear because we do not know the eigenfunctions of the reflection integral, and superposition does not hold: because the integral is nonlocal, a surface variation with two modes of different wavelength cannot be described by the addition of two modes with different $`q`$.
## II Effects of Dirt
A layer of dirt on the surface of the snow changes its properties. We model both the optical and insulating effects of dirt, and fit the theory to melting data measured by Driedger . These data allow measurement of a crucial parameter in the model, and the good agreement between theory and experiment show that we have captured the important effects of dirt. The essential features are that thin dirt speeds ablation, because it increases absorption, while thick dirt insulates the snow, slowing ablation. This basic behavior leads to the two different regimes of instability .
Dirt looks black because it absorbs light. The presence of dirt effectively decreases the surface albedo and therefore increases the fraction of absorbed light. We assume light has a probability of being absorbed that is constant per unit thickness of dirt. The fraction of light not absorbed by the dirt is $`e^{s/s_e}`$, where $`s`$ is the dirt thickness and $`s_e`$ the extinction length in the dirt—typically of order the characteristic dirt particle size. Therefore dirt modifies the albedo according to
$$\alpha _d=\alpha e^{s/s_e}.$$
(17)
Note that absorption by the dirt layer is not isotropic—more light will be absorbed near grazing incidence, decreasing the reflection even more. The qualitative effect of dirt remains the same however, and thus we neglect this anisotropy. Increased absorption through a lower effective albedo hastens snow ablation.
But the dirt also slows ablation. In the presence of an insulating dirt layer, the temperature at the surface of the snow is decreased below the ambient temperature, and more heat is required to ablate a given amount of snow. Suppose an amount of heat $`L`$ is necessary to ablate a unit area of clean snow. How much additional heat is required in the presence of a dirt layer? At steady state the temperature satisfies
$$^2T=0.$$
(18)
When the radius of curvature of the surface is large compared to the dirt thickness (the important limit for growth of perturbations) we can treat the snow surface as planar, leading to variations in $`T`$ in the $`z`$ direction only. The boundary conditions are: At the dirt-air interface ($`z=0`$), the temperature must be equal to the ambient temperature. The temperature gradient at the surface due to heat flux into the dirt from the air is $`T^{}(z=0)=P/\kappa `$, where $`P`$ is the incident power flux and $`\kappa `$ the thermal conductivity of the dirt. Thus we find that the temperature at the snow surface is less than $`T(z=0)`$ by an amount $`\mathrm{\Delta }T=Ps/\kappa `$. An extra amount of heat $`\mathrm{\Delta }Q=C\mathrm{\Delta }T`$ is needed to raise the snow temperature up to its value in the absence of dirt, where $`C`$ is the heat capacity of the snow. Thus the effective latent heat for a dirt thickness $`s`$ is
$$L_d=L+\frac{CPs}{\kappa }$$
(19)
Both $`L`$ and $`C`$ depend on the ambient temperature $`T`$. However, the dependence is sufficiently weak that we can neglect it.
Combining these two effects we find that the snow ablation velocity for a flat surface covered with dirt is
$$m(s)=\frac{I}{L}g(s)$$
(20)
where $`g`$ is a dimensionless function of the dirt thickness. In this model,
$$g(s)=\frac{1\alpha e^{s/s_e}}{1+\gamma s(1\alpha e^{s/s_e})}$$
(21)
where we have defined the dimensionless measure of the insulating value of dirt:
$$\gamma =\frac{s_eCI}{L\kappa }.$$
(22)
The non-monotonic behavior of this curve—positive slope for small $`s`$ and negative for large $`s`$—is the important qualitative result. Note that in the absence of dirt the ablation rate is as expected:
$$m(s=0)=\frac{I}{L}(1\alpha )$$
(23)
A fit to the data of Driedger is shown in Figure 6. Driedger measured melting rates of a flat surface for different dirt thickness. The plot shows $`m/m(s=0)`$ versus dimensionless dirt thickness. We picked $`s_e=1`$ mm from Driedger’s measurement of the dirt particle size. Fitting the data to Equation 21 allows us to determine the dimensionless insulation coefficient $`\gamma `$. For $`\alpha 0.5`$ (from other measurements ) the fit gives $`\gamma =0.047`$. We can also estimate the parameter $`\gamma `$ using other data (see below). The estimate gives $`\gamma `$ close to the value obtained from this fit.
This model and the experiment of Driedger are in the regime where solar radiation is the dominant heat source. The discussion of Rhodes et. al. points out that the ablation curve changes when sensible heating is important. In fact, if radiation is negligible the curve will monotonically decrease as the dirt thickness increases, because light absorption effects disappear in this limit. It is straightforward to adjust the model to include other sources of heating. Measurements of the type Driedger performed, compared to the type of model presented here, could in principle give information on the relative importance of radiant and sensible heating.
### A Dynamics of Dirt
As the snow surface ablates, the dirt layer on it moves (Figure 4). We assume the particles are sufficiently small that the snow moves purely normal to the surface. The sideways ($`x`$-direction) velocity of a piece of dirt is
$$v=\dot{h}h^{}$$
(24)
The thickness of the dirt $`s(x)`$ must obey a conservation equation $`\dot{s}+(vs)=0`$, since we assume dirt is neither deposited on nor removed from the surface. The evolution equation for the thickness of dirt is thus
$$\frac{s}{t}=\frac{}{x}(vs)=(\dot{h}h^{}s)^{}$$
(25)
When the surface of the snow is flat ($`h^{}=0`$) the velocity of the dirt $`v=0`$. Thus the tops of peaks and the bottoms of valleys are equilibrium points. The peaks are stable equilibria where dirt becomes concentrated, while valleys are unstable (Figure 4).
### B Model
We now rewrite the model equations incorporating dirt. We have equations for the height of the surface $`h`$, the dirt thickness $`s`$, and the incident power $`P`$.
$`\dot{h}`$ $`=`$ $`{\displaystyle \frac{P(x)}{L}}{\displaystyle \frac{1}{1+\frac{C}{\kappa L}Ps}}+Dh^{\prime \prime }`$ (26)
$`\dot{s}`$ $`=`$ $`\left(\dot{h}h^{}s\right)^{}`$ (27)
The only sources of heat flux $`P`$ we will consider are direct and reflected radiation.
$$\frac{P(x)}{L}=(1\alpha e^{s/s_e})\frac{I}{L}+\frac{\alpha e^{s/s_e}I}{\pi L}_{\begin{array}{c}\text{ line of}\\ \text{ sight}\end{array}}\frac{dx_1(\mathrm{\Delta }hh^{}(x)\mathrm{\Delta }x)}{\mathrm{\Delta }x^2+\mathrm{\Delta }h^2}$$
(28)
We use the same reference ablation rate as in Section I: $`I/L=10^3\text{cm}\text{sec}^1`$. However, the presence of dirt introduces a new length scale in the problem: the length scale for light absorption by the dirt. We choose to nondimensionalize in terms of this length, since the physically important regimes of thin and thick dirt are measured relative to this thickness. When Driedger measured diameters of ash particles on a glacier, 90 percent of the particles had diameters between 0.25 and 1.0 mm. We therefore choose $`s_e=1`$ mm as the order of magnitude extinction length for dirt absorption; this choice is supported by the good fit to the data.
The dimensionless timescale comes from combining the ablation rate and length scale: $`\overline{t}_d=Ls_e/I=100`$ seconds. This is the time for a depth $`s_e`$ of snow to melt in bright sun. Glacial debris has $`\kappa 2\times 10^4`$ erg $`\text{cm}^1\text{sec}^1^{}`$K . This allows us to estimate the dimensionless parameter $`\gamma =s_eCI/(L\kappa )=0.03`$. Note that the thermal conductivity and the specific heat depend on the density, wetness, etc. The fit to Driedger’s data (Figure 6) gives a value of $`\gamma 0.047`$, somewhat larger than this estimate. We interpret this as a measurement of the dirt thermal conductivity $`\kappa `$, and therefore use the implied value $`\kappa =1.3\times 10^4`$ erg $`\text{cm}^1\text{sec}^1^{}`$K. The nondimensionalized diffusion constant is $`D\overline{t}_d/s_e^2=0.25`$.
For sublimation the time scale $`\overline{t}_d700`$ seconds and the dimensionless diffusion constant $`D\overline{t}_d/s_e^21.75`$; the dimensionless parameter $`\gamma `$ similarly decreases by a factor of 7.
The nondimensionalized equations are
$`\dot{h}`$ $`=`$ $`{\displaystyle \frac{P}{1+\gamma Ps}}+D^2h`$ (29)
$`\dot{s}`$ $`=`$ $`\left(\dot{h}h^{}s\right)^{}`$ (30)
$`P`$ $`=`$ $`r(1\alpha e^s)+{\displaystyle \frac{\alpha e^sr}{\pi }}{\displaystyle _{\begin{array}{c}\text{ line of}\\ \text{ sight}\end{array}}}{\displaystyle \frac{dx_1(\mathrm{\Delta }hh^{}(x)\mathrm{\Delta }x)}{\mathrm{\Delta }x^2+\mathrm{\Delta }h^2}}`$ (33)
The dimensionless control parameters are $`r`$, the solar light intensity; and $`s`$, the initial dirt thickness. Here we have introduced the parameter $`r`$:
$$r=\frac{I}{L}10^3\text{sec/cm}$$
(34)
to examine the effects of varying the light intensity away from the typical value.
### C Linear Analysis
Here we analyze the stability of equations (2933), including effects of dirt. There are two important regimes: when the initial dirt thickness is small compared to $`s_e`$, the dirt acts to modify the reflection-driven instability. We find that the instability is suppressed by the absorption of the dirt layer, exponentially in the dirt thickness. In this regime, dirt can also a induce travelling, dispersive instability of the snow surface. Qualitatively, this dispersion arises from the coupling of dirt motion to absorption. Dirt migrates to the highest point on the surface—but then the thicker dirt increases the ablation of that peak, and it ablates until it is no longer a local maximum. The existence of these waves is an experimentally testable prediction which has not, to my knowledge, been discussed before.
The other limit is when the dirt thickness is large compared to $`s_e`$. The effective albedo $`\alpha e^s0`$. Therefore the dirt instability is independent of light reflections; the “light” therefore acts simply as a source of heat. The instability is driven by dirt insulating the snow. The characteristic length and time scale of the instability depends only on the thermal properties of the dirt. Within this insulation-dominated regime, the behavior of the instability depends on whether $`s1/(\gamma r)`$ or $`s1/(\gamma r)`$—see below. Thus there are three different regimes of behavior, depending the dirt thickness.
As mentioned above, under different weather conditions uniform heating from the air may be more important than radiant heating. In this case any amount of dirt slows ablation of the snow , and the insulation-driven instability is the only one possible. This can be included in the model by removing the dirt-dependent absorption of light.
We will perform a linear perturbation analysis: we assume that variations of the dirt thickness $`\mathrm{\Delta }s`$ are always small. However, the initial uniform dirt thickness $`s_o`$ may be large or small relative to $`s_e`$; this initial thickness determines the limit of instability.
### D Thin Dirt Limit
Here we consider the limit $`s_os_e`$, meaning the initial uniform dirt thickness is small compared to the extinction length.
There are in general two modes of dirt modulation (Figure 7): the symmetric mode with constant thickness and the antisymmetric mode with $`\mathrm{\Delta }s=2ϵ\mathrm{cos}qx`$. The symmetric mode, because it has constant thickness, it simpler to analyze. Note that constant dirt thickness is unstable: any modulation in the dirt thickness tends to grow.
#### 1 Symmetric Mode
Because the symmetric mode has constant thickness, it insulates the snow surface uniformly. Therefore, no thick dirt instability can arise from the symmetric mode. But the symmetric mode affects the reflection-driven instability. We look for solutions of the form
$$h=mt+\delta e^{\omega t}\mathrm{cos}qx$$
(35)
where $`m(s)`$ is the ablation rate of a flat surface covered with dirt, calculated above. The dirt thickness $`s_o=`$ constant. We expand the equations to first order in $`\delta `$. The resulting dispersion relation is
$$\omega =\frac{\alpha re^{s_o}q}{\pi (1+(1\alpha )\gamma rs_o)^2}Dq^2.$$
(36)
Compare this to the clean snow dispersion relation, Equation 14. The first term (proportional to $`q`$) contains the factor $`e^{s_o}`$. This term decreases exponentially with increasing dirt thickness. For $`s_o`$ much larger than one, this term is so small that the instability practically does not exist. The factor $`(1+(1\alpha )\gamma rs_o)^2`$ in the dispersion relation results from uniform insulation by the dirt layer.
The most unstable mode $`q_{}`$ is
$`q_{}`$ $`=`$ $`{\displaystyle \frac{\alpha re^{s_o}}{2\pi D(1+(1\alpha )\gamma rs_o)^2}}`$ (37)
$`\omega _{}`$ $`=`$ $`q_{}^2D`$ (38)
Figure 8 shows how dirt cuts off the instability, with fixed light intensity $`r=1`$. When $`s_o1`$, the wavelength is close to the wavelength in the absence of dirt. However, the absorption of light by dirt becomes important for $`s_o>0.1`$ and the wavelength increases exponentially. As the wavelength increases, the growth rate of the instability decreases, and the instability becomes less readily observed.
#### 2 Antisymmetric Mode
The antisymmetric mode involves variations in the thickness of the dirt. We must solve for the coupling between snow ablation and dirt motion. The solution is of the form
$`h`$ $`=`$ $`mt+\delta e^{\omega t}\mathrm{cos}qx`$ (39)
$`s`$ $`=`$ $`s_o+2ϵe^{\omega t}\mathrm{cos}qx`$ (40)
where $`s`$ is the uniform dirt thickness at $`t=0`$. Upon linearization, Equation 30 for the motion of dirt relates the perturbation amplitudes
$$\frac{ϵ}{\delta }=\frac{ms_oq^2}{2\omega }.$$
(41)
The dispersion relation, to second order in $`q`$, is
$$\omega =[1\pm \sqrt{f}]\frac{\alpha re^{s_o}q}{2\pi (1+\gamma r^{}s_o)^2}[1\frac{1}{\sqrt{f}}]\frac{Dq^2}{2}$$
(42)
where $`f`$ is, defining $`w=1/(1+(1\alpha )\gamma rs_o)`$ and recalling $`m=(1\alpha e^{s_o})rw`$ is the dimensionless melting rate as a function of dirt thickness,
$$f=(\alpha re^{s_o}\pi ^1w^2)^2(1+\frac{4s_om(m^2\gamma \alpha re^{s_o}w^2)}{(\alpha re^{s_o}\pi ^1w^2)^2}$$
(43)
In the limit $`s_o0`$, this dispersion relation is identical to the symmetric mode. However, for increasing dirt thickness it contains effects from the dirt modulation. The term $`f`$ can be negative, leading to an oscillatory component to $`\omega `$. Thus dirt can cause the instability to travel on the snow surface, in a region of phase space shown in Figure 9. For the typical solar brightness $`r=1`$, any dirt thickness $`s_o>0.008`$ will induce travelling; therefore, most dirty snow surfaces should show this behavior. Qualitatively, this arises from the coupling of dirt motion to absorption. Dirt migrates to the highest point on the surface—but then the thicker dirt increases the ablation of that peak, and it ablates until it is no longer a local maximum. The positive and negative roots in the dispersion relation correspond to left and right moving modes. The existence of these travelling instabilities is an experimentally testable prediction.
Note that the equation is not well-behaved for $`f=0`$. When $`f=0`$ the terms in the equation coupling motion of dirt to ablation vanish; the dispersion relation reduces to the expression for the symmetric mode above.
When $`f`$ is negative, we can find the fastest growing wavelength by looking at the real part of $`\omega `$:
$`q_{}`$ $`=`$ $`{\displaystyle \frac{\alpha e^{s_o}r}{2\pi D(1+(1\alpha )\gamma rs_o)^2}}`$ (44)
$`\omega _{}`$ $`=`$ $`{\displaystyle \frac{D}{2}}q_{}^2`$ (45)
### E Thick Dirt Limit
The equations are considerably simplified in the limit of thick dirt $`s_o1`$. The effective albedo $`\alpha e^{s_o}0`$. Therefore the dirt instability is independent of any reflections; the quasi-linearized equations are truly linear in this limit. The thick-dirt instability is driven purely by dirt motion coupled to slower ablation under a thicker dirt layer. This instability is the linear precursor to the dirt cones of Figure 3.
Note that if light is not an important source of heat, the “thick dirt limit” is actually valid for all dirt thicknesses.
Replacing $`\alpha e^{s_o}0`$, the symmetric mode disappears. The background ablation rate $`m=r/(1+\gamma rs_o)`$. The dispersion relation is, to second order in $`q`$,
$$\omega =\pm \frac{\sqrt{\gamma s_om^3}}{2}q\frac{D}{2}q^2$$
(46)
Here no imaginary component to the dispersion relation is present; it is a straightforward linear instability with one growing mode. The most unstable wavenumber is
$$q_{}=\frac{1}{2D}\sqrt{\frac{\gamma r^3s_o}{(1+\gamma rs_o)^3}}$$
(47)
For a fixed value of the heat input $`r`$, the most unstable wavelength scales differently at small and large $`s_o`$:
$`\lambda _{}`$ $``$ $`s_o^{1/2}\text{for}s_o1/(\gamma r)`$ (48)
$``$ $`s_o\text{for}s_o1/(\gamma r)`$ (49)
The location of the minimum wavelength is determined by the dimensionless parameter $`\gamma `$, which represents how well the snow insulates per unit thickness. Therefore, even for optically thick dirt $`s_o1`$ there is a change in the behavior, depending on the value of $`s_o`$ compared to the insulation parameter. Since typically $`\gamma r=0.05`$, these limits are consistent.
There is an optimal $`s_o1/(\gamma r)2`$ cm where the wavelength is smallest. Figure 10 illustrates this: it shows the unstable wavelength vs. dirt thickness for the typical $`r=1`$, with the optimal $`s_o202`$ cm. Comparing this figure to the thin-dirt instability, we see that when $`s_o>1`$ the wavelength will initially decrease, then increase beyond $`s_o=20`$. The growth rate of this instability will be greatest where the wavelength is smallest.
## III Discussion: Comparison to Experiment
This paper has presented work on a simple theory to describe the initial formation ablation structures such as suncups, penitentes and dirt cones. We have tried to make the model as simple as possible while including the essential physics. As we have shown, most parameters in the equations can be calculated or measured in experiments, allowing predictions with no free parameters. The exception is the effective diffusion coefficient $`D`$, which we estimate using the value for light diffusion. However, we have not realistically treated the small-scale scattering of light in these schematic results.
At this point, the only quantitative comparison between this model and experiment is the prediction of ablation rate of a flat snow surface, compared with the data of Driedger in Figure 6. This measurement allows us to extract the dimensionless constant governing dirt insulation. The good agreement indicates we have captured the important effects of dirt.
The linear stability analysis of the equations shows the two types of instability described in the literature. The model predicts the dependence of the most unstable wavelength and characteristic growth rate on the experimental control parameters, predictions which could be tested. We argue that for little or no surface dirt, light reflection drives the instability. This instability is exponentially suppressed by a dirt layer, consistent with field observations. We predict travelling modes induced by a modulated dirt layer in this regime. The existence of such travelling modes is an experimentally testable new phenomenon.
In the presence of a thick layer of dirt, our analysis finds the insulation-driven instability, as expected. Here we predict an optimal dirt thickness where the instability is most easily observed, which depends on the thermal properties of the dirt.
The visually striking structures in the field are the larger structures: penitentes and dirt cones. Understanding the nonlinear regime of the model presented here is therefore of interest, and will be the subject of a future paper. The scale of both penitentes and dirt cones is typically larger than the size of smaller-amplitude structures. One way to explain this, which has been suggested from observations, is that large structures grow at the expense of small ones. Such coarsening behavior is also apparent in preliminary work on the nonlinear regime of the model presented here.
The most obvious problem with the results here is that we have considered variation of the surface height in only one direction. Checking whether the results are the same for a realistic 2D surface is a necessary extension of this work. A better understanding of the small-scale cutoff is also important. In particular, we need to understand how using different representations of the short-scale physics affect the numerical predictions (of the fastest-growing wavelength, for example).
Because the model here is simplified, we have left out some physical effects which may be important in the experiment. Our treatment of light reflection considered single reflections only, which may be a bad approximation with the albedo is close to 1 (large amount reflected). In the field, the sun of course is not always high overhead—the variation of the angle of incident light over the course of the day might change the shapes. Other possibly important effects which can occur in field situations include other sources of heat transfer to the surface, gravity, and the deposition/removal of dirt. Better comparison with lab or field experiments should indicate which of these effects are most important to include.
Acknowledgements: I am grateful to John Wettlaufer and Norbert Untersteiner for comments on this paper. I thank Eric Nodwell and Tom Tiedje for discussing their work on this problem with me. I also wish to thank Michael Brenner, Daniel Fisher, David Weitz, Martine Benamar, David Lubensky, and David Nelson for helpful discussions, questions, and criticism. This work was supported by the NSF under grant DMS9733030.
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# 1 𝐷⁺→𝜋⁻𝜋⁺𝜋⁺ Dalitz plot distribution. Since we have two identical particles this distribution was symmetrized.
Experimental evidence for a light and broad scalar resonance in $`D^+\pi ^{}\pi ^+\pi ^+`$ decay .
E. M. Aitala,<sup>9</sup> S. Amato,<sup>1</sup> J. C. Anjos,<sup>1</sup> J. A. Appel,<sup>5</sup> D. Ashery,<sup>14</sup> S. Banerjee,<sup>5</sup> I. Bediaga,<sup>1</sup> G. Blaylock,<sup>8</sup> S. B. Bracker,<sup>15</sup> P. R. Burchat,<sup>13</sup> R. A. Burnstein,<sup>6</sup> T. Carter,<sup>5</sup> H. S. Carvalho,<sup>1</sup> N. K. Copty,<sup>12</sup> L. M. Cremaldi,<sup>9</sup> C. Darling,<sup>18</sup> K. Denisenko,<sup>5</sup> S. Devmal,<sup>3</sup> A. Fernandez,<sup>11</sup> G. F. Fox,<sup>12</sup> P. Gagnon,<sup>2</sup> C. Gobel,<sup>1</sup> K. Gounder,<sup>9</sup> A. M. Halling,<sup>5</sup> G. Herrera,<sup>4</sup> G. Hurvits,<sup>14</sup> C. James,<sup>5</sup> P. A. Kasper,<sup>6</sup> S. Kwan,<sup>5</sup> D. C. Langs,<sup>12</sup> J. Leslie,<sup>2</sup> B. Lundberg,<sup>5</sup> J. Magnin,<sup>1</sup> A. Massafferri,<sup>1</sup> S. MayTal-Beck,<sup>14</sup> B. Meadows,<sup>3</sup> J. R. T. de Mello Neto,<sup>1</sup> D. Mihalcea,<sup>7</sup> R. H. Milburn,<sup>16</sup> J. M. de Miranda,<sup>1</sup> A. Napier,<sup>16</sup> A. Nguyen,<sup>7</sup> A. B. d’Oliveira,<sup>3,11</sup> K. O’Shaughnessy,<sup>2</sup> K. C. Peng,<sup>6</sup> L. P. Perera,<sup>3</sup> M. V. Purohit,<sup>12</sup> B. Quinn,<sup>9</sup> S. Radeztsky,<sup>17</sup> A. Rafatian,<sup>9</sup> N. W. Reay,<sup>7</sup> J. J. Reidy,<sup>9</sup> A. C. dos Reis,<sup>1</sup> H. A. Rubin,<sup>6</sup> D. A. Sanders,<sup>9</sup> A. K. S. Santha,<sup>3</sup> A. F. S. Santoro,<sup>1</sup> A. J. Schwartz,<sup>3</sup> M. Sheaff,<sup>17</sup> R. A. Sidwell,<sup>7</sup> A. J. Slaughter,<sup>18</sup> M. D. Sokoloff,<sup>3</sup> J. Solano,<sup>1</sup> N. R. Stanton,<sup>7</sup> R. J. Stefanski,<sup>5</sup> K. Stenson,<sup>17</sup> D. J. Summers,<sup>9</sup> S. Takach,<sup>18</sup> K. Thorne,<sup>5</sup> A. K. Tripathi,<sup>7</sup> S. Watanabe,<sup>17</sup> R. Weiss-Babai,<sup>14</sup> J. Wiener,<sup>10</sup> N. Witchey,<sup>7</sup> E. Wolin,<sup>18</sup> S. M. Yang,<sup>7</sup> D. Yi,<sup>9</sup> S. Yoshida,<sup>7</sup> R. Zaliznyak,<sup>13</sup> and C. Zhang<sup>7</sup>
(Fermilab E791 Collaboration)
<sup>1</sup> Centro Brasileiro de Pesquisas Físicas, Rio de Janeiro, Brazil, <sup>2</sup> University of California, Santa Cruz, California 95064, <sup>3</sup> University of Cincinnati, Cincinnati, Ohio 45221, <sup>4</sup> CINVESTAV, Mexico City, Mexico, <sup>5</sup> Fermilab, Batavia, Illinois 60510, <sup>6</sup> Illinois Institute of Technology, Chicago, Illinois 60616, <sup>7</sup> Kansas State University, Manhattan, Kansas 66506, <sup>8</sup> University of Massachusetts, Amherst, Massachusetts 01003, <sup>9</sup> University of Mississippi-Oxford, University, Mississippi 38677, <sup>10</sup> Princeton University, Princeton, New Jersey 08544, <sup>11</sup> Universidad Autonoma de Puebla, Puebla, Mexico, <sup>12</sup> University of South Carolina, Columbia, South Carolina 29208, <sup>13</sup> Stanford University, Stanford, California 94305, <sup>14</sup> Tel Aviv University, Tel Aviv, Israel, <sup>15</sup> Box 1290, Enderby, British Columbia, V0E 1V0, Canada, <sup>16</sup> Tufts University, Medford, Massachusetts 02155, <sup>17</sup> University of Wisconsin, Madison, Wisconsin 53706, <sup>18</sup> Yale University, New Haven, Connecticut 06511
August, 2000
## Abstract
From a sample of $`1172\pm 61`$ $`D^+\pi ^{}\pi ^+\pi ^+`$ decay, we find $`\mathrm{\Gamma }(D^+\pi ^{}\pi ^+\pi ^+)/\mathrm{\Gamma }(D^+K^{}\pi ^+\pi ^+)=0.0311\pm 0.0018_{0.0026}^{+0.0016}`$. Using a coherent amplitude analysis to fit the Dalitz plot of these decays, we find strong evidence that a scalar resonance of mass $`478_{23}^{+24}\pm 17`$ MeV/$`c^2`$ and width $`324_{40}^{+42}\pm 21`$ MeV/$`c^2`$ accounts for approximately half of all decays.
The three-body decays of charm mesons often proceed as quasi-two-body decays with resonant intermediate states. In our companion paper, we find that $`D_s^+f_0(980)\pi ^+`$ accounts for approximately half of the three-pion decay rate and $`D_s^+f_0(1370)\pi ^+`$ accounts for more than half of what remains, clearly establishing the dominance of isoscalar resonances in producing the three-pion final state. In this paper we present a study of the singly Cabibbo-suppressed decay $`D^+\pi ^{}\pi ^+\pi ^+`$. We determine the ratio of decay rates $`\mathrm{\Gamma }(D^+\pi ^{}\pi ^+\pi ^+)/\mathrm{\Gamma }(D^+K^{}\pi ^+\pi ^+)`$ and study the $`D^+\pi ^{}\pi ^+\pi ^+`$ Dalitz plot to determine the structure of its density distribution. We find that allowing an amplitude for an additional scalar state, with mass and width unconstrained, improves our fit substantially. The mass and width of the resonance found by this fit are $`478_{23}^{+24}\pm 17`$ MeV/$`c^2`$ and $`324_{40}^{+42}\pm 21`$ MeV/$`c^2`$. Referring to this $`\pi ^+\pi ^{}`$ resonance as the $`\sigma (500)`$, we find that $`D^+\sigma (500)\pi ^+`$ accounts for about half of the total decay rate.
Other experiments have presented controversial evidence for low-mass $`\pi \pi `$ resonances in partial wave analyses , with ambiguous results for the characteristics of such particles. Theoretically, light scalar and isoscalar resonances are predicted in models for spontaneous breaking of chiral symmetry such as the Nambu-Jona-Lasinio linear $`\sigma `$ model and its QCD extension . Also, these particles have important consequences for the quark model , for quark-gluon models , for understanding low energy $`\pi \pi `$ interactions , and for understanding the $`\mathrm{\Delta }I=1/2`$ rule.
This study is based on the sample of $`2\times 10^{10}`$ events recorded in Fermilab experiment E791, in which 500 GeV/$`c`$ $`\pi ^{}`$-nucleon interactions were observed using an open geometry spectrometer. The final analysis makes no direct use of particle identification; it is solely based on tracking and vertex reconstruction capabilities. To reduce background, we required a 3-prong decay (secondary) vertex to be well-separated from the production (primary) vertex, and located well outside of the target foils and other solid material. The momentum vector of the $`D`$ candidate had to point back to the primary vertex. A more detailed description of the final sample selection criteria and some additional details are provided in the companion paper where the resulting $`\pi ^{}\pi ^+\pi ^+`$ invariant mass distribution is shown as Fig. 1.
We have, in addition to the combinatorial background, three kinds of charm backgrounds: the reflection of the $`D^+K^{}\pi ^+\pi ^+`$ decay, located below 1.85 GeV/c<sup>2</sup> in this spectrum; the decay $`D^0K^{}\pi ^+`$ plus one extra track (mostly from the primary vertex); and the decay chain $`D_s^+\eta ^{}\pi ^+`$, $`\eta ^{}\rho ^0(770)\gamma `$, $`\rho ^0(770)\pi ^+\pi ^{}`$. The last two reflections populate the whole $`\pi ^{}\pi ^+\pi ^+`$ analyzed spectrum. We use Monte Carlo (MC) simulations and data to determine both the shape and the size of each type of charm background. The combinatorial background in the $`\pi ^{}\pi ^+\pi ^+`$ invariant mass spectrum is represented by an exponential function.
We fit the $`\pi ^{}\pi ^+\pi ^+`$ invariant mass distribution shown in Fig.1 of reference as the sum of $`D^+`$ and $`D_s^+`$ signals plus background. Each signal is described as the sum of two Gaussians with a common centroid but different widths, all these parameters determined by the fit. The background is represented by a function with four terms described above and in . The fit yields 1172 $`\pm `$ 61 $`D^+`$ events and 848 $`\pm `$ 44 $`D_s^+`$ events.
We measure the branching ratio for $`D^+\pi ^{}\pi ^+\pi ^+`$ relative to that of $`D^+K^{}\pi ^+\pi ^+`$. The $`K^{}\pi ^+\pi ^+`$ signal, selected with the same criteria used for the $`\pi ^{}\pi ^+\pi ^+`$ signal, is 34790 $`\pm `$ 232 events. The absolute efficiency, $`\epsilon `$, for each decay mode is approximately $`3\%`$. From Monte Carlo studies we determine that the ratio $`\epsilon (D^+\pi ^{}\pi ^+\pi ^+)/\epsilon (D^+K^{}\pi ^+\pi ^+)=1.08\pm 0.02`$. Note that we use the decay matrix element found in this analysis and an appropriate one for $`D^+K^{}\pi ^+\pi ^+`$ to determine the efficiencies. We also weight the MC production model to match our data. We find the relative branching ratio to be
$$\frac{\mathrm{\Gamma }(D^+\pi ^{}\pi ^+\pi ^+)}{\mathrm{\Gamma }(D^+K^{}\pi ^+\pi ^+)}=0.0311\pm 0.0018_{0.0026}^{+0.0016}.$$
(1)
The first error is statistical and the second is systematic. Uncertainties in the $`\pi ^{}\pi ^+\pi ^+`$ background shape and the levels of some contributions dominate the systematic error. This result can be compared with the measurements reported by E691 , $`0.035\pm 0.007\pm 0.003`$, by WA82 , $`0.032\pm 0.011\pm 0.003`$, and by E687 , $`0.043\pm 0.003\pm 0.003`$.
To study the resonant structure of the decay $`D^+\pi ^{}\pi ^+\pi ^+`$, we consider the 1686 candidates with invariant mass between 1.85 and 1.89 GeV. The integrated signal-to-background ratio in this range is about 2:1. Fig. 1 shows the Dalitz plot for these events. The horizontal and vertical axes are the squares of the $`\pi ^+\pi ^{}`$ invariant masses, and the plot has been symmetrized with respect to the two $`\pi ^+`$’s.
To study the resonant structure in Fig. 1, we use MINUIT to extract the parameters. We do this by maximizing the log (Likelihood) $``$ for several models of signal and background. For each model we compute $``$ in terms of signal and background probability distribution functions (PDF’s) of the $`\pi ^{}\pi ^+\pi ^+`$ invariant mass, $`M`$, and the Lorentz invariants $`s_{12}m_{12}^2`$ and $`s_{13}m_{13}^2`$ (in our convention the odd-charged pion is labeled particle $`1`$). Writing $`𝒫_S`$ and $`𝒫_B`$ for the signal and background PDF’s, $`=_{j=1}^{1686}[𝒫_S+𝒫_B]_j`$. We take $`𝒫_S=\frac{1}{N_S}g(M)\epsilon (s_{12},s_{13})𝒜^2`$, with
$$𝒜=a_0e^{i\delta _0}𝒜_0+\underset{n=1}{\overset{N}{}}a_ne^{i\delta _j}𝒜_n(s_{12},s_{13}).$$
(2)
In this equation $`N_S`$ is the normalization constant, $`\epsilon (s_{12},s_{13})`$ is the net efficiency, $`g(M)`$ is a Gaussian function describing the signal $`\pi ^{}\pi ^+\pi ^+`$ invariant mass spectrum. The fit parameters are the coefficient magnitudes, $`a_n`$, and the phases, $`\delta _n`$.
The non-resonant amplitude, $`𝒜_0`$, is represented by a constant. Each resonant amplitude, $`𝒜_n(n1)`$, is written as a product of four terms.
$$𝒜_n=^JF_n\times {}_{}{}^{J}_{n}^{}\times BW_n$$
(3)
The first term is form factor for the $`n^{th}`$ resonance We assume the form factor for the $`D`$ decay to be one. $`{}_{}{}^{J}_{n}^{}`$ is a term which accounts for angular-momentum conservation and depends on the spin J of the resonance. The final term is a relativistic Breit-Wigner function given by:
$$BW_n=\frac{1}{m^2m_0^2+im_0\mathrm{\Gamma }_n(m)}$$
(4)
with
$$\mathrm{\Gamma }(m)=\mathrm{\Gamma }_0\frac{m_0}{m}\left(\frac{p^{}}{p_0^{}}\right)^{2J+1}\frac{{}_{}{}^{J}F_{n}^{2}(p^{})}{{}_{}{}^{J}F_{n}^{2}(p_0^{})}.$$
(5)
In Eqs. 3 and 4, $`m`$ is the invariant mass of the two pions forming a spin-J resonance. The functions $`{}_{}{}^{J}F`$ are the Blatt-Weisskopf damping factors : $`{}_{}{}^{0}F=1`$ for spin 0 particles, $`{}_{}{}^{1}F=1/\sqrt{1+(rp^{})^2}`$ for spin 1 and $`{}_{}{}^{2}F=1/\sqrt{9+3(rp^{})^2+(rp^{})^4}`$ for spin 2. The parameter $`r`$ is the radius of the resonance ($`3fm`$) and $`p^{}=p^{}(m)`$ the momentum of decay particles at mass $`m`$, measured in the resonance rest frame, $`p_0^{}=p^{}(m_0)`$, where $`m_0`$ is the resonance mass. The spin part of the amplitude $`{}_{}{}^{J}_{n}^{}`$ is defined equal to 1 for a spin-0 resonance, $`2𝐩_\mathrm{𝟑}𝐩_\mathrm{𝟐}cos\theta `$ for spin-1 and $`\frac{4}{3}(𝐩_\mathrm{𝟑}𝐩_\mathrm{𝟐})^2(3cos^2\theta 1)`$, where $`𝐩_\mathrm{𝟑}`$ is the 3-momentum of the unlike-charge pion and $`𝐩_\mathrm{𝟐}`$ is the 3-momentum of the other like-charge pion, both measured in the resonance rest frame; and $`\theta `$ is the angle between pions 2 and 3. Finally, each resonant amplitude is Bose symmetrized: $`𝒜_n=𝒜_n[(\mathrm{𝟏𝟐})\mathrm{𝟑}]+𝒜_n[(\mathrm{𝟏𝟑})\mathrm{𝟐}]`$.
The background distribution is given by $`𝒫_B=b(M)_{i=1}^3\frac{b_i}{N_{B_i}}_i(s_{12},s_{13})`$; $`b(M)`$ is the function describing the background distribution in the $`\pi ^{}\pi ^+\pi ^+`$ spectrum, $`b_i`$ are the relative amount of each background type and $`N_{B_i}`$ are the corresponding normalization constants.
The three components of the background distribution are the combinatorial background, assumed to be uniform before any acceptance effects, and the $`D^0K^{}\pi ^+`$ and $`D_s^+\eta ^{}\pi ^+`$ reflections. The relative background fractions are $`80\pm 6\%`$, $`4\pm 1\%`$, and $`16\pm 6\%`$, respectively. The shape, location, and size of the charm background were obtained using MC simulations and previously determined $`D^0`$ and $`D_s`$ production rates relative to $`D^+`$ in our data sample. All parameters used for the background description are fixed during the fit.
The dominance of the above three background contributions was checked in several tests. The analysis was repeated with more stringent selection criteria, with various levels of Čerenkov-counter requirements on the pions, and with varied background levels in the fit. All test results were consistent with the quoted final results. In addition, Monte Carlo simulations were used to study specific charm decay channels and a generic sample of charm decays. The latter was examined to look for Dalitz-plot structure in the events which passed our final selection criteria. No structures, other than those noted above, were significant.
In a first model, which we will refer to as Fit 1, the signal PDF includes a non-resonant amplitude and amplitudes for $`D^+`$ decaying to a $`\pi ^+`$ and any of five established $`\pi ^+\pi ^{}`$ resonances: $`\rho ^0(770)`$, $`f_0(980)`$, $`f_2(1270)`$, $`f_0(1370)`$, and $`\rho ^0(1450)`$. In the case of the $`f_0(980)`$ and $`f_0(1370)`$, we used the parameters of Ref. and not those of Ref. . The fit extracts the magnitudes and phases of each of the amplitudes along with the error matrix for these parameters. We calculate the decay fraction for each amplitude as its intensity, integrated over the Dalitz plot, divided by the integrated intensity of the signal’s coherently summed amplitudes. The Fit 1 results are listed in the first column of Table 1. In this model, the non-resonant, the $`\rho ^0(1450)\pi ^+`$ and the $`\rho ^0(770)\pi ^+`$ amplitudes dominate. The qualitative features of this fit are similar to those reported by E691 and E687.
To assess the quality of the fit, we developed a fast-MC program which produces binned Dalitz-plot densities accounting for signal and background PDF’s, including detector efficiency and resolution. Comparing the binned Dalitz-plot-density distribution generated by MC events using the magnitudes and phases of the amplitudes given in Fit 1 with that for the data, we produce the $`\chi ^2`$ distribution for the difference in densities and observe a concentration of a large $`\chi ^2`$ in the low $`\pi ^+\pi ^{}`$ mass ($`m_{\pi ^+\pi ^{}}`$) region. The $`\chi ^2`$ summed over all bins is 254 for 162 degrees of freedom ($`\nu `$), which corresponds to a confidence level less than $`10^5`$, assuming Gaussian errors. Since the two $`m_{\pi ^+\pi ^{}}^2`$ projections are nearly independent, we display the sum of $`s_{12}`$ and $`s_{13}`$ in Fig. 2a for the data and for the fast-MC.
The small value of the confidence level casts doubt on the validity of the model used. While the projection of the MC onto the $`\pi ^+\pi ^{}`$ mass<sup>2</sup> axis describes the data in the $`\rho ^0(770)`$ and $`f_0(980)`$ regions well, there is a discrepancy at lower mass, suggesting the possibility of another amplitude.
To investigate the possibility that another $`\pi ^+\pi ^{}`$ resonance contributes an amplitude to the $`D^+\pi ^{}\pi ^+\pi ^+`$ decay, we add a sixth resonant amplitude to the signal PDF. We allow its mass and width to float and assume a scalar angular distribution. This fit (Fit 2) converges and finds values of $`478_{23}^{+24}`$ MeV$`/c^2`$ for the mass and $`324_{40}^{+42}`$ MeV$`/c^2`$ for the width. We will refer to this possible state as the $`\sigma (500)`$. The corresponding results, including the systematic errors , are collected in the second column of Table 1.
In Fit 2, the $`\sigma (500)`$ amplitude produces the largest decay fraction, 46%, with a relatively small statistical error, 9%. The non-resonant fraction, which at $`(39\pm 10)\%`$ was the largest in the original fit, is now only $`(8\pm 6)\%`$. When we project this model onto the Dalitz plot, the $`\chi ^2/\nu `$ becomes 138/162. The projection of this model onto the $`\pi ^+\pi ^{}`$ invariant mass squared distribution, shown in Fig. 2b, describes the data well, including the accumulation of events near 0.2 GeV$`{}_{}{}^{2}/c^4`$. We tested the accuracy of the fit’s error estimates by producing hundreds of fast-MC samples using the parameters of Fit 2 and then fitting the samples. The central values of all the parameters are reproduced accurately, and the width of the $`\sigma (500)`$ fraction distribution is 0.12, slightly larger than MINUIT’s estimate, 0.09.
When comparing the two models we used the fast-MC to simulate an ensemble of samples for each model. For each sample we calculated $`\mathrm{\Delta }w=2(\mathrm{ln}_i\mathrm{ln}_\sigma )`$, where $`_\sigma `$ and $`_i`$ are the likelihood functions evaluated with the parameters from the fit with and without the $`\sigma \pi ^+`$ amplitude, respectively. In the data, $`\mathrm{\Delta }w=118`$; in the fast MC with a $`\sigma `$, $`\mathrm{\Delta }w=108`$; in the fast MC with no $`\sigma `$, $`\mathrm{\Delta }w=106`$. In both MC experiments, the rms deviation for $`\mathrm{\Delta }w`$ is about 20, indicating a strong preference for the additional amplitude, as is indicated also by the difference in $`\chi ^2/\nu `$ for the two models.
We consider the systematic errors associated with the values of the fixed parameters in the fit. The most important ones come from uncertainties in the background model (the background shape, composition, and level), in particular the $`D_s^+\eta ^{}\pi ^+`$ reflection which populates the same region as the $`D^+\rho ^0\pi ^+`$ component. We also account for the uncertainties in the parameters describing the acceptance function.
To better understand our data, we also fit it with vector, tensor, and toy models for the sixth (sigma) amplitude, allowing the masses, widths, and relative amplitudes to float freely. The vector and tensor models test the angular distribution of the signal. The toy model tests the phase variation expected of a Breit-Wigner amplitude by substituting a constant relative phase. The vector resonance model converges to poorly defined values of the mass and width: $`805\pm 194`$ MeV/$`c^2`$ and $`1438\pm 903`$ MeV/$`c^2`$; the tensor model to more poorly defined values: $`2350\pm 683`$ MeV/$`c^2`$ and $`690\pm 1033`$; and the toy model to $`434\pm 11`$ MeV/$`c^2`$ and $`267\pm 37`$ MeV/$`c^2`$. As a test of the models, we again project the vector, tensor, and toy models onto the Dalitz plot and obtain $`\chi ^2/\nu `$ = 188/162, 148/162, and $`152/162`$, respectively. For these models, we also find $`\mathrm{\Delta }w`$ = 66, 13, and 15 where MC experiments predict $`\mathrm{\Delta }w64`$, 56, and 38 when the data is generated with the scalar parameters and the negatives of those values when the MC data is generated according to the vector, tensor, and toy model parameters. The rms widths of the MC distributions are 15, 15, and 11 units respectively. These statistical tests strongly exclude the vector model. They clearly prefer the scalar model to the tensor and toy models. Note that the central value for the tensor mass is well above threshold for $`D^+`$ decay and the negative width is an indication that no physically meaningful tensor resonance fits the data. In the toy model, the extra amplitude interferes strongly with a large non-resonant amplitude, leading to an unphysically large sum of resonant fractions.
In summary, from 1172 $`\pm `$ 61 $`D^+\pi ^{}\pi ^+\pi ^+`$ we have measured $`\mathrm{\Gamma }(D^+\pi ^{}\pi ^+\pi ^+)/\mathrm{\Gamma }(D^+K^{}\pi ^+\pi ^+)=0.0311\pm 0.0018_{0.0026}^{+0.0016}`$. In an amplitude analysis of a sample with S:B $``$ 2:1 we find strong evidence that a scalar resonance with mass $`478_{23}^{+24}\pm 17`$ MeV/$`c^2`$ and width $`324_{40}^{+42}\pm 21`$ MeV/$`c^2`$ produces a decay fraction $`50\%`$. Alternative explanations of the data fail to describe it as well. The prominence of an amplitude for an isoscalar plus a $`\pi ^+`$ in this decay accords well with our observation that the amplitude for an isoscalar plus a pion produce a large majority of all $`D_s^+\pi ^{}\pi ^+\pi ^+`$ decays.
We gratefully acknowledge the assistance of the staffs of Fermilab and of all the participating institutions. This research was supported by the Brazilian Conselho Nacional de Desenvolvimento Científico e Tecnológico, CONACyT (Mexico), the Israeli Academy of Sciences and Humanities, the U.S. Department of Energy, the U.S.-Israel Binational Science Foundation, and the U.S. National Science Foundation. Fermilab is operated by the Universities Research Association, Inc., under contract with the U.S. Department of Energy.
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# 1 Conjectured phase diagram of overlap-Dirac fermions for Nₗ≥2. The number of lattice pions in each phase is indicated.
July 2000 TAUP–2632–00
UTCCP-P-89
Overlap-Dirac fermions with a small hopping parameter
Maarten Golterman<sup>a</sup> and Yigal Shamir<sup>b</sup>
> <sup>a</sup>Center for Computational Physics, University of Tsukuba,
> Tsukuba, Ibaraki 305-8577, JAPAN
> and
> Department of Physics, Washington University, St. Louis, MO 63130, USA
> maarten@aapje.wustl.edu
> <sup>b</sup>School of Physics and Astronomy, Beverly and Raymond Sackler Faculty of Exact Sciences, Tel-Aviv University, Ramat Aviv, 69978 ISRAEL
> shamir@post.tau.ac.il
ABSTRACT
> We consider overlap-Dirac fermions at non-zero bare coupling and for a small hopping parameter, or, equivalently, large $`|M|`$ with $`M`$ the domain-wall height. We prove the existence of a phase at large positive $`M`$ where the abelian axial group $`U_A(1)`$ is a symmetry, and the corresponding pseudo-scalar is an exact Goldstone boson. We also provide a conjecture for the phase diagram of asymptotically free gauge theories with overlap-Dirac fermions. In particular, we suggest that, for large gauge coupling, the massive-fermion phase at negative $`M`$ possibly extends to all $`M<4`$.
Permanent address
1. Introduction. In recent years it has become clear that the chiral limit can be separated from the continuum limit in QCD-like lattice gauge theories \[1-15\]. Domain-wall fermions \[1-3\] provide a simple way of approaching this limit, and have become popular in numerical simulations (see ref. for a review). On the theoretical side, a key role in the formulation of the chiral limit is played by the algebraic Ginsparg-Wilson relation (see ref. for a review). An extensively studied solution of the Ginsparg-Wilson relation is the overlap-Dirac operator , which may be regarded as a certain limiting case of domain-wall fermions .
A common feature of all numerical simulations which attempt to approach the chiral limit is that they show a pattern of increasing chiral symmetry violations when the bare coupling is increased. Keeping such symmetry violations under control in the most economic way is important, and analytic results valid beyond the weak-coupling limit can provide valuable clues as to what is the best way of achieving this goal.
In this paper we consider a euclidean $`SU(N_c)`$ lattice gauge theory with $`N_l`$ copies of overlap-Dirac fermions and with a finite bare coupling $`g`$. For a small hopping parameter $`\kappa `$, the theory may be in a massive or in a massless phase (depending on whether $`A=1`$ or $`A=1`$ in eq. (3) below). Our main result is that, in the massless phase, the theory may be reformulated such that the abelian axial group $`U_A(1)`$ becomes a manifest symmetry of both the action and the measure, without spoiling the locality of the theory. As a result, there are $`N_l^2`$ lattice pions in that massless phase. The term lattice pion here denotes a pseudo-scalar state which is an exact Goldstone boson for finite $`g`$.
To avoid additional complications which are besides our main point, we mainly consider lattice theories where the limit $`g0`$ defines a confining continuum theory. We then expect chiral symmetry to be spontaneously broken in the continuum limit as well . (Basically this means that the number of flavors is small enough compared to $`N_c`$; see below for a comment on the case $`N_c=3`$.) If the continuum theory has $`N_f`$ flavors, the number of pions is $`N_f^21`$. The small-$`\kappa `$ phase referred to above thus has one extra (lattice) pion compared to QCD with $`N_f=N_l`$. This indicates some sort of species doubling, and we will argue that, if the continuum limit is taken inside that phase, the number of flavors is actually $`N_f=16N_l`$, as for naive fermions.
Our work was motivated by two recent strong-coupling calculations . The results of those papers will be compared with ours in the last section.
This paper is organized as follows. In Sect. 2 we review some relevant properties of the overlap-Dirac operator, discussing in particular the physical significance of Lüscher’s axial transformations . In Sect. 3 we prove the main result, and in Sect. 4 we comment on the role of the index of the overlap-Dirac operator. In Sect. 5 we discuss the continuum limit in some more detail. In Sect. 6 we turn the existing results into a conjecture on the phase diagram of overlap fermions. Finally, in Sect. 7 we compare our results with those of refs. , and list some issues for future research. Locality of the reformulated action is proved in the Appendix.
2. The overlap-Dirac operator. The Ginsparg-Wilson relation is (we work in units of the lattice spacing)
$$D\gamma _5+\gamma _5D=D\gamma _5D.$$
(1)
The overlap-Dirac operator which satisfies the Ginsparg-Wilson relation is defined by
$$D=1+\gamma _5H/|H|,$$
(2)
$$H=\gamma _5X,X=A+\kappa \left(\underset{\mu }{}\gamma _\mu C_\mu B\right),$$
(3)
where $`X`$ is the usual Wilson-Dirac operator, and
$$(C_\mu )_{x,y}=\frac{1}{2}(\delta _{x+\widehat{\mu },y}U_{x,\mu }\delta _{x\widehat{\mu },y}U_{y,\mu }^{}),$$
(4)
$$B_{x,y}=\frac{1}{2}\underset{\mu }{}(\delta _{x+\widehat{\mu },y}U_{x,\mu }+\delta _{x\widehat{\mu },y}U_{y,\mu }^{}).$$
(5)
Since $`D`$ is unchanged if $`H`$ is multiplied by an arbitrary positive number, we use this freedom to set $`A=\pm 1`$ in eq. (3). The hopping parameter $`\kappa `$ is positive by convention. (A sign flip of $`\kappa `$ can be undone by the transformation $`\psi _x(1)^{x_1+x_2+x_3+x_4}\psi _x`$ and similarly for $`\overline{\psi }_x`$.) We take the Wilson parameter to be $`r=1`$, but the discussion can easily be generalized to other values.
The action $`S(\psi ,\overline{\psi })=\overline{\psi }D\psi `$ is invariant under Lüscher’s gauge-field dependent axial transformation
$`\delta \psi `$ $`=`$ $`T\widehat{\gamma }_5\psi ,\widehat{\gamma }_5=\gamma _5(1D)=H/|H|,`$
$`\delta \overline{\psi }`$ $`=`$ $`\overline{\psi }T\gamma _5.`$ (6)
Here $`T`$ is a $`U(N_l)`$ generator acting on the flavor indices. Since $`(\widehat{\gamma }_5)^2=1`$, Lüscher’s axial transformations together with the usual vector transformations generate a $`U_L(N_l)\times U_R(N_l)=U_V(1)\times U_A(1)\times SU_L(N_l)\times SU_R(N_l)`$ symmetry of the lattice action. This statement is true for any $`\kappa `$. To avoid confusion we recall that if $`\kappa =O(1)`$ the fermion measure is in general not invariant under Lüscher’s $`U_A(1)`$, leading to the expected axial anomaly and a massive $`\eta ^{}`$ particle.
The symmetry (6) does not always have the physical significance of an axial symmetry. Let us first recall the quark spectrum described by a single overlap-Dirac field. As usual, the quark spectrum is determined by the free theory. We then have for $`B_0`$ and $`C_0`$, in momentum space,
$$C_{0\mu }(p)=i\mathrm{sin}p_\mu ,B_0(p)=\underset{\mu }{}\mathrm{cos}p_\mu .$$
(7)
It is easy to check that at the corners $`p_c`$ of the Brillouin zone ($`p_{c\mu }=0`$ or $`\pi `$, all $`\mu `$), and only there, $`D_0(p_c)`$ is equal to either 0 or 2. A given corner of the Brillouin zone gives rise to a massless quark field in the continuum limit if and only if $`D_0(p_c)=0`$.
In the context of domain-wall or overlap-Dirac fermions, the customary parametrization of the Wilson-Dirac operator in eq. (3) is
$$X=4M+\underset{\mu }{}\gamma _\mu C_\mu B.$$
(8)
Comparing eqs. (3) and (8), we see that $`\kappa =1/|4M|`$, and that $`A=1`$ ($`A=1`$) corresponds to $`M<4`$ ($`M>4`$). When $`n`$ components of the four-momentum are equal to $`\pi `$ and the rest are zero, one has $`B_0(p_c)=42n`$. Referring to the parametrization of eq. (8), it follows that for $`M<0`$ all corners have $`D_0(p_c)=2`$, while for $`M>8`$ all corners have $`D_0(p_c)=0`$. For $`M<0`$ there are no massless quarks, and when $`M`$ is increased above the values $`0,2,4,6`$ and $`8`$, the numbers of massless quarks that are added to the spectrum are $`1,4,6,4`$ and $`1`$ respectively. For $`M>8`$ there are 16 massless quarks. Note that the points $`M=0,2,\mathrm{},8,`$ represent discontinuities in the spectrum.
Suppose now that $`D_0(p_c)=0`$ at some corner of the Brillouin zone. Near that corner Lüscher’s transformation (6) reduces to an ordinary axial transformation (cf. eq. (10) below) when acting on the corresponding massless-quark state. We recall that the physical axial charge is equal or opposite to the lattice axial charge depending on whether $`n`$ is even or odd, respectively .
The other possibility is that $`D_0(p_c)=2`$ at a corner. In this case, eq. (6) reduces to $`\delta \psi =T\gamma _5\psi `$, $`\delta \overline{\psi }=\overline{\psi }T\gamma _5`$. Because of the minus sign in the $`\psi `$ transformation rule, this is no longer an axial transformation. There exists another version of Lüscher’s transformations where the $`\overline{\psi }`$ and $`\psi `$ rules look more symmetric, given by $`\delta \psi =T\gamma _5(1\frac{1}{2}D)\psi `$, $`\delta \overline{\psi }=\overline{\psi }T(1\frac{1}{2}D)\gamma _5`$. In this form, the transformation still reduces to an ordinary axial transformation when $`D_0(p_c)=0`$. But for $`D_0(p_c)=2`$ this becomes $`\delta \psi =\delta \overline{\psi }=0`$. In other words, for $`D_0(p_c)=2`$ the transformation does not act at all on states with a momentum close to $`p_c`$.
One can summarize the situation by saying that, unlike an ordinary axial symmetry, Lüscher’s symmetry by itself does not imply the existence of massless quarks. But if massless quarks exist, it acts on them as an ordinary axial symmetry. In fact, these requirements single out the transformations (6) almost uniquely. For consider taking $`\delta \psi =T\gamma _5(1𝒪)\psi `$ for some $`𝒪`$, with $`𝒪=0`$ for $`p=0`$, but possibly non-zero for other momenta in order to accommodate the removal of doublers. We also want to form the same Lie algebra as in the continuum, so we need $`(\gamma _5(1𝒪))^2=1`$. From this it follows immediately that $`𝒪`$ satisfies the Ginsparg-Wilson relation.
3. The massless small-$`\kappa `$ phase. The properties of $`D`$ for a small hopping parameter $`\kappa `$ rest on the fact that $`X=\pm 1+O(\kappa )`$ (cf. eq. (3)). It follows immediately that $`D=2O(\kappa )`$ for $`A=1`$ while $`D=O(\kappa )`$ for $`A=1`$ (these statements refer to the eigenvalues of $`D`$; we recall that $`|D|2`$ always, see ref. ). In the case $`A=1`$, since $`D=2O(\kappa )`$ regardless of $`g`$, all correlation lengths in the fermion sector must be finite in lattice units, and moreover tend to zero for $`\kappa 0`$. Hence $`A=1`$ corresponds to a massive-fermion phase. (This is the usual situation in a hopping expansion. It can be established using the methods of ref. , or, more explicitly, by applying the techniques of the Appendix to obtain bounds on the kernel of $`D^1`$ for $`A=1`$ and small $`\kappa `$.)
In this section we will show that $`A=1`$ corresponds to a massless phase with $`N_l^2`$ lattice pions. The first important observation is that the leading, $`O(\kappa )`$, term in the expansion of $`D`$ is proportional to the naive-fermion operator. Also, we have already seen that for small $`\kappa `$ (large $`M`$) overlap-Dirac fermions undergo the same, maximal, doubling in the continuum limit as naive fermions. (However, overlap-Dirac and naive fermions have a different massless spectrum (of gauge invariant states) at finite $`\kappa `$ and $`g`$, see below. Note that the $`\kappa 0`$ limit is singular for normalized expectation values: in this limit there are massless fermions, whereas setting $`\kappa =0`$ gives rise to no propagation at all. The hopping expansion in the $`A=1`$ phase is thus qualitatively different from that for Wilson fermions, where, as in the $`A=1`$ phase, it is an expansion around an infinite fermion-mass theory.)
In order to analyze the situation for a small but finite hopping parameter it is convenient to introduce new variables
$$q=(2D)\psi ,\overline{q}=\overline{\psi }.$$
(9)
In terms of the new variables, Lüscher’s transformation (6) reduces to the ordinary axial transformation
$$\delta q=T\gamma _5q,\delta \overline{q}=\overline{q}T\gamma _5.$$
(10)
The partition function is rewritten as
$`Z=\mathrm{det}(D)`$ $`=`$ $`{\displaystyle 𝒟\psi 𝒟\overline{\psi }\mathrm{exp}(S(\psi ,\overline{\psi }))}`$ (11)
$`=`$ $`{\displaystyle 𝒟q𝒟\overline{q}\mathrm{det}(2D)\mathrm{exp}(S^{}(q,\overline{q}))},`$
where
$$S^{}(q,\overline{q})=\overline{q}D(2D)^1q.$$
(12)
It is easily verified that $`D(2D)^1`$ anticommutes with $`\gamma _5`$, as it must, in view of eq. (10). The action $`S^{}(q,\overline{q})`$ is thus invariant under ordinary axial transformations, as well as under the usual hyper-cubic rotations.
The transition to the new variables cannot be done for arbitrary $`\kappa `$ and $`A`$. When $`\kappa `$ exceeds some critical value, $`(2D)^1`$ becomes singular, and in the free fermion limit there are poles in the action $`S^{}(q,\overline{q})`$. Moreover, in non-trivial topological sectors $`S^{}(q,\overline{q})`$ is undefined since there exist stable eigenmodes with eigenvalues 0 and 2. Finally, when $`A=1`$, poles appear in $`S^{}(q,\overline{q})`$ already in the hopping expansion, and therefore the discussion below is not applicable in the massive phase.
When $`A=1`$, the transition to the new variables turns out to be a powerful tool, since in this case eq. (12) is a local action for small enough $`\kappa `$. In more detail, for a range $`\kappa <\kappa _0`$ of the hopping parameter $`D(2D)^1`$ is bounded and has an exponentially decaying kernel (as is the standard practice in this context, the last statement defines the notion of locality used in this paper). Boundedness is obvious since $`D=O(\kappa )`$, while locality of $`D(2D)^1`$ is established in the Appendix for $`\kappa <\kappa _0`$ (we note that the “admissibility” constraint on the gauge field is not necessary for small $`\kappa `$). As for the factor $`\mathrm{det}(2D)`$ in eq. (11), it is not expected to change the universality class (an argument similar to that in the appendix shows that $`\mathrm{tr}\mathrm{log}(2D)`$ is local).
Thus, we find that for $`\kappa <\kappa _0`$ the theory can be consistently formulated in terms of a local action $`S^{}(q,\overline{q})`$, where the action and, obviously, the measure are invariant under ordinary vector and axial transformations. The fact that we are dealing here with the simple axial transformation of eq. (10) is important. As mentioned in the introduction, we restrict the discussion to those cases where the limit $`g0`$ defines a confining continuum theory. We then have confinement for any $`g`$, and the standard arguments that confinement implies chiral symmetry breaking apply . In the present context, the formation of a $`\overline{q}q`$ condensate was confirmed in the strong coupling limit in a $`1/N_c`$ expansion (see also ref. ).
Therefore, the lattice theory has $`N_l^2`$ axial generators, and the pseudo-scalar Goldstone bosons must be in one-to-one correspondence with those. We thus conclude that there are $`N_l^2`$ lattice pions for $`\kappa <\kappa _0`$ and any (finite) $`g`$. The (confining) small-$`\kappa `$ massless phase may actually be defined as the part of the phase diagram with $`N_l^2`$ lattice pions. We expect this phase to extend beyond the region $`\kappa <\kappa _0`$ and, in fact, beyond the region where $`(2D)^1`$ is bounded, see below. We comment in passing that in the case of naive fermions the lattice symmetry is bigger , and hence the number of lattice pions is larger as well. This naive-fermion symmetry requires that the action has only odd-neighbor couplings.
In summary, we have established the existence of a class of overlap-Dirac theories where for $`A=1`$ and small $`\kappa `$ there is only one phase for all values of $`g`$. This phase is characterized by confinement and chiral symmetry breaking, and has $`N_l^2`$ lattice pions.
4. The index and $`U_A(1)`$. The question of whether $`U_A(1)`$ is a symmetry or not may be approached from another direction. The integrated Ward identity associated with a $`U_A(1)`$ transformation reads
$$\delta 𝒪=\mathrm{Tr}(\widehat{\gamma }_5)𝒪,$$
(13)
where $`𝒪=𝒪(\psi ,\overline{\psi },U_\mu )`$ and the axial variation denoted by $`\delta `$ is defined in eq. (6). (To make the above statement well defined we may assume that we work in a finite volume.) It is easy to see that $`\mathrm{Tr}(\widehat{\gamma }_5)=0`$ for small $`\kappa `$ for all gauge field configurations, regardless of the sign of $`A`$ in eq. (3). First, $`\mathrm{Tr}(\widehat{\gamma }_5)`$ is (proportional to) the index of $`D`$ . Then, a non-zero index requires the simultaneous existence of eigenvectors with eigenvalues zero and two . But, for $`A=1`$, $`D=O(\kappa )`$ and there can be no eigenvalue equal to two. For $`A=1`$, $`D=2O(\kappa )`$ and there can be no zero eigenvalue.
What we learn is that the small-$`\kappa `$ global symmetry of both the $`A=1`$ and the $`A=1`$ phases is $`U_L(N_l)\times U_R(N_l)`$. However, as already explained above, there is an important difference between the two phases. In the $`A=1`$ phase the conserved axial generators may be taken to be the ordinary axial generators associated with eq. (10). We have seen that this implies the existence of massless quarks and (if there is confinement) spontaneous chiral symmetry breaking. In contrast, in the $`A=1`$ phase we have only Lüscher’s symmetry at our disposal, and as explained in Sect. 2, this is not incompatible with the absence of massless-fermion states.
5. Relation to the continuum limit. The previous results provide valuable information for the task of mapping out the $`(M,g)`$ phase diagram of overlap-Dirac fermions. As already mentioned, the hopping expansion is an expansion in $`1/|4M|`$. The massless ($`A=1`$) phase at small $`\kappa `$ corresponds to $`M>M_{c+}(g)`$, while the massive ($`A=1`$) phase corresponds to $`M<M_c(g)`$.
We now wish to determine the end points of these two critical lines in the continuum limit. We will argue that the massless phase at large positive $`M`$ extends down to $`M=8`$ for $`g0`$, whereas the (massive-fermion) phase at large negative $`M`$ extends up to $`M=0`$ provided $`N_l2`$. The region $`O(g)<M<8O(g)`$ at small $`g`$ is filled with a phase supporting $`N_l^21`$ lattice pions.
Let us begin with the last statement. For $`0<M<8`$ there are massless quarks, and since Lüscher’s symmetry (6) acts on those as an axial symmetry, spontaneous symmetry breaking should occur, and the corresponding Goldstone bosons should exist. In this range no massless state in the $`U_A(1)`$ channel is expected because the index of $`D`$ can be non-zero. Since there are $`N_l^21`$ conserved axial generators for non-zero $`g`$, this must also be the number of lattice pions. A corollary is that the massless large-$`M`$ phase with $`N_l^2`$ lattice pions cannot end at any $`M<8`$ for $`g0`$. Recall that, according to the free-field (or weak-coupling) analysis, the largest value where a phase transition takes place in the $`g0`$ limit is $`M=8`$. Assuming that that analysis exhausts all possible phase transition points for $`g0`$, it follows that the large-$`M`$ phase must end at $`M=8`$. Similarly, the massive-fermion phase with no lattice pions at negative $`M`$ (which is actually a pure glue phase) should end at $`M=0`$ for $`g0`$. The last statement is true except for $`N_l=1`$, where both the negative-$`M`$ and intermediate-$`M`$ regions do not support lattice pions, and therefore they may be analytically connected at non-zero $`g`$.
For small $`g`$ and $`M>8`$ there are massless states at all the corners of the Brillouin zone, namely 16 quark fields per each overlap-Dirac fermion. The lattice $`U_A(1)`$ symmetry corresponds to a flavor non-diagonal axial symmetry in the continuum limit (see Sect. 2). The total number of pions in the continuum limit is $`(16N_l)^21`$, whereas for finite $`g`$ there are only $`N_l^2`$ pions. The difference is explained by states which are approximate Goldstone bosons for small $`g`$, and which become massless only in the continuum limit. This bears some resemblance to staggered fermions, where only one pion is exactly massless at finite lattice spacing, while the rest become massless only in the continuum limit.
If one relaxes the assumption we have made in the introduction about the fermion content, one can find cases where asymptotic freedom is lost and/or there is no chiral symmetry breaking for small $`g`$. In such cases there has to be (at least) one phase transition as a function of $`g`$ even at large (positive) $`M`$. For the physical case $`N_c=3`$, asymptotic freedom at large $`M`$ is lost for $`N_l=2`$ ($`N_f=32`$), and it is not clear if there is (confinement and) chiral symmetry breaking even for $`N_l=1`$ ($`N_f=16`$) in the continuum limit. At strong coupling one always has confinement, and we expect our conclusion to hold on the large-$`g`$ side of such a transition. Also, whatever replaces confinement and chiral symmetry breaking at small $`g`$ should still be consistent with having $`16N_l^2`$ massless fermions for $`g=0`$. Note that this issue does not affect the phase diagram in the region of interest for Lattice QCD, with $`0<M<2`$.
6. Conjectured phase diagram of overlap-Dirac fermions. In this section we put forward a conjecture for the $`(M,g)`$ phase diagram of overlap-Dirac fermions. We have already discussed three out of four boundaries of the phase diagram. We invoke a mean-field argument to cover the last, $`g=\mathrm{}`$, boundary, and then consider the simplest finite-$`(M,g)`$ interpolation.
The conjectured phase diagram is depicted in Fig. 1 for $`N_l2`$. It contains the two phases with zero and $`N_l^2`$ lattice pions whose existence was established for large-negative and large-positive $`M`$ respectively. There is also a single phase at intermediate values of $`M`$, characterized by the existence of $`N_l^21`$ lattice pions. The phase transition lines emanating from the points $`M=2,4,6`$ must be there because the quark spectrum in the continuum limit changes at those points. As explained earlier, the regions surrounding these points at small $`g>0`$ all support $`N_l^21`$ lattice pions. In our view the most plausible scenario is that these regions are analytically connected at non-zero $`g`$, but this does not have to be the case. The phase transitions are discontinuous at $`g=0`$, and therefore they are discontinuous for small non-zero $`g`$, by continuity. We expect them to be discontinuous for all $`g`$. We remind the reader that this phase diagram is applicable provided there is a confining continuum theory in the limit $`g0`$ for all $`M`$.
If $`D`$ satisfies the Ginsparg-Wilson relation, so does $`2D`$. The passage from $`D`$ to $`2D`$ is effected by $`XX`$, and the phase diagram for $`2D`$ is obtained by the replacement $`M8M`$. (In the $`D`$-phase diagram, however, the critical lines may not transform into each other under $`M8M`$, and in that sense Fig. 1 may be misleading.)
The symmetries of the three phases are the following. (Below $`GH`$ denotes that the global symmetry is $`G`$ and that $`H`$ is the symmetry of the vacuum.) The phase with $`N_l^2`$ lattice pions is characterized by $`U_L(N_l)\times U_R(N_l)U_V(N_l)`$. The phase with $`N_l^21`$ lattice pion is characterized by $`U_V(1)\times SU_L(N_l)\times SU_R(N_l)U_V(N_l)`$. Finally, the phase with no lattice pions has the global symmetry $`U_L(N_l)\times U_R(N_l)`$ but no massless fermions, and no spontaneous symmetry breaking.
As can be seen from Fig. 1, we believe that the two small-$`\kappa `$ phases meet at $`M=4`$ for $`g=\mathrm{}`$. If this is true, the phase with $`N_l^21`$ lattice pion does not extend to $`g=\mathrm{}`$. (It could be that the triple point is at $`g=\mathrm{}`$. Note that one is usually interested in taking the continuum limit inside the middle phase because only in the interval $`0<M<2`$ there is one quark per one overlap fermion in the continuum limit.) Let us define $`U_{x\mu }`$ in a gauge-invariant way, for instance as the fourth root of the average plaquette. Then we have that $`U_{x\mu }0`$ for $`g\mathrm{}`$ and, if a mean-field analysis is reliable, the $`g\mathrm{}`$ limit should behave like $`\kappa 0`$. For $`M>4`$ this means that we are in the massless small-$`\kappa `$ phase with $`A=1`$, which has $`N_l^2`$ lattice pions, and for $`M<4`$ we are in the $`A=1`$ massive phase, with no lattice pions.
If correct, the phase diagram of Fig. 1 leads to an important observation. For $`0<M<4`$, as we increase $`g`$ at fixed $`M`$, we eventually move into the massive phase, thus loosing all the massless quarks (even though the action is still invariant under Lüscher’s symmetry)! Remembering the connection with domain-wall fermions , this means that under the same conditions domain-wall fermions will support (no massless quarks and) no massless pions even if the extent of the fifth dimension is taken to be arbitrarily large.
The phase diagram of Wilson fermions is known to contain the Aoki phase , where parity and vector-like symmetries are broken by a flavor non-singlet pseudo-scalar condensate. The question may arise whether a similar phenomenon can take place in the present context as well. The answer is no. As explained in ref. , the Aoki phase is related to those terms in the effective chiral lagrangian (for finite lattice spacing) that break axial symmetries explicitly. There are no such terms in the present situation since one has a chiral symmetry at finite lattice spacing as well. Using the lattice chiral symmetries one can rotate any non-singlet pseudo-scalar condensate into the usual singlet scalar one $`\overline{q}q`$. An Aoki phase may, however, arise if an explicit quark mass term is added to the action. This was investigated numerically in the context of domain-wall fermions in ref. .
7. Discussion. Finally, we compare our results with previous work. Our results are in agreement with those of Ichinose and Nagao who studied the massless phase to second order in the hopping-parameter expansion . The fate of the $`U_A(1)`$ pseudo-scalar state was left open in ref. , and it was conjectured that it may eventually pick up a non-zero mass (in a higher order in $`1/N_c`$). As we showed rigorously in this paper, in fact $`U_A(1)`$ is an exact symmetry for a small hopping parameter, and the corresponding pseudo-scalar is an exact Goldstone boson for any value of $`g`$. This result is consistent with the observation that, for small $`\kappa `$, the lattice $`U_A(1)`$ symmetry becomes a flavor non-diagonal axial symmetry in the continuum limit.
The hamiltonian strong-coupling analysis of Brower and Svetitsky was done for domain-wall fermions with a continuous fifth coordinate ($`a_5=0`$) in the limit $`L_5\mathrm{}`$, where $`L_5`$ is the size of the fifth dimension. We recall that in this double limit domain-wall fermions are expected to reduce to overlap-Dirac fermions . This work is not limited to a small hopping parameter. At zeroth order in the expansion in $`1/g^2`$ they find a massive phase for $`M<3`$ and a gap-less phase for $`M>3`$. This leading-order result supports the mean-field argument of Sect. 6 that there is only one critical point at $`g=\mathrm{}`$. (At that order there are no spatial couplings, and for each site the problem reduces to a free one-dimensional hamiltonian acting on the fifth coordinate. While the critical value in ref. is $`M=3`$, rather than $`M=4`$, this is merely a technical difference stemming from the fact that time is taken to be continuous in ref. .)
The order-$`1/g^2`$ result of ref. is summarized, for $`M>3`$, by an effective low-energy hamiltonian in $`3+1`$ dimensions which contains two terms, and (in their notation) reads $`H_{\mathrm{s}\mathrm{s}}^{\mathrm{eff}}+H_{\mathrm{site}}^{\mathrm{eff}}`$. The symmetry of $`H_{\mathrm{s}\mathrm{s}}^{\mathrm{eff}}`$ is that of naive fermions. Since, at that order, $`H_{\mathrm{s}\mathrm{s}}^{\mathrm{eff}}`$ has only nearest-neighbor couplings, this can be explained by the fact that the $`O(\kappa )`$ term in the overlap-Dirac operator is proportional to the naive fermion one. At order $`1/g^4`$, $`H_{\mathrm{s}\mathrm{s}}^{\mathrm{eff}}`$ will contain also next-to-nearest neighbor couplings, and this may reduce the global symmetry to $`U_L(N_l)\times U_R(N_l)`$. If we would ignore the second term, $`H_{\mathrm{site}}^{\mathrm{eff}}`$, this would be in agreement with our results.
The second term in the effective hamiltonian, $`H_{\mathrm{site}}^{\mathrm{eff}}`$, explicitly breaks all axial symmetries, having the same structure as an explicit mass term for the quarks. Although mathematically there is no direct conflict between this (hamiltonian) result and our (euclidean) result, physically the two results seem to be in conflict if both $`g`$ and $`M`$ are large (and therefore both results should be valid). We hope to resolve this issue in the future. To this end, it may be useful to carry out a euclidean strong-coupling analysis, for example using the method of ref. .
In conclusion, in this paper we have analyzed overlap-Dirac fermions with a small hopping parameter. While there are issues that require further work, a concrete picture of the phase diagram is beginning to emerge. In the future we hope to generalize the discussion to domain-wall fermions (with a non-zero $`a_5`$ and a finite $`L_5`$), thus making closer contact with present-day numerical simulations. However, as we argued here, already in the limit where Lüscher’s chiral symmetry is exact, it is possible that a phase without massless quarks exist at fixed $`0<M<4`$ and large gauge coupling.
Acknowledgements. We wish to thank Mike Creutz, Ikuo Ichinose, Ben Svetitsky and Marvin Weinstein for helpful discussions. MG would like to thank the Institute for Nuclear Theory at the University of Washington for its hospitality. This research is supported by the United-States – Israel Binational Science Foundation, and Y.S. is supported in part by the Israel Science Foundation, while MG is supported in part by the US Department of Energy.
Appendix. Exponential localization. The Legendre expansion of $`(H^2)^{1/2}`$ (cf. eq. (3)) was used in ref. to prove the exponential localization of $`D`$. This expansion is very informative, and may be used to obtain exponential-localization bounds on functions of $`D`$ as well, in particular on $`D(2D)^1`$. The Legendre expansion is given by
$$(H^2)^{1/2}=c\underset{k0}{}t^k\varphi _k(Z).$$
(14)
Here $`\varphi _k`$ are Legendre polynomials, normalized such that $`\mathrm{max}\varphi _k(z)=1`$ for $`1z1`$. With slight adaptation the other ingredients are defined as follows. Let $`v_0=\mathrm{max}(H^2)`$, $`u_0=\mathrm{min}(H^2)`$, where extremization is done over the entire gauge-field configuration space. We take a (positive and) small enough hopping parameter $`\kappa `$ so that $`u_0>0`$. Then
$$Z=\frac{u_0+v_02H^2}{v_0u_0},$$
(15)
$$\frac{1}{2}(t+t^1)=\frac{v_0+u_0}{v_0u_0},$$
(16)
with $`0<t<1`$ and
$$c=\left(\frac{4t}{v_0u_0}\right)^{1/2}.$$
(17)
Since $`|B|4`$ and (for $`A=1`$)
$$H^2=1+2\kappa B+O(\kappa ^2),$$
(18)
we have $`v_0=1+8\kappa +O(\kappa ^2)`$, $`u_0=18\kappa +O(\kappa ^2)`$, $`t=4\kappa +O(\kappa ^2)`$ and $`c=1+O(\kappa )`$. We may then write (for $`A=1`$)
$$D=(c1)+c\underset{k1}{}t^k\left(\varphi _k(Z)bY\varphi _{k1}(Z)\right),$$
(19)
where $`b=\kappa /t`$ and $`Y=_\mu \gamma _\mu C_\mu B`$, cf. eq. (3). We comment that the Legendre expansion may be set up using any $`0<uu_0`$ and $`vv_0`$. The “best” values defined above, $`u_0`$ and $`v_0`$, guarantee that $`t/\kappa =O(1)`$. This is a natural relation because, as described below, $`t`$ effectively plays the role of a hopping parameter.
We now turn to the operator
$$D(2D)^1=\underset{n1}{}(D/2)^n.$$
(20)
In order to obtain a bound on the kernel corresponding to eq. (20) it is useful to regard $`t`$ as an independent parameter. Doing so, we obtain the expansion
$$D(2D)^1=\underset{k0}{}t^k𝒟_k,$$
(21)
where $`𝒟_k`$ is defined by substituting eq. (19) into eq. (20) and collecting all terms involving an explicit factor of $`t^k`$. At order $`t^k`$, the operator encountered in the expansion of $`D`$ (eq. (19)) allows for at most $`2k`$ hoppings . The same statement applies to the new kernels: $`𝒟_k(x,y)=0`$ for $`|xy|>2k`$ where $`|xy|=_\mu |x_\mu y_\mu |`$ (the “taxi-driver distance”). Therefore
$$D(2D)^1(x,y)=\underset{2k|xy|}{}t^k𝒟_k(x,y).$$
(22)
What is still needed is a bound on $`𝒟_k`$. We first observe that
$$|\varphi _k(Z)bY\varphi _{k1}(Z)|1+8b.$$
(23)
In view of this bound we are lead to consider
$$E=d+c(1+8b)\underset{k1}{}t^k,$$
(24)
where $`d=|c1|`$, as well as
$$E(2E)^1=\underset{n1}{}(E/2)^n=\underset{k0}{}t^k_k.$$
(25)
Again $`_k`$ is defined by collecting all terms involving $`t^k`$ using eq. (24). The double geometric series is easily summed, giving
$$E(2E)^1=\frac{d+(c+8bcd)t}{2d(2+c+8bcd)t},$$
(26)
and an explicit expression for $`_k`$ follows by re-expanding eq. (26) in powers of $`t`$.
Following the construction we see that, thanks to inequality (23), each term in $`𝒟_k`$ is bounded by a corresponding term in $`_k`$. The latter is obtained if every factor of $`\varphi _k(Z)bY\varphi _{k1}(Z)`$ is replaced by $`1+8b`$ and every factor of $`c1`$ by $`d`$. Using eq. (22) we thus obtain a bound (for $`|xy|2`$)
$`|D(2D)^1(x,y)|`$ $``$ $`{\displaystyle \underset{2k|xy|}{}}t^k_k=`$ (27)
$`=`$ $`{\displaystyle \frac{ds^{[|xy|/2]}+t(c+8bcd)s^{[|xy|/21]}}{2d(2+c+8bcd)t}},`$
where $`s=(2+c+8bcd)t/(2d)`$ and $`[l]`$ is the smallest integer $`n`$ such that $`nl`$. Recalling the value of $`t=4\kappa +O(\kappa ^2)`$ (cf. eq. (16)), the desired exponential bound on the kernel of $`D(2D)^1`$ is thus given by eq. (27) provided we choose $`\kappa <\kappa _0`$, where $`\kappa _0`$ is the smallest value where the denominator in eq. (27) is zero. The denominator is strictly positive for $`0<\kappa <\kappa _0`$ since both $`t`$ and $`d`$ are $`O(\kappa )`$.
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# Phase Clocks for Transient Fault Repair
## 1 Introduction
Measuring time is widely recognized as an important system service and greatly simplifies the construction of many distributed algorithms. The reason, simply put, is that deductions about the progress of concurrent activities, made by measuring elapsed time, effectively substitute for communication and protocols that directly monitor such progress. Of course this technique can only be used to the extent that a distributed system is synchronous, matching its progress with elapsed time. Yet so attractive is the use of time to simplify algorithm construction, that even in asynchronous systems, researchers seek to simulate synchrony , introduce logical clocks , and/or logical time as programming tools.
One illustration of logical time in an asynchronous system is the organization of a computation into phases. The basic property of a phased computation is that a process does not enter phase $`(k+1)`$ until each related process has completed phase $`k`$. The case where all processes are related is equivalent to barrier synchronization, and the case where the relation between processes is specified by a graph corresponds to a *phase clock*. Many implementations of phased computation simply use a counter, called a *clock*, to represent the current phase number of a process. Consider the graph relation between processes to be a network communication topology, where the graph has diameter $`𝒟`$ and distance between processes $`p`$ and $`q`$ is denoted by $`\text{dist}_{pq}`$. A phase clock invariantly relates phase numbers and distance as follows: any process $`p`$ has $`\text{clock}_p=k+d`$ only if $`\text{clock}_qk`$ holds for each process $`q`$ satisfying $`\text{dist}_{pq}=d`$ (notice that $`k=1`$ is just the basic property mentioned above). Thus if $`\text{clock}_p=k+d`$, holds at some state, we deduce that $`\text{clock}_q=k`$ holds currently or held at some previous state. This is a useful timing property because programs can use phase clocks for inferences about nonlocal information relayed through neighboring processes. For example, process $`p`$ could use its clock to infer termination of a broadcast operation, rather than use explicit termination detection, by waiting for sufficiently many increments to $`\text{clock}_p`$ (assuming that the broadcast operation is geared to the phase clock).
Phased computation is a reasonable discipline for many activities of a distributed system, including procedures invoked as part of fault diagnosis and repair. The fault domain for this paper is the model of *transient faults*, which corrupt local process states and communication registers, but do not damage a system’s control logic. It is therefore feasible for a system to self-diagnose and restore variables corrupted by a transient fault to values that enable correct system function. One of the difficulties in using phase clocks to control distributed repair activities is that faults may corrupt local clock values. The phase clock protocol presented in this paper not only repairs clock values corrupted by a transient fault, but does so in a manner that enables the system to use the phase clock for other repair activities.
#### Contributions.
This paper presents a distributed phase clock, called the *repair timer*, specialized for the task of transient fault repair in a distributed system. The repair timer is *time adaptive*, meaning that it satisfies desired accuracy and progress properties within $`O(k)`$ time after any transient fault event corrupting $`k`$ processes. The repair timer differs from standard phase clocks because it starts at zero and halts after repair is complete (behaving somewhat like an egg-timer); this enables direct inspection of elapsed repair time, which standard phase clocks do not provide<sup>1</sup><sup>1</sup>1To see why this is not trivial, suppose some time-adaptive phase clock were available, and consider measuring elapsed repair time by recording the start time of repair in some local variable; such a local variable could, however, have an erroneous value due to a transient fault. Since transient faults do not provide any signal at the start of repair, a process cannot locally decide whether its local variables are accurate or not.. The repair timer is also a self-stabilizing algorithm, able to restore all variables to a legitimate state following any transient fault event or combination of transient failures. Finally, the paper presents composition theorems to show how the repair timer is useful for the timing of fault repair procedures in a distributed system.
#### Related work.
Many recent works are motivated by what is seen as pessimism in the model of self-stabilization, which does not discriminate between cases of severe transient faults and minor transient faults. In addition to the desired robustness of self-stabilization, fast stabilization of output variables has been recently demonstrated in a number of algorithms and some general methods to achieve time adaptivity or local self-stabilization .
Self-stabilizing phase clocks are given in . None of these constructions guarantee fast stabilization for cases of limited transient faults, and all appear to require lengthy stabilization time (proportional to the diameter of the communication graph) in some cases where only a single process variable is corrupted by a transient fault. Requirements for a repair timer are described in , which is a precursor to this paper.
#### Contents.
Section 2 presents the computation and system model for the paper. Section 3 presents the algorithm for the repair timer and Section 4 verifies the self-stabilization and time adaptive properties of the algorithm. To illustrate the use of the repair timer, Section 5 describes two designs incorporating the repair timer as a component in a system. The paper’s concluding remarks are the subject of Section 6. Proofs of technical lemmas have been moved to the paper’s Appendix.
## 2 Distributed System
The system consists of a fixed set of $`n`$ processes that communicate by reading and writing shared registers. Communication between processes is limited to a network represented by an undirected, connected graph: for any pair of processes $`(p,q)`$ there exist a pair $`(\text{Register}_{pq},\text{Register}_{qp})`$ if and only if there is an edge between $`p`$ and $`q`$ in the communication graph. Process $`p`$ is the only writer of $`\text{Register}_{pq}`$ and $`q`$ is the only reader of $`\text{Register}_{pq}`$. A process cannot read the registers it writes. Registers thus approximate message passing with bounded buffers, and a self-stabilizing simulation of link registers using messages is described in . A register can have numerous fields used to write values of different local variables (just as numerous local values can be transmitted in fields of one message).
If $`(p,q)`$ is an edge in the communication graph, then $`p`$ and $`q`$ are called *neighbors*, which is denoted by $`p𝒩_q`$ or equivalently, $`q𝒩_p`$. The diameter of the communication graph is $`𝒟`$. The distance between any pair $`(p,q)`$ in the graph is denoted by $`\text{dist}_{pq}`$. The term *region* refers to a connected component of the graph that has some property of interest.
Each process is an autonomous, finite-state computing entity. We use conventional imperative programming notation and concepts to describe the operation of a process, so each process has a program counter and program statements that manipulate variables. A subset of these variables are called *output variables*, which directly support the system’s intended function.
A *configuration of p* is a specification of values, one for each of process $`p`$’s variables, the value of $`p`$’s program counter, and a value for each register that $`p`$ writes. A (system) *state* is a vector of process configurations, one configuration for each process in the system. Any function from the set of all states to the set $`\{\text{true},\text{false}\}`$ is called a *state predicate*.
A *process step* is either a register operation (and corresponding advancement of the program counter) or some modification of internal and output variables (and program counter) of that process. A *computation* is an infinite sequence of states so that each consecutive pair of states corresponds to a process step and the sequence of states includes an infinite number of steps of each process. We thus assume that computations are fair; more precisely, we assume weak fairness in that no process is prevented from executing steps in a computation. We use the term *computation segment* to denote a finite, contiguous subsequence of a computation.
The program of each process specifies a *cycle*, which consists of three parts: (i) a process reads the registers written by each of its neighbors, (ii) the process possibly assigns values to its variables, and (iii) the process writes registers for each of its neighbors. The definition of a cycle is a convenient and simple abstraction for measuring the progress of a process in a computation.
The system is designed to accomplish some task represented by a state predicate $`_O`$. Whether or not $`_O`$ holds at a given state is solely determined by the values of output variables. A predicate $``$ is called a *legitimacy predicate* iff $``$ is a system invariant and $`_O`$. A state $`\sigma `$ is *output-legitimate* if $`_O`$ holds at $`\sigma `$, and is *legitimate* if $``$ holds at $`\sigma `$. It is often preferable to specify legitimacy (or output legitimacy) in terms of the behavior of processes rather than explicitly specifying a state predicate. A formal definition of legitimacy in terms of behaviors is possible, but to streamline the presentation, the state-based definition is used in this paper. Where process behavior is important in this paper, we verify separately that the system exhibits the desired behavior.
Because each iteration of a process program specifies a cycle, time is conveniently measured in asynchronous rounds, which are defined inductively. A *round* of a computation, with respect to an initial state $`\sigma `$, is a computation segment originating with $`\sigma `$ of minimum length containing at least one complete cycle (from reading registers to writing registers) of each process. The *first round* of a computation consists of a round with respect to the initial state of the computation, and *round k* of a computation, $`k>1`$, consists of a round with respect to the first state following round $`k1`$.
A round is, roughly speaking, one unit of “parallel time” in the system. A notion similar to a round is commonly used to analyze the complexity of message-passing protocols by normalizing message delay to the maximum message delay . For analysis in this paper, the notion of a round is further refined. An $`\text{R}_p^d`$-round starting from a state $`\sigma `$ is a computation segment of minimal length containing at least one complete cycle of each process in the set $`\{q|\text{dist}_{pq}d\}`$. A round is thus equivalent to an $`\text{R}_p^𝒟`$-round for any choice of $`p`$.
A system is *self-stabilizing* if every computation contains a legitimate state (that is, for any initial state, the system eventually reaches a legitimate state). The *stabilization time* is the worst case number of rounds in the prefix of a computation that does not contain a legitimate state. Proving that a system is self-stabilizing entails demonstrating that a predicate $``$ is invariant, implies $`_O`$, and that every computation contains some state satisfying $``$.
A fault event is a non-computational operation that modifies variables, program counters, and/or registers. More formally, a *fault event* can be any pair of states (whereas a consecutive pair of states in a computation is a process step). Computations do not include fault events; a system history could be a sequence of states consisting of computation segments punctuated by fault events. Reasoning about fault repair proceeds with respect to each computation segment, since the system cannot anticipate whether or not another fault will occur.
A state $`\sigma `$ is *k-faulty* if $`k`$ is the minimum number of process configurations in $`\sigma `$ that, if appropriately changed, transform $`\sigma `$ into a legitimate state. The number $`k`$ thus corresponds to the Hamming distance from $`\sigma `$ to the nearest legitimate state. There may be numerous ways to transform $`\sigma `$ to a legitimate state by changing $`k`$ process configurations, some in which process $`p`$’s configuration changes, and others where the transformation does not change $`p`$’s configuration. It is convenient to resolve this ambiguity by some unique, deterministic choice of which processes should change configurations to obtain a legitimate state from $`k`$-faulty state $`\sigma `$. With such a deterministic choice, process configurations of $`\sigma `$ can be labelled *faulty* or *nonfaulty* depending on whether they should change or not. This deterministic choice can further be refined to label variables and register fields as either faulty or nonfaulty. How such a deterministic choice should be implemented turns out not to be an issue in the sequel; for the repair timer given in Section 3 there is an unambiguous definition of a faulty process configuration and for the interface proposed in Section 5 it is only required that if a faulty process configuration neighbors a nonfaulty process configuration, then the presence of a fault can be detected (which for many systems is the case even by reversing the choice of which of these two neighboring configurations is considered to be faulty).
The main emphasis of this paper is time-adaptive, stable repair of output variables, meaning that a system should stabilize its output variables to satisfy $`_O`$ from any $`k`$-faulty initial state after at most $`O(k)`$ rounds. Formally, a system is *time adaptive* if each computation starting from any $`k`$-faulty initial state $`\sigma `$ contains an output-legitimate state $`\sigma ^{}`$, within $`O(k)`$ rounds following $`\sigma `$, such that every state following $`\sigma ^{}`$ in the computation is output-legitimate. Given this emphasis, it is convenient to extend the terminology for faults: a process $`p`$ is faulty (nonfaulty) in a computation iff $`p`$’s configuration is faulty (nonfaulty) at the initial state.
## 3 Algorithm
One of the difficulties in using phase clocks to control distributed repair activities is that faults may corrupt local clock values. Indeed, repair of the clocks values is a primary concern of this paper, and the usual timing properties of phase clocks must be modified to cope with faults. Two goals for such modifications are: (a) clock values of processes not affected by faults can reliably be used for inferences about nonlocal information; (b) the response effort of the system is proportional to the scope of the fault.
Goal (a) seems relatively simple to satisfy, since the clocks of nonfaulty processes have predictable values. However for a standard phase clock, there are ambiguous cases of faulty situations. Suppose neighboring clocks have values $`x`$ and $`x2`$ and only one of these two is a faulty value; there is no obvious way of distinguishing which of these two is faulty. The approach taken in this paper is to use a specialized phase clock for fault repair called a *timer*. Whereas phase clocks advance throughout system computation, the timer stops advancing when repair is complete. Thus each timer clock reaches a prescribed value $`𝒯`$ when the system state is fully repaired. If neighboring clocks have values $`𝒯`$ and $`𝒯2`$, then we may conclude that the value $`𝒯2`$ is due to a fault.
#### Variable Conventions.
The variables appearing in Figure 1 are local variables of process $`p`$. A number of proof arguments are statements relating variables of different processes, and subscripts are used to distinguish variable ownership (for instance, $`\text{clock}_q`$ is owned by $`q`$). Similarly, the predicates defined in Figure 1 are subscripted in definitions and proof arguments (such as $`\text{gap}_p`$ for process $`p`$).
Statement S1 copies four register fields to four local variables, $`\text{x,y,r,s}`$. Call these variables the *image* variables. Implicitly the code of Figure 1 defines a mapping from each image variable to a register field and a corresponding “base” variable of a neighboring process (written by statement S7). We say that each image variable is *based on* a variable of a neighboring process, meaning that the value of an image variable is copied (via register communication) from the variable upon which it is based. Variable $`\text{x}_p[q]`$, for example, is based on $`\text{clock}_q`$. Register fields are also images that are statically based on variables.
The meaning of time adaptivity described in Section 2 depends on declaring some of the process variables to be output variables. For the repair timer, let $`\text{clock}_p`$ be the output variable of process $`p`$. The output correctness for clock variables is the subject of Section 4.2.
#### Program Conventions.
The statements S1S7 given in Figure 1 describe one complete cycle of the repair timer for process $`p`$. Therefore, after executing S7, process $`p`$ executes S1 to start the next cycle. The group of statements S3S6 constitute a multiway if statement; in any cycle, at most one of S3S6 are executed.
Statements S2S6 specify internal calculations for process $`p`$, since they manipulate local variables. In a computation, we suppose that each of these statements specifies one computation step. Statements S1 and S7 specify each $`|𝒩_p|`$ computation steps, since a step can read or write at most one register. Ordering of the read and write operations of S1 and S7 is unimportant to the algorithm.
#### Algorithm Structure.
To understand the algorithm of Figure 1 it is useful to first ignore statements S3S6 and focus attention on the w variables. Notice that statement S2 will reduce $`\text{w}_p`$ if any of the registers read by S1 imply a value $`\text{w}_p2`$ or smaller for any neighboring w variable. The global effect of many processes executing S2 can thus be “convergence to the least w” over a number of rounds. The result of executions of S2 will, in general, lead to a situation where neighboring process w variables differ by at most 1, which is one of the properties of a phase clock. The wEcho condition of S2 allows any process with a globally minimal w variable to increment its w variable after all neighbors acknowledge its current value, via the s image variables (which occurs within two rounds). Therefore the set of w variables apparently enjoy both properties of a phase clock — that neighboring w variables differ by at most 1 and increase continually (until the upper bound of $`3𝒟+1`$ occurs) in a computation.
Why not simply use the w variables for repair timing and dispense with the logic of S3S6? The answer lies in the additional constraint we impose for faulty initial states. For repair purposes, it is not enough for clocks to be in phase and increment, they should also be accurate, meaning that the value of a clock should be a measure of how long computation has progressed after the initial detection of a fault. The w variables do not have this property. For instance, a faulty initial state could have $`𝒟`$ as the initially smallest w variable, so that all subsequent states have w variables overstating the repair time at least by $`𝒟`$. An attempt to fix this problem would be some statement similar to S4, that would reset w to zero whenever neighboring w variables differ by more than 1. It is easy to construct examples of computations where such an attempted fix will fail because w variables are reset to zero infinitely often. This kind of idea can work, however, if any w variable were guaranteed to be reset to zero at most once in a computation, and that is the basic idea behind statements S3S6, which reset a clock variable to zero at most once in a computation. Although w variables do not enjoy the accuracy needed for repair timing, they provide a useful “reset layer” for the clock adjustments of the algorithm.
###### Definition 1
*A state is *timer-final* if every register field and image variable value is equal to the value of its base variable, and $`(p::\text{clock}_p=𝒯\text{w}_p=3𝒟+1)`$. A process configuration is *timer-final* if all its variables and register fields have values corresponding to a timer-final state. We define predicate $`_T`$ to hold for a state iff that state is timer-final.*
The value $`𝒯`$ used in the algorithm and Definition 1 is a constant adequate for the fault tolerance of the repair timer and for the application of the timer, as discussed in Section 5. The proof of self-stabilization of the repair timer requires only that $`𝒯11𝒟`$.
#### Verification.
The verification of desired repair timer properties is divided into two stages. First, the algorithm is considered as an isolated component, so that faulty states are those states deviating from Definition 1. Section 4.1 is devoted to a proof that the repair timer self-stabilizes to a timer-final state. Section 4.2 presents the proofs that apply to $`k`$-faulty initial states, showing that the repair timer achieves desired accuracy after $`O(k)`$ time following the initial state.
The second stage of verification is concerned with integration of the repair timer as a component of a system. The timer is a tool for time adaptive repair. Discussion of how the timer is used is deferred to Section 5, where it is explained that the timer is a service with only one operation, namely to start the timer by assigning $`\text{clock}0`$; thereafter, the clock should increment as a phase clock. Although a system state’s legitimacy depends on variables of all system components, the simple interface between the timer and other system components makes it reasonable to consider fault tolerance properties of the timer in isolation, which motivates the two stage approach to verification.
## 4 Stabilization and Adaptivity
### 4.1 Self-Stabilization
Each process writes its communication registers in every cycle from its variables. Therefore, following the first round of any computation, all register fields are equal to current or previous values of the corresponding base variables. Following the second round, each image variable has a value previously written from the corresponding base variable. Moreover, following the third round of a computation, the third and fourth fields of $`\text{Register}_{qp}`$ contain values previously written by $`p`$ and then copied by $`q`$. It is convenient to assume that register fields correspond to values previously written in the computation, so we call a computation *based* if it is the suffix, starting from round three or higher, of another computation.
Statements S4 and S5 have the only assignments that may reduce the value of clock variables. We call a computation (or computation segment) *reset-free* if no process executes S4 or S5 in that computation. A computation is called *rising* if it is the suffix of a based, reset-free computation such that each process has read its registers at least once in the based, reset-free computation prior to the first state of the suffix. Rising computations enjoy the useful property that at all states, the value contained in $`\text{x}_p[q]`$ is a lower bound on the current value of $`\text{clock}_q`$. (This property follows because the computation is reset-free and each process previously read registers and assigned to its x variables while the computation was reset-free.)
###### Definition 2
$$b_{pq}(p𝒩_q)(|\text{clock}_p\text{clock}_q|<2\text{x}_p[q]\text{clock}_q|\text{clock}_p\text{x}_p[q]|<2)$$
*A state is *smooth* if $`(p,q::b_{pq})`$. A set of processes $`P`$ forms a *smooth region* if the subgraph of the communication topology induced by $`P`$ is connected and $`(p,q:p,qP:b_{pq})`$.*
###### Lemma 1
*In a rising computation, $`(b_{pq}b_{qp})`$ is an invariant for any pair of processes $`p`$ and $`q`$.*
###### Lemma 2
*Let $`\sigma `$ be the first state of a rising computation segment such that for $`p𝒩_q`$, both $`\text{clock}_p`$ and $`\text{clock}_q`$ have incremented at least once in the computation segment. Then $`(b_{pq}b_{qp})`$ holds at state $`\sigma `$.*
###### Lemma 3
*If each clock variable has incremented at least once prior to state $`\sigma `$ in a rising computation segment, then $`\sigma `$ is smooth.*
###### Lemma 4
*Smoothness is invariant for a based computation; within $`O(𝒟)`$ rounds following a smooth state, a based computation contains a timer-final state.*
###### Lemma 5
*If $`\text{clock}_p`$ is less than $`𝒯`$ and less than or equal to all neighboring clock values at the initial state of a based, reset-free computation segment, and $`\text{w}_p=3𝒟+1`$ holds at the initial state, and this computation segment contains at least two rounds, then $`\text{clock}_p`$ increments at least once in the computation segment.*
###### Lemma 6
*Let the initial state of a based computation satisfy $`(p::\text{clock}_p7𝒟\text{w}_p=3𝒟+1)`$. The computation contains a state where $`(q::\text{clock}_q=10𝒟+1)`$; the first state satisfying $`(q::\text{clock}_q=10𝒟+1)`$ is a smooth state.*
###### Lemma 7
*Let the initial state of a based computation satisfy $`(p::\text{clock}_p7𝒟\text{w}_p=3𝒟+1)`$. Within $`O(𝒟)`$ rounds, the computation contains a smooth state.*
###### Lemma 8
*Consider a based computation such that $`(\text{clock}_r=0\text{w}_r=0)`$ holds for some process $`r`$ in the initial state. Within $`𝒟`$ rounds there is a state satisfying $`(p::\text{clock}_p3𝒟\text{w}_p3𝒟)`$.*
###### Lemma 9
*Consider a based computation such that $`(\text{clock}_r=0\text{w}_r=0)`$ holds for some process $`r`$ in the initial state. Within $`O(𝒟)`$ rounds there is a smooth state or there is a state satisfying $`(p::\text{clock}_p7𝒟\text{w}_p=3𝒟+1)`$.*
###### Theorem 1
*The timer stabilizes to a timer-final state (satisfying $`_T`$) in $`O(𝒟`$) rounds.*
* The invariance of $`_T`$ is verified by observing that none of S1S6 change any variable value at a timer-final state. Convergence is demonstrated by a sequence of claims about an arbitrary computation $`A`$. Let $`B`$ be a suffix of $`A`$ beginning following the second round of $`A`$; by definition, $`B`$ is a based computation. We consider two cases for $`B`$.
Case: $`B`$ contains no step executing S4 within $`O(𝒟)`$ rounds. By arguments similar to those given in the proof of Lemma 9, some state of $`B`$ satisfies $`(p::w_p=3𝒟+1)`$ within $`O(𝒟)`$ rounds and continues to hold at least until S4 executes. Let $`C`$ be a suffix of $`B`$ satisfying $`(p::w_p=3𝒟+1)`$ at its initial state. Observe that $`C`$ is based and reset-free for $`O(𝒟)`$ rounds, so Lemma 5 is applicable to $`C`$. Within $`O(2𝒯)`$ rounds of $`C`$, $`(p::\text{clock}_p=𝒯)`$ holds, and the state satisfies $`_T`$.
Case: $`B`$ contains some step executing S4 within $`O(𝒟)`$ rounds. Execution of S4 results in a state satisfying the premise of Lemma 9. Therefore $`B`$ either contains a smooth state within $`O(𝒟)`$ rounds, or contains a state satisfying $`(p::\text{clock}_p7𝒟\text{w}_p=3𝒟+1)`$ within $`O(𝒟)`$ rounds. The latter possibility is the premise for Lemma 7, which shows that a smooth state is subsequently obtained within an additional $`O(𝒟)`$ rounds, so with either possibility, $`B`$ contains a smooth state within $`O(𝒟)`$ rounds. Lemma 3 implies that $`B`$ contains a timer-final state within $`O(𝒟)`$ rounds following a smooth state.
### 4.2 Time Adaptivity
The desired fault tolerance of the timer consists, informally, of the following two properties. (1) Within $`k`$ rounds from a $`k`$-faulty initial state, every clock is accurate, that is, if $`\text{clock}_p=t`$ for $`t<𝒯`$, it should be that $`p`$ has incremented $`\text{clock}_p`$ as a phase clock $`t`$ times during the repair procedure. (2) Each faulty process clock is reset to zero and subsequently increments as a phase clock, incrementing to $`k`$ within $`O(k)`$ rounds.
Property (1) provides the accuracy needed so that a process can safely wait for distant information to be reliable. Property (2) assures that such distant information arrives in a timely fashion. Because faults may damage clock and other timer variables, Theorem 2 below provides a conditional form of (1), necessarily relaxed to accommodate unusual initial states. Also, some unusual cases of initial states require a conditional form of (2), provided by Theorem 3.
A system state is faulty if it does not satisfy the definition of legitimacy. In considering the timer in isolation, a state is $`k`$-faulty if no fewer than $`k`$ process configurations require change to obtain a timer-final state. However, a complete definition of system legitimacy depends on components other than the timer, so a limited notion of fault is appropriate for the timer.
###### Definition 3
*A set of processes $`P`$ is *unperturbed* at state $`\sigma `$ if $`P`$ forms a smooth region, $`(p:pP:\text{clock}_p>𝒯𝒟)`$, and $`(p,q:pPq𝒩_pqP:\text{clock}_p=𝒯\text{x}_p[q]𝒯1)`$. A process $`p`$ is *unperturbed* at $`\sigma `$ if there exists an unperturbed region containing $`p`$; process $`p`$ is *perturbed* if there exists no unperturbed region containing $`p`$. State $`\sigma `$ is $`k`$-*perturbed* iff $`k`$ is the number of perturbed processes at $`\sigma `$.*
The motivation for this definition derives from the ambiguity of certain clock values and nondeterminism of asynchronous computation. Some proofs are simplified using Definition 3, which defines a perturbed process to be a weakening of a faulty process configuration (a nonfaulty process configuration is unperturbed, but the converse may not hold). It follows that if the timer algorithm satisfies desired properties (1)–(2) within $`k`$ rounds from any $`k`$-perturbed state, then similar properties also hold for any $`k`$-faulty initial state. Definition 3 is not useful if $`k=0`$, so in the sequel any reference to $`k`$-perturbed state is assumed to imply $`k>0`$.
###### Definition 4
*Within a computation, a variable $`\text{clock}_p`$ is *d-accurate* at a state $`\sigma `$ if $`\text{clock}_p>𝒯𝒟`$ holds, or if $`\text{clock}_p𝒯𝒟`$ implies, for $`0m𝒟`$, that the number of $`\text{R}_p^m`$-rounds completed prior to state $`\sigma `$ is at least $`(\text{clock}_pmd)`$, and that for every process $`q`$, the value of $`\text{clock}_q`$ has incremented at least $`(\text{clock}_p\text{dist}_{pq}d)`$ times prior to state $`\sigma `$ in the computation. For a computation initiating from a $`k`$-perturbed state, a state $`\sigma `$ is *time-accurate* if for unperturbed $`p`$, $`\text{clock}_p`$ is $`d`$-accurate for $`d=2\mathrm{min}(k,𝒟)`$, and for perturbed $`p`$, $`\text{clock}_p`$ is $`d`$-accurate for $`d=5\mathrm{min}(k,𝒟)`$.*
Definition 4 falls short of the desired precision of property (1), but satisfies safety concerns for many situation of repair timing because a $`d`$-accurate clock provides a lower bound on the number of cycles that distant processes have completed during repair. For instance, a repair application could depend on a distributed procedure that terminates after $`m`$ clock increments in a non-faulty environment; this application could wait for $`d+m`$ clock increments if the repair timer ensures only $`d`$-accurate clock variables. Unfortunately, an initially faulty state can have arbitrary values in faulty process clock variables, making it impossible to instantly have time accuracy. Theorem 2 given at the end of this section states that time accuracy is guaranteed from any $`k`$-faulty initial state, provided $`k<n`$, after at most $`\mathrm{min}(k,𝒟)`$ rounds of computation.
###### Lemma 10
*Any process $`p`$ executes S4, resetting $`\text{clock}_p`$ and $`\text{w}_p`$, at most once in any computation.*
Lemma 10 is a corollary of arguments given in the proofs of Lemmas 8 and 9. It is useful to know that processes execute S4 at most once because any reset step subsequent to S4 is therefore due to S5. Arguments in the proof of Lemma 9 show that w values increase if S4 does not execute, and this idea can be used to establish the eventual increase of clock values.
###### Lemma 11
*Let $`\sigma `$ be a result of $`p`$ executing S4. Then for any process $`q`$ satisfying $`\text{dist}_{pq}=t\mathrm{min}(k,𝒟)`$ there occurs a state $`\sigma ^{}`$, within $`t`$ rounds following $`\sigma `$, such that $`\text{clock}_q3t\text{w}_q3t`$; and if there is a path consisting of unperturbed processes from $`p`$ to $`q`$, then a state $`\sigma ^{\prime \prime }`$ occurs within $`t`$ rounds following $`\sigma `$ such that $`\text{clock}_qt\text{w}_qt`$.*
Lemma 11 considers a level of detail not discussed in the proof of Lemma 8, which supposes based computations. Lemma 11 can also be extended to distances beyond $`k`$, shown in the following.
###### Lemma 12
*In any computation starting from a $`k`$-perturbed initial state, for each unperturbed process $`p`$ satisfying $`\text{dist}_{pq}=t`$ with respect to some perturbed process $`q`$, the following holds: process $`p`$ executes S4 within $`4+\mathrm{min}(𝒟,k+t)`$ rounds.*
###### Lemma 13
*In any computation beginning from a $`k`$-perturbed state, any process $`p`$ satisfying $`\text{dist}_{pq}=t`$ from some perturbed process $`q`$ does not execute S4 after round $`4+\mathrm{min}(2𝒟,t+2k)`$.*
###### Lemma 14
*Let $`\sigma `$ be a result of $`p`$ executing S4. Then for any process $`q`$, within $`t`$ rounds following $`\sigma `$ there occurs a state $`\sigma ^{}`$ such that $`\text{w}_q\mathrm{min}((t\text{dist}_{pq})/2,3𝒟+1)`$ is invariant for the computation beginning with $`\sigma ^{}`$.*
###### Lemma 15
*Let $`\sigma `$ be a result of $`p`$ executing S4. Then for any process $`q`$, within $`t+2`$ rounds following $`\sigma `$ there occurs a state satisfying $`\text{clock}_q\mathrm{min}(((t2)\text{dist}_{pq})/2,𝒯)`$.*
###### Theorem 2
*Any computation starting from a $`k`$-perturbed initial state, $`k<n`$, contains a time-accurate state $`\sigma `$ after at most $`\mathrm{min}(k,𝒟)`$ rounds following the initial state, and all states following $`\sigma `$ are time-accurate states.*
* Provided $`k<n`$, arguments in the proof of Lemma 12 show that for each perturbed region $`R`$, some process $`r`$ executes S4 within the first round, where $`r`$ satisfies either $`rR`$ or $`r𝒩_q`$ for some $`qR`$. Lemma 11 then implies that within $`\mathrm{min}(k,𝒟)`$ additional rounds, each $`pR`$ satisfies $`\text{clock}_p3\mathrm{min}(k,𝒟)`$. Each unperturbed process clock variable remains larger than $`𝒯𝒟`$ until S4 is executed, which resets the clock to zero. Thus within $`\mathrm{min}(k,𝒟)`$ rounds, each $`\text{clock}_p`$ is either larger than $`𝒯𝒟`$ or is at most $`3\mathrm{min}(k,𝒟)`$. After $`\mathrm{min}(k,𝒟)`$ rounds, unperturbed processes can decrease clock variables to zero, but such a decrease does not falsify the conditions for a time-accurate state. Therefore, to show time accuracy, it suffices to show that increments to $`\text{clock}_p`$ imply corresponding increments have executed at distant processes.
After a process $`p`$ executes S4, it does not increment $`\text{clock}_p`$ until cEcho holds. If an unperturbed $`q`$ is a neighbor of $`p`$, then $`p`$ does not increment $`\text{clock}_p`$ until $`q`$ has reset $`\text{clock}_q`$ and updated the image variables and register fields so that $`p`$ observes cEcho. It is a simple induction to show that $`\text{clock}_p`$ cannot increase to a value $`t`$ unless $`q`$ has incremented $`\text{clock}_q`$ at least $`(t1)`$ times. Now consider a minimum length path $`P`$ of processes, of length $`d`$, from $`p`$ to some process $`r`$, such that each process in $`P`$ is unperturbed. By a double induction, on $`t`$ and $`d`$, it follows that $`\text{clock}_p`$ cannot increase from zero to $`t`$ unless each process $`qP`$ has incremented $`\text{clock}_q`$ at least $`t\text{dist}_{pq}`$ times. The same argument shows that processes of $`P`$ complete at least the same number of $`\text{R}_p^d`$-rounds in the period where $`\text{clock}_p`$ increases from zero to $`t`$.
Returning to the event of $`p`$ executing S4, we now consider the case of perturbed $`q𝒩_p`$. As observed in the proof of Lemma 11, it is possible that $`p`$ can increment $`\text{clock}_p`$ twice before $`q`$ completes a cycle because corrupt values in the initial state enable the cEcho and wEcho conditions. Furthermore, $`p`$ can increment $`\text{clock}_p`$ a third time before $`q`$ increments its clock because $`q`$ completes a cycle to enable $`\text{cEcho}_p`$. However in the case of such a third successive increment by $`p`$, $`\text{clock}_p>\text{clock}_q`$ and $`b_{pq}b_{qp}`$ hold as a consequence. Thereafter, we reason about the interaction between $`p`$ and $`q`$ as for unperturbed neighbors (note that any subsequent executions of S5 by $`p`$ or $`q`$ validate this argument, since we reason about the highest value attained for clock variables after $`p`$’s initial three increments). Therefore, the value of $`\text{clock}_p`$ does not increase to $`t`$ unless $`q`$ has incremented $`\text{clock}_q`$ at least $`t3`$ times. Again, we may consider a minimum length path $`P`$ of processes, of length $`d`$, from $`p`$ to some process $`r`$, such that each process in $`P`$ is perturbed (with the possible exception of $`p`$). By a double induction, on $`t`$ and $`d`$, it follows that $`\text{clock}_p`$ cannot increase from zero to $`t`$ unless each process $`qP`$ has incremented $`\text{clock}_q`$ at least $`t3\text{dist}_{pq}`$ times. Similar arguments show the completion of the appropriate number of $`\text{R}_p^d`$-rounds while $`\text{clock}_p`$ increases from zero to $`t`$.
Notice that in the case of a perturbed path of processes, accuracy can diminish by two extra clock units per unit of distance, whereas in the case of an unperturbed path, accuracy corresponds precisely to distance. These observations combined can be used to verify that in any minimum length path $`P`$ from $`p`$ to $`r`$, after $`p`$ executes S4, the value of $`\text{clock}_p`$ increases to $`t`$ only if for each $`qP`$, the value of $`\text{clock}_q`$ has incremented at least $`t\text{dist}_{pq}2m`$ times, where $`m`$ is the number of perturbed processes in the subpath of $`P`$ from $`p`$ to $`q`$. Since $`m\mathrm{min}(k,𝒟)`$, time accuracy is verified for $`p`$.
The arguments above show that time accuracy holds for all unperturbed processes within $`\mathrm{min}(k,𝒟)`$ rounds and that any subsequent state is $`2\mathrm{min}(k,𝒟)`$-accurate for unperturbed processes. For perturbed processes, similar reasoning applies. Instead of relying on S4 to establish the baseline clock value, we use instead a value bound by the construction given in Lemma 11’s proof. Within $`\mathrm{min}(k,𝒟)`$ rounds, there is a state $`\sigma ^{}`$ where perturbed $`p`$ has a clock value of at most $`3j`$, and $`j<\mathrm{min}(k,𝒟)`$ is the distance to some unperturbed process that executes S4 in the first round. The value of $`\text{clock}_p`$ cannot increase from $`3j`$ to $`3j+t`$ unless process $`q`$ has incremented its clock at least $`t\text{dist}_{pq}2m`$ times, where $`m`$ is at most $`\mathrm{min}(k,𝒟)`$. Therefore when $`\text{clock}_p=x`$ at some state following $`\sigma ^{}`$, we infer that $`\text{clock}_q`$ has incremented at least $`x3\mathrm{min}(k,𝒟)\text{dist}_{pq}2\mathrm{min}(k,𝒟)`$ times, which verifies time accuracy for unperturbed processes.
Theorem 2 addresses desired property (1) set out at the beginning of the section. Property (2) specifies that each faulty process clock be reset to zero and then advance as a phase clock. For the same reason that (1) has been weakened to the time accuracy of Definition 4, we weaken (2) to require only that each perturbed process be reset to some value in the range $`[0,3\mathrm{min}(k,𝒟)]`$ within $`k`$ rounds following the $`k`$-faulty initial state, and thereafter increments as a phase clock. Theorem 2 implies that subsequent increases to clock values satisfy a distance property relating the value of a clock to the number of increments of other clock variables. The following theorem states the weakened form of (2).
###### Theorem 3
*For any computation starting from a $`k`$-faulty initial state, $`k<n`$, each perturbed process clock is at most $`3\mathrm{min}(k,𝒟)`$ within $`\mathrm{min}(k,𝒟)`$ rounds and increases to value $`((t4)\mathrm{min}(k,𝒟))/2`$ within $`t`$ rounds; and each unperturbed process clock similarly increases to $`(t4)/2`$ within $`t`$ rounds after resetting by S4.*
* Lemma 11 directly shows that perturbed processes assign clock variables to at most $`3\mathrm{min}(k,𝒟)`$ within the first $`\mathrm{min}(k,𝒟)`$ rounds. Lemma 15 establishes that $`p`$ increases its clock to at least $`m=\mathrm{min}(((t2)\text{dist}_{pq})/2`$ after $`t+2`$ rounds following the execution of S4. Lemma 12 establishes that for each perturbed region, some process $`p`$ either within or neighboring the perturbed region executes S4 in the first round. Lemma 15 establishes that processes within a given distance increase their clock values as $`\text{clock}_p`$ increases. Any process $`q`$ within a perturbed region containing or neighboring $`p`$ is at most distance $`\mathrm{min}(k,𝒟)`$ from $`p`$; simplifying the bound of Lemma 15 using $`\mathrm{min}(k,𝒟)`$ as a distance upper bound yields a lower bound of $`\text{clock}_q((t2)\mathrm{min}(k,𝒟))/2`$ after $`t+2`$ rounds.
## 5 Embedded Timer
This section discusses use of the repair timer as a component in a system. Whereas Section 4 investigated properties of the repair timer in isolation, the results of this section are essentially composition theorems stating conditions under which the repair timer can be used as a tool to enable time-adaptive fault tolerance in a system.
Consider a system that uses the repair timer as one of its components. The term *core system* is used in this section to refer to all system components outside the repair timer; in other words, the entire system consists of the core system plus the repair timer. The elements of a process configuration (variables and registers) can be partitioned into those belonging to the repair timer and those belonging to the core system. The *timer projection* of a state is formed by removing all elements from each process configuration not relevant to the repair timer (that is, only clock, w, related image variables and register fields are retained). A *core projection* is formed by removing all repair timer elements from the state.
###### Requirement 1
*Output legitimacy $`_O`$ of the system is defined solely in terms of the core projection, that is, no repair timer variable is an output variable. Core system legitimacy, given by the predicate $`_C`$, is also defined with respect to the core projection; predicate $`_C`$ is independent of repair timer variables or register fields. The legitimacy predicate for the system is $`_C_T`$.*
The interface between core system and repair timer is illustrated in Figure 2. Communication between these two components occurs in each process, but is restricted to two methods: the core system can reset the clock and w variables, and the core system may read the current clock value. Henceforth the term *double-reset* is used to denote the assignment $`\text{clock},\text{w}0,0`$. Both S4’s assignment of Figure 1 and the core system’s assignment illustrated in Figure 2 are double-reset assignments.
So that results from Section 4 are applicable to the composite system, each process invokes the repair timer (statements S2S6) once in each cycle. Figure 1 includes S1 and S7 to present the repair timer in isolation, however in the context of a system invoking the repair timer, these two statements would be subsumed by statements reading registers at the beginning of a process cycle and writing registers at the end of a cycle.
A process configuration can be faulty with respect to the repair timer elements, the core system elements, or a combination of both elements. If a state $`\sigma `$’s core projection violates $`_C`$ then $`\sigma `$ is said to be *core-faulty*; if $`\sigma `$’s timer projection violates $`_T`$ then $`\sigma `$ is *timer-faulty*. While Definition 1 provides the basis for a precise characterization of a faulty repair timer, the situation for a general system can be ambiguous, as observed in Section 2.
###### Requirement 2
*If $`p`$’s process configuration is not core-faulty and $`\text{Register}_{qp}`$ is faulty at $`\sigma `$, then the presence of a fault at $`\sigma `$ can be detected from the variables of $`p`$ and the contents of $`\text{Register}_{qp}`$.*
In many cases it is not difficult to design a system satisfying Requirement 2, in spite of the ambiguity of a faulty process configuration — the requirement only specifies that $`p`$ detect the *presence* of a fault, and $`p`$ is not required to determine the fault’s location (fault identification remains ambiguous). Depending on the particular computation, $`p`$ may not detect a fault. For instance, $`q`$ may repair its configuration, changing the contents of $`\text{Register}_{pq}`$, before $`p`$ reads the register.
The importance of Requirement 2 is that nonfaulty $`p`$ has the capability to detect a fault, retain the current values of its output variables, and initiate repair procedures. Moreover, $`p`$ can “contain” the fault because it reacts before copying values from $`\text{Register}_{pq}`$ and transmitting them to other processes.
###### Requirement 3
*Each cycle of a process invokes the repair timer. If, after reading registers at the start of a cycle, a fault can be inferred (as described in Requirement 2) for process $`p`$, and if $`(\text{clock}>𝒯𝒟)`$, then $`p`$ executes a double-reset. No other statements of the core system change the clock or w variables; any number of statements of the core system may read the clock variable. The legitimacy predicate for the core system does not depend on the clock or w variables of the repair timer.*
###### Requirement 4
*If any process $`p`$ executes a double-reset resulting in a state $`\sigma `$, then within $`𝒯7𝒟`$ rounds following $`\sigma `$, the core system component of the state is legitimate.*
Requirement 4 means, for most core systems, that the core system stabilization time $``$ satisfies $`𝒯7𝒟`$. In essence, this is a constraint on $`𝒯`$, which is added to the constraint $`𝒯11𝒟`$ given in Section 3.
###### Lemma 16
*If the core system is self-stabilizing with stabilization time $``$ and satisfies Requirements 14, then the system is self-stabilizing with stabilization time $`+O(𝒯)`$, and the double-reset assignment executes at most once for each process in any computation.*
The proof of Lemma 16 rests on the independence of the core system and the repair timer, as specified by Requirement 3, and the fact that the core system stabilizes before there is any possibility of executing a second double-reset by any process. The requirements do not, however, preclude the design of the core system from depending on repair timer properties. For instance, proving stabilization time $``$ for the core system may depend on timer accuracy, since the core system can read clock variables during convergence to $`_C`$, and timer accuracy can be used in some circumstances to measure the progress of distributed algorithms and to allow processes to wait for such algorithms to stabilize.
More interesting than using the repair timer for stabilization is the use of the repair timer to enable time adaptive repair of output variables. The remainder of this section illustrates the use of the repair timer in two designs. Design 1 is a time adaptive system, repairing output variables in $`O(\mathrm{min}(k,𝒟))`$ rounds from any $`k`$-faulty initial state. The design requires that the core system use a sequence of repair procedures, following an idea developed in . Output variables of nonfaulty processes may change to illegitimate values during convergence, but all output variables satisfy $`_O`$ within $`O(\mathrm{min}(k,𝒟))`$ rounds and continue to satisfy $`_O`$ thereafter. Design 2 is not fully time adaptive, but illustrates another use of the repair timer: the system can repair output variables in $`O(r)`$ rounds from any $`k`$-faulty initial state, $`kr`$, and no nonfaulty process changes an output variable during repair.
###### Design 1
*The core system has $`𝒟`$ independent repair procedures, denoted $`\text{repair}^i`$ for $`1i𝒟`$. Each of the repair procedures uses its own set of variables, including variables that are intended to be copied to the core system’s output variables. Let $`\text{output}^i`$ denote the set of variables of $`\text{repair}^i`$ that correspond to the system’s output variables, and let $`\text{repair}_p^i`$ denote process $`p`$’s portion of $`\text{repair}^i`$. We suppose that the core system also prepares a set of variables $`\text{output}_C`$ intended to be copied to output variables. Each repair procedure is invoked in every process cycle and $`\text{repair}^i`$ is self-stabilizing to a predicate $`^i`$ within $``$ rounds. When $`_C`$ holds, the $`\text{output}^i`$ variables are equal to the system’s output variables, for $`1i𝒟`$, and $`\text{output}^C`$ is also equal to the system’s output. <br>Procedure $`\text{repair}^i`$ has the property that, if the initial state is $`j`$-faulty, for $`ji`$, then for all $`p`$, within $`hi`$ rounds, there occurs a state $`\sigma `$ such that for all $`p`$, the variables of $`\text{output}^i`$ satisfy $`_O`$ (modulo renaming or copying their values to the system outputs) at $`\sigma `$ and all subsequent states. To exploit the repair timer, we suppose a stronger convergence property for $`\text{repair}^i`$, namely that $`\text{output}_p^i`$ variables stabilize within $`hi`$ of the $`\text{R}_p^i`$-rounds. <br>Given any $`k`$-faulty initial state satisfying $`k𝒟`$, certain repair procedures agree on values for output variables: for $`ik`$, after $`\text{repair}^i`$ stabilizes $`\text{output}^i`$, any procedure $`\text{repair}^{\mathrm{}}`$ for $`\mathrm{}>i`$ stabilizes $`\text{output}^{\mathrm{}}`$ to the same values that $`\text{output}^i`$ has. Moreover, the core system stabilizes $`\text{output}^C`$ to the same values contained in the stabilized $`\text{output}^k`$ variables. The stabilized values of output sets are also constrained by distance from a fault: for any nonfaulty $`p`$ such that the minimum distance from $`p`$ to a faulty process is $`d`$, where $`d>k`$, then *all* of the $`\text{output}_p`$ sets stabilize to the same values already contained in $`p`$’s output variables. <br>The $`\text{output}_p`$ sets are copied to the output variables of process $`p`$ as follows. In each cycle, if $`\text{clock}_p=𝒯`$, then $`p`$ copies $`\text{output}_p^C`$ to its output variables. Otherwise, in each cycle, $`p`$ copies $`\text{output}_p^i`$ to its output variables where $`i`$ is the largest value satisfying $`1i𝒟`$ and $`(h+5)i\text{clock}_p`$. No $`\text{output}_p`$ set is copied to the output variables if $`(h+5)𝒟<\text{clock}_p<𝒯`$ holds.*
###### Theorem 4
*If a system for Design 1 satisfies Requirements 14 and $`=O(𝒟)`$, then within $`O(\mathrm{min}(k,𝒟))`$ rounds following any $`k`$-faulty initial state, the system output-stabilizes to $`_O`$.*
* Consider a $`k`$-faulty initial state. If $`k𝒟`$, then from $`=O(𝒟)`$ and Theorem 1, the system stabilizes to $``$ and hence $`_O`$ in $`O(𝒟)`$ rounds, which proves the conclusion. The remaining case is $`k<𝒟`$ for a $`k`$-faulty initial state. For this case, we first show that any faulty process clock is time accurate within the first $`k`$ rounds. Requirement 3 ensures that some process within distance $`k`$ from any faulty process executes a double-reset in the first round, and Theorem 2 implies subsequent time accuracy within $`k`$ rounds. The same argument implies that each nonfaulty process within distance $`k`$ from a faulty process has a time-accurate clock after at most $`k`$ rounds. All nonfaulty processes have time-accurate clock variables throughout the computation.
Design 1 specifies that some nonfaulty processes do not change their output variables by any repair procedure, so the proof obligation is to show that faulty processes and those nonfaulty processes within distance $`k`$ to a faulty process stabilize their output variables in $`O(k)`$ time. After $`k`$ rounds, all such processes have time-accurate clock variables. By definition of time accuracy for a $`k`$-faulty initial state, a time-accurate $`\text{clock}_p`$ variable with value $`t`$ implies that the number of $`\text{R}_p^k`$-rounds preceding in the computation is at least $`t5k`$. Procedure $`\text{repair}^k`$ converges within $`hk`$ of the $`\text{R}_p^k`$-rounds, so after time accuracy holds, a $`\text{clock}_p`$ value of $`hk+5k=(h+5)k`$ implies that variables of $`\text{output}_p^k`$ can be copied to $`p`$’s output variables. The conditions of Design 1 also justify copying $`\text{output}_p^j`$ to the output variables when $`\text{clock}_p(h+5)j`$ for $`k<j<𝒟`$.
Having established the safety of copying output sets to the output variables, the remaining obligation is to show that all such copying either completes within $`O(k)`$ rounds or that any subsequent copying will not affect $`_O`$. Theorem 3 implies that all processes within distance $`k`$ to a faulty process will, after time accuracy holds, increase their clock variables to $`(h+5)k`$ within $`O(k)`$ rounds and will not subsequently decrease their clock values below this value. Therefore, within $`O(k)`$ rounds, all processes within distance $`k`$ to a faulty process assign their output variables, while those processes further than distance $`k`$ from a fault do not assign their output variables to falsify $`_O`$ by any step of the computation.
The $`\text{repair}^i`$ procedures of Design 1 are independent, meaning that they do not share any of the variables they modify. Because the variables of $`\text{repair}^C`$ are inactive for nonfaulty processes during the period of stabilization, they are a resource for faulty processes: values from nonfaulty $`\text{output}_p^C`$ can be disseminated to other processes and used for the stabilization of $`\text{repair}^i`$ procedures. For details on this technique, illustrated in a synchronous computation model, the reader is referred to .
###### Design 2
*The core system uses procedure $`\text{repair}^r`$ with a set of variables denoted $`\text{output}^r`$ that are equal to the system output variables at a legitimate state. Procedure $`\text{repair}^r`$ stabilizes the $`\text{output}^r`$ variables to satisfy $`_O`$ within $`hr`$ time from any initial state that is $`j`$-faulty for $`jr`$; each faulty process $`p`$ stabilizes $`\text{output}_p^r`$ after at most $`hr`$ of the $`\text{R}_p^r`$-rounds occur. The $`\text{output}_p^r`$ sets are copied to output variables of process $`p`$ as follows. In each cycle, if $`\text{clock}_p(h+5)r`$, then $`p`$ copies $`\text{output}_p^r`$ to its output variables; for all other values of $`\text{clock}_p`$ process $`p`$ leaves its output variables unchanged. The $`\text{repair}^r`$ procedure stabilizes $`\text{output}_p^r`$ to values already contained in $`p`$’s output variables for any nonfaulty $`p`$.*
###### Theorem 5
*If a system for Design 2 satisfies Requirements 13, then within $`O(r)`$ rounds following any $`k`$-faulty initial state for $`kr`$, the system output-stabilizes to $`_O`$. No step modifies output variables of nonfaulty processes to values differing from those specified by $`_O`$.*
Theorem 5 can be verified by reasoning similar to the proof of Theorem 4. Design 2 is not self-stabilizing and Theorem 5 does not specify Requirement 4 as a condition. The fault tolerance of this design is limited to $`r`$ faulty processes.
## 6 Concluding Remarks
It is challenging to construct a system that can repair variables inflicted by transient faults. A reasonable methodology for such system construction is based on tools for fault detection and repair, and these tools must themselves satisfy properties of time adaptivity and stabilization. This paper presented a phase clock algorithm specialized for the task of fault repair. The designs presented in Section 5 show how the repair timer can be composed with other system components.
Although time adaptivity and self-stabilization are major themes for this paper, the repair timer can be useful even when neither full stabilization nor fast stabilization is needed, because it is convenient to reason about the progress of repair procedures by measuring elapsed time (which would otherwise be complicated due to possible corruption of time-measurement variables). An observer of the system located at process $`p`$ could monitor repair progress by repeatedly examining $`\text{clock}_p`$, possibly delaying critical activity until repair is complete.
Use of the repair timer can add overhead to repair procedures because each cycle of repair invokes the timer, and the clock variable only increments in relation to rounds. It could be that actual repair only involves a small subset of processes, but a clock variable will not, in general, increment $`t`$ times unless *all* processes at distance $`d`$ have completed $`td`$ cycles — including processes that are not involved in the repair. Thus the measurement of repair time in rounds could be overly pessimistic and cause processes to wait longer than necessary before they infer that repair is complete. Another slowing of repair timing results if a loose upper bound on the network diameter is used for $`𝒟`$ (an upper bound is typically proposed for dynamic networks) since $`𝒯`$, the “resting value” for the repair timer, is determined by the value $`𝒟`$.
## 7 Appendix: Proofs
* The essence of the proof is that neither S3 nor S6 increment a clock to a value two greater than any neighbor. Since Definition 2 involves x variables, the effect of statement S1 requires examination. Reading a register to assign an x variable only increases the accuracy of the image variable; in particular, given $`(b_{pq}b_{qp})`$ as a precondition, S1 does not falsify this condition, because $`p`$ and $`q`$ have clock values differing by at most one in the precondition. Therefore it suffices to verify that any change to $`\text{clock}_p`$ or $`\text{clock}_q`$ also satisfies the lemma. In a reset-free computation, only S3 and S6 change a clock variable. If S3 executes, incrementing $`\text{clock}_p`$, we have $`\text{clock}_p\text{x}_p[q]`$ as a precondition. Since $`b_{pq}`$, we have $`\text{clock}_p\{\text{clock}_q,\text{clock}_q+1\}`$ also as a precondition; thus the increment to $`\text{clock}_p`$ results in a state satisfying $`|\text{clock}_p\text{clock}_q|<2`$, verifying $`b_{pq}`$. The postcondition also satisfies $`b_{qp}`$, since the change to $`\text{clock}_p`$ does not alter the relation between $`\text{clock}_q`$ and $`\text{x}_q[p]`$. A similar argument applies to S6, and also to the case of $`q`$ incrementing its clock.
* By definition of a rising computation, each process has a lower bound on neighboring clock variables in its x variable, because clock values cannot decrease in a reset-free computation. Suppose $`p`$ is the first of $`(p,q)`$ to increment its clock. A precondition for this step is $`\text{clock}_p\text{x}_p[q]`$, which implies $`\text{clock}_p\text{clock}_q`$, which in turn implies $`\text{x}_q[p]\text{clock}_q`$. Consider two cases for this last inequality, (i) $`\text{x}_q[p]<\text{clock}_q`$ or (ii) $`\text{x}_q[p]=\text{clock}_q`$. For (i), process $`q`$ cannot increment $`\text{clock}_q`$, and this situation will persist until $`p`$ increments its clock sufficiently many times so that $`\text{clock}_p\text{clock}_q`$. It is straightforward to verify that $`p`$ does not increment $`\text{clock}_p`$ beyond $`\text{clock}_q+1`$, so for case (i) the first increment to $`\text{clock}_q`$ establishes $`(b_{pq}b_{qp})`$. For (ii), we deduce from the inequalities above $`(\text{clock}_q\text{clock}_p\text{clock}_p\text{clock}_q)`$ holds as precondition to $`p`$’s first increment step, and the inequalities with regard to the x variables are similar. So for case (ii), $`(b_{pq}b_{qp})`$ holds directly.
* By Lemma 2, neighboring $`(p,q)`$ establish $`(b_{pq}b_{qp})`$ at or before $`\sigma `$; by Lemma 1 such processes continue to satisfy this property for the remainder of the reset-free computation segment.
* Because we consider a based computation, and not a rising computation in this lemma, the invariance of $`(b_{pq}b_{qp})`$ stated in Lemma 1 is not applicable. Note that $`(p::\neg \text{gap}_p)`$ holds at a smooth state. The invariance of smoothness is therefore verified from the conditions of S3S6, since no gap exists at a smooth state and S3 preserves smoothness. It is also simple to verify that the least clock value, if smaller than $`𝒯`$, increments within two rounds from a smooth state, hence at most $`2𝒯=O(𝒟)`$ rounds are needed to obtain a state satisfying $`(p::\text{clock}_p=𝒯)`$. A similar argument shows that all w variables converge to $`3𝒟+1`$ within $`O(𝒟)`$ rounds.
* In the first round, $`p`$ reads neighboring clock values and detects local minimality. If $`p`$ increments in this round, the lemma holds; and if $`p`$ does not increment, it writes its clock and detects cEcho in the next round, and local minimality implies $`p`$ will increment $`\text{clock}_p`$ either by S3 or S6.
* Observe that $`(p::\text{w}_p=3𝒟+1)`$ holds at least until some clock exceeds $`𝒯𝒟`$ so that S4 can execute. And $`\text{w}_p=3𝒟+1\text{wBig}_p`$, so process $`p`$ does not execute the assignment of S5. This implies that the computation is reset-free until some clock obtains the value exceeding $`𝒯𝒟`$. Lemma 5 implies that each minimal clock value increments in any pair of rounds, which implies that the maximum of the set of clock values eventually grows as the computation proceeds. Let $`\text{clock}_p`$ be the first clock to attain the value $`8𝒟+1`$ at state $`\alpha `$. Thus $`(q:q𝒩_p:\text{clock}_q8𝒟)`$ holds prior to $`\alpha `$. More generally, it follows by induction that $`(q:\text{dist}_{qp}=k>0:\text{clock}_q8𝒟k)`$. Therefore, each clock value has incremented at least once prior to $`\alpha `$. Let $`\text{clock}_q`$ be the first clock to attain the value $`9𝒟+1`$ at state $`\beta `$. At state $`\beta `$, each process has incremented its clock twice in a reset-free computation, implying that each process has read all of its registers at least once in this reset-free computation. Therefore the computation segment beginning with $`\beta `$ is by definition a rising computation segment (at least until some clock exceeds $`𝒯𝒟`$). Now let $`\text{clock}_r`$ be the first clock to attain the value $`10𝒟+1`$ at state $`\gamma `$. At state $`\gamma `$, each process has incremented its clock at least once in a rising computation, and by Lemma 3, $`\gamma `$ is a smooth state.
* Lemma 6 shows that the computation contains a smooth state, so the obligation here is to show the $`O(𝒟)`$ time bound. By Lemma 5 each minimal clock value increments at least once in any two consecutive rounds, so within $`20𝒟+2`$ rounds, some clock attains the value $`10𝒟+1`$, establishing a smooth state.
* The proof begins with a claim on the first $`k`$ rounds of the based computation: within $`k`$ rounds there is a state satisfying
(1) $`(q:\text{dist}_{qr}=k:\text{w}_q2k\text{clock}_q2k)`$
(2) $`(q,j:j<k\text{dist}_{qr}=j:\text{w}_q3kj\text{clock}_q3kj)`$
The claim is shown by induction. The first state of the computation satisfies the claim for $`k=0`$ as the base case. Suppose the claim holds for $`k\mathrm{}`$ and consider two processes $`q`$ and $`s`$ such that $`\text{dist}_{rq}=\mathrm{}`$, $`s𝒩_q`$, and $`\text{dist}_{rs}=\mathrm{}+1`$. Let $`\sigma `$ be a state satisfying (1)–(2) for $`k=\mathrm{}`$. By the *cEcho* condition of S3, process $`q`$ does not increment $`\text{clock}_q`$ beyond $`2\mathrm{}`$ until process $`s`$’s image fields in $`\text{Register}_{sq}`$ have the appropriate values. In fact, these register fields may initially have the appropriate values, which would allow $`q`$ to increment clock and w variables to $`2\mathrm{}+1`$ by S2S3. However process $`q`$ cannot subsequently increment to $`2\mathrm{}+2`$ until the *cEcho* condition holds, which requires a cycle by $`s`$ (and all other neighbors). Process $`s`$ therefore observes $`y_s[q]2\mathrm{}+1`$ in its cycle and assigns at most $`2\mathrm{}+2`$ to its w and clock variables. Since $`\sigma `$ occurs at least by round $`\mathrm{}`$, the bound of $`2\mathrm{}+2`$ for $`s`$ variables applies within round $`\mathrm{}+1`$, which establishes (1) of the claim.
Condition (2) is also shown by induction. For $`k=0`$, the base case, (2) holds vacuously. Now suppose (2) holds for $`k\mathrm{}`$ and consider two processes $`q`$ and $`s`$ such that $`\text{dist}_{rq}=\mathrm{}`$, $`s𝒩_q`$, and $`\text{dist}_{rs}=\mathrm{}+1`$. Condition (1) places an upper bound on variables at distance $`\mathrm{}+1`$ from process $`r`$ within round $`\mathrm{}+1`$. Therefore $`\text{clock}_s2(\mathrm{}+1)`$ within round $`\mathrm{}+1`$. In moving from round $`\mathrm{}`$ to $`\mathrm{}+1`$, we consider the possibilities for process $`q`$ and $`\text{clock}_q`$. If $`\text{clock}_q`$ and $`\text{clock}_s`$ differ by more than one and process $`q`$ executes a cycle, then S5 resets $`\text{clock}_q`$; before any further change to $`\text{clock}_q`$ occurs, the cEcho condition requires a full cycle by $`s`$, which validates (2) up to distance $`\mathrm{}+1`$ within round $`\mathrm{}+1`$. If $`\text{clock}_q`$ and $`\text{clock}_s`$ are equal or differ by one, then $`\text{clock}_q`$ could increment. Observe here that no clock or w variable can increment beyond one more than any neighboring value; by another inductive argument, no clock or w variable increments beyond $`j`$ more than any corresponding variable at distance $`j`$. Therefore $`\text{clock}_q`$ does not increment beyond $`(2\mathrm{}+2)+1`$ so long as $`\text{clock}_s2\mathrm{}+2`$. This observation is generalized by (2) for $`k=\mathrm{}+1`$ within round $`\mathrm{}+1`$. Note that we have assumed that any clock increment is due to S3 and not S6 in this argument; this assumption is justified by (1), since $`\text{w}<3𝒟+1`$, which disables execution of S6.
* Let $`\sigma `$ be a state satisfying $`(q::\text{clock}_q3𝒟)`$. By Lemma 8 such a state $`\sigma `$ occurs with $`𝒟`$ rounds of the based computation. So long as every clock is at most $`𝒯𝒟`$, no step subsequent to $`\sigma `$ decreases a w variable; and if no w variable is reset by S4 in a consecutive pair of rounds, then the minimum value of the set of w variables either increases by that pair of rounds or all w variables already have the maximum $`3𝒟+1`$ value (we consider a consecutive pair of rounds to ensure that *wEcho* will hold for S2). Therefore, if no clock variable attains the value $`7𝒟+1`$ within $`2(3𝒟+1)`$ rounds, all w variables equal $`3𝒟+1`$ and the lemma holds. On the other hand, if some clock does attain the value $`7𝒟+1`$, we shall deduce that all w values equal $`3𝒟+1`$, which also proves the lemma. The argument rests on the following claim: at all states subsequent to $`\sigma `$ satisfying $`(p::\text{clock}_p7𝒟)`$, the implication $`\text{clock}_p3𝒟+k\text{w}_pk`$ holds for every $`p`$ and $`0k3𝒟+1`$. This claim is verified by induction on $`k`$. For $`k=0`$ the result is immediate from the domain of w variables. Now consider $`k>0`$ and suppose the claim holds for $`k1`$. Let $`q`$ be the first process to assign $`\text{clock}_q3𝒟+k`$. If the assignment occurs by S6 then $`w=3𝒟+1`$ and the claim holds; if the assignment occurs by S3, then each neighbor of $`q`$ has a clock value of $`3𝒟+(k1)`$, hence by hypothesis each neighboring w variable is at least $`k1`$, and $`\text{w}_qk1`$ by the same hypothesis. The result is that the same cycle assigning $`\text{clock}_q3𝒟+k`$ also assigns $`\text{w}_q`$ to be at least $`k`$. Similar arguments treat the general case for $`q`$ (not necessarily the first) assigning $`3𝒟+k`$ to $`\text{clock}_q`$, verifying that $`\text{w}_qk`$ as a result. To complete the lemma, consider the first state $`\delta `$ where some $`\text{clock}_q`$ has value $`7𝒟+1`$. By the induction argument given in the proof of Lemma 8, any clock at distance $`j`$ from $`\text{clock}_q`$ has had a value of at least $`7𝒟j`$ prior to state $`\delta `$. Therefore every clock has contained a value of at least $`6𝒟+1`$ prior to $`\delta `$, implying that each w variable is at least $`3𝒟+1`$ prior to $`\delta `$. The state immediately preceding $`\delta `$ thus satisfies proof obligation.
* by induction on $`t`$. For $`t=0`$ let $`\sigma ^{}=\sigma `$ to satisfy the base case. For $`t>0`$, we have $`\text{clock}_q3t\text{w}_q3t`$ by hypothesis. By the Echo conditions of S2, S3 and S6, the clock and w values of $`q`$ remain at most $`3t`$ until all neighbors either (i) complete cycles that observe these values and write corresponding images to output registers or (ii) happen to have these values already in their output registers.
Considering (i), for $`r𝒩_q`$ satisfying $`\text{dist}_{pr}=t+1`$, the execution of S2 assures $`\text{w}_r3t+1`$ within one round, and $`\text{clock}_r`$ is at most $`3t+1`$ if $`r`$ observes no gap, or assigned some value at most $`w_r`$ otherwise; either case verifies the inductive hypothesis for $`t+1`$. These considerations for (i) also verify the second part of the lemma, which concerns a path of unperturbed processes, and the same hypothesis with $`3t`$ replaced by $`t`$.
Considering (ii), process $`q`$ may increment $`\text{clock}_q`$ and $`\text{w}_q`$ because $`r𝒩_q`$ happens already to have values corresponding to $`\text{clock}_q`$ and $`\text{w}_q`$ in its output register fields. In this case, $`q`$ may increment its variables to at most $`3t+1`$ immediately. Furthermore process $`r`$ may initially have its program counter at S7, about to write its image variables in such a way that $`q`$ can observe the cEcho condition (even though $`r`$ would not actually read and write in a full cycle). Therefore, if $`r`$ executes S7, process $`q`$ can increment variables again to at most $`3t+2`$. However, here a cEcho condition will not be satisfied at $`q`$ until all neighbors complete full cycles, so $`q`$’s variables cannot exceed $`3t+2`$ until $`r`$ completes a cycle. When $`r`$ does complete a cycle, by the reasoning above for (i) we deduce that $`\text{clock}_r3t+3`$ and $`\text{w}_r3t+3`$ for $`r𝒩_q`$.
* Note that the lemma holds trivially if the initial state is $`n`$-perturbed. For the case $`k<n`$ we use induction on $`t`$ and nested induction on $`k`$ and suppose a based computation. For the base case $`t=0`$ consider $`p𝒩_q`$. Since $`q`$ is perturbed, there is a path $`P`$ from $`p`$ to some perturbed $`r`$ (possibly through $`q`$) of $`k+2`$ or fewer processes, which is not smooth. Because $`\text{clock}_p=𝒯`$, some neighboring pair of processes along path $`P`$ has the property that one clock exceeds $`𝒯𝒟`$ while the other is less than $`𝒯𝒟`$. Therefore some process in path $`P`$ executes S4 in the first round. By the arguments of Lemma 11 it follows that $`p`$ executes S4 within $`k+2`$ rounds. This completes the base case, but reasoning similar to the nested induction also applies for $`t>0`$. Finally, because the initial state may not justify a based computation, two additional rounds are added to conclude a $`k+t+4`$ bound.
* Lemma 10 states that a process executes S4 at most once in a computation, so it suffices to show that $`p`$ either does not execute S4 or executes S4 within the first $`4+\mathrm{min}(𝒟,t+k)`$ rounds. If $`p`$ is unperturbed, Lemma 12 implies the result. If $`p`$ is perturbed, then for some perturbed region $`P`$ containing $`p`$, there is an unperturbed $`q`$ neighboring some process of $`P`$ that executes S4 within the first $`4+\mathrm{min}(𝒟,k)`$ rounds by Lemma 12. Applying Lemma 11 we deduce that $`\text{clock}_p3\mathrm{min}(𝒟,k)`$ holds after $`\mathrm{min}(𝒟,k)`$ additional rounds, and by arguments of Lemmas 8 and 9 process $`p`$ does not execute S4 in the remainder of the computation. Therefore, for perturbed $`p`$, the distance from $`p`$ to a perturbed process is $`t=0`$ and after $`4+\mathrm{min}(𝒟,k)+\mathrm{min}(𝒟,k)`$ rounds, process $`p`$ does not execute S4.
* by induction on $`t`$. The base case $`t=0`$ trivially follows from the domain of w variables, which have non-negative values. The same observation concerning the domain of w variables simplifies the proof obligation to the case $`\text{dist}_{pq}t`$. It is useful also to observe base cases for $`t=1`$ and $`t=2`$, since by the end of round two the computation is based, which simplifies reasoning for higher rounds. For $`t=1`$ the verification is again trivial by the domain of w variables. For $`t=2`$, it is required to show that by the end of round two, $`\text{w}_p1`$. In fact any change to $`\text{w}_p`$ is an increase from its original value of zero, and at least one increment occurs because $`\text{wEcho}_p`$ is observed by $`p`$ within two rounds following $`\sigma `$. No subsequent reduction to $`\text{w}_p`$ results in a value less than one, since $`\text{wMin}_p`$ is at least zero at all states. This verifies the base case for $`t=2`$.
Now suppose the hypothesis $`\text{w}_q(t\text{dist}_{pq})/2`$ for every $`q`$ such that $`\text{dist}_{pq}t`$ at some state $`\sigma ^{}`$. Note that no such process $`q`$ subsequently executes S4 in the computation, by Lemma 11; therefore any subsequent change to $`\text{w}_q`$ occurs by S2. If S2 assigns $`\text{w}_q`$ a value at least $`((t+1)\text{dist}_{pq})/2`$ in the round following $`\sigma ^{}`$, or if $`\text{w}_q`$ already has such a value and does not decrease, then the induction step is verified. Therefore we consider the possibility that $`\text{w}_q`$ either remains unchanged or decreases below $`(t\text{dist}_{pq})/2`$ by execution of S2. A decrease only occurs if $`\text{w}_q>\text{wMin}_q+1`$, so a decrease below $`(t\text{dist}_{pq})/2`$ is only possible if there is a neighbor $`r𝒩_q`$ satisfying $`\text{y}_q[r](t\text{dist}_{pq})/22`$, which would in turn imply that such a value existed in $`\text{w}_r`$ in the previous round. But by hypothesis, $`\text{w}_r(t\text{dist}_{pr})/2`$, and since $`r𝒩_q`$ the value of $`\text{w}_r`$ is at least $`(t\text{dist}_{pq}\pm 1)/2`$, which contradicts $`\text{y}_q[r](t\text{dist}_{pq})/22`$. Therefore such a decrease to $`\text{w}_q`$ cannot occur.
The remaining case to consider is that $`\text{w}_q=(t\text{dist}_{pq})/2`$ and does not change in the round following $`\sigma ^{}`$. Here there are two cases for $`t`$ and $`q`$, either $`(t\text{dist}_{pq})`$ is even or it is odd. If $`(t\text{dist}_{pq})`$ is even, then $`(t\text{dist}_{pq})/2`$ is equal to $`((t+1)\text{dist}_{pq})/2`$ and the hypothesis for $`(t+1)`$ is proved — the value of $`\text{w}_q`$ can remain unchanged in the round following $`\sigma ^{}`$ and satisfy the hypothesis. If, however, $`(t\text{dist}_{pq})`$ is odd, then $`\text{w}_q`$ is required to increment to verify the hypothesis for $`(t+1)`$. Observe that if $`(t\text{dist}_{pq})`$ is odd, then $`(t\text{dist}_{pq})/2`$ is equal to $`((t1)\text{dist}_{pq})/2`$, so we infer that $`\text{w}_q=(t\text{dist}_{pq})/2`$ held at round $`(t1)`$ (here we assume the hypothesis not only for $`t`$, but $`(t1)`$ as well, which is permissible because base cases for $`t=1`$ and $`t=2`$ have been verified). Therefore by round $`(t+1)`$, process $`q`$ observes $`\text{wEcho}_q`$ and increments $`\text{w}_q`$, which verifies the hypothesis for $`(t+1)`$.
* by induction on $`t`$, for $`t0`$. Note that round $`t+2`$ occurs in a based computation, since within two rounds following $`\sigma `$ the computation is based. The base case for induction is shown for $`t=0`$ and $`t=1`$, since the main induction step relies on two previous rounds of a based computation. For $`t1`$, since every clock variable is at least zero, the base cases are verified directly by the domain of clock variables — which are at least zero at any state.
Note that for any $`t`$, $`t2\text{dist}_{pq}`$ trivially satisfies the conclusion because clock variables are always at least zero; therefore in the remainder of the proof we consider only the case of $`q`$ and $`t`$ satisfying $`t2>\text{dist}_{pq}`$. Now suppose the hypothesis holds for $`t1`$ and $`t2`$, $`t2`$, aiming to show that the hypothesis also holds for $`t`$, that is, that clock variables satisfy the specified lower bound by the end of round $`t+2`$.
By Lemmas 11 and 12, by round $`t`$, any process in the set $`R=\{r|\text{dist}_{pr}t3\}`$. has either executed S4 or will not do so throughout the remainder of the computation. Therefore in round $`t+2`$, any reduction to $`\text{clock}_r`$ for $`rR`$ could only occur by S5. Lemma 14 establishes that $`\text{w}_r((t+1)\text{dist}_{pr})/2`$ holds invariantly following round $`t+1`$. So if process $`r`$ executes S5, the result satisfies $`\text{clock}_r((t2)\text{dist}_{pr})/2`$, which would verify the inductive hypothesis for $`r`$ and round $`t+2`$. If $`r`$ does not execute S5 in round $`x+1`$, then consider two cases for $`r`$.
Case: $`t\text{dist}_{pr}`$ is even. Observe that $`((t2)\text{dist}_{pr})/2`$ differs from $`((t3)\text{dist}_{pr})/2`$, meaning that the obligation is to show that $`\text{clock}_r`$ is either at least $`((t2)\text{dist}_{pr})/2`$ by the end of round $`t+1`$, or that $`\text{clock}_r`$ increments during round $`t+2`$. If the former holds, the hypothesis is proved, so suppose $`\text{clock}_r=((t3)\text{dist}_{pr})/2`$ at the end of round $`t+1`$. Because $`t\text{dist}_{pr}`$ is even, $`\text{clock}_r((t4)\text{dist}_{pr})/2`$ by hypothesis for $`t2`$. But this implies that during round $`t+1`$, the value of $`\text{clock}_r`$ either did not change or was reduced by S5. However a reduction by S5 would satisfy the hypothesis for $`t`$ as well, because of Lemma 14’s bound on w variables. The only remaining possibility is that $`\text{clock}_r`$ does not change in round $`t+1`$, implying that $`r`$ observes cEcho during round $`t+2`$. Therefore, if $`\text{clock}_r\text{cMin}_r`$ when $`r`$ observes cEcho, then $`\text{clock}_r`$ will increment either by S3 or S6. To show that $`r`$ does indeed observe cEcho, we use the hypothesis for $`t1`$ and each $`q𝒩_r`$. If $`\text{dist}_{pq}\text{dist}_{pr}`$, then by round $`t+1`$ (and throughout round $`t+2`$) the relation $`\text{clock}_q\text{clock}_r`$ holds at least until $`r`$ increments its clock. If $`\text{dist}_{pq}=\text{dist}_{pr}+1`$, then $`((t2)\text{dist}_{pr})/2`$ and $`((t2)\text{dist}_{pq})/2`$ are equal, and again the relation $`\text{clock}_q\text{clock}_r`$ holds at least until $`r`$ increments its clock.
Case: $`t\text{dist}_{pr}`$ is odd. A similar detailed argument can be given for this case, but there is a simpler approach: $`((t2)\text{dist}_{pr})/2`$ and $`((t3)\text{dist}_{pr})/2`$ are equal, so the hypothesis for $`t1`$ and $`r`$ directly suffice to verify the hypothesis for $`t`$.
* In any computation, either some process executes a double-reset or no process does so. In the latter case, the core system component stabilizes within $``$ rounds, and the repair timer concurrently reaches the timer-final condition within $`O(𝒯)`$ rounds by Theorem 1. This demonstrates $`+O(𝒯)`$ stabilization time if no double-reset occurs; the same argument applies to the case where any double-reset occurs by S4 and not by the core system. Lemma 10 implies that a double-reset occurs at most once for each process in this case.
Now consider the possibility that the core system executes a double-reset at least once in a computation. All such assignments cease after the base system stabilizes, which occurs within $``$ rounds, so the system stabilization time is $`+O(𝒯)`$. To show that any process executes a double-reset at most once, we demonstrate that the core system stabilizes before $`\text{clock}>𝒯𝒟`$ holds at any process, since Requirement 3 prevents repeated resets of the clock so long as $`\text{clock}𝒯𝒟`$.
If any double-reset assignment occurs, then within $`𝒟`$ rounds thereafter, a state $`\sigma `$ occurs such that each clock is at most $`3𝒟`$ by Lemma 11, and also within $`𝒟`$ rounds, time-accuracy holds and is invariant thereafter by Theorem 2. Although Theorem 2 is conditioned on $`k<n`$ for a $`k`$-perturbed initial state, its proof arguments are valid for the case of an $`n`$-faulty initial state, provided some process executes a double-reset in the first round. While we do not suppose that a double-reset occurs in the first round, the state preceding the first double-reset can be considered as the initial state for the subsequent computation, so that Theorem 2’s results apply for the suffix computation. Time accuracy for the extreme case of an $`n`$-faulty initial state implies for $`\text{clock}=t`$ that at least $`(t𝒟5𝒟)=t6𝒟`$ rounds have transpired. Therefore, if $`𝒯𝒟X+6𝒟`$, where $`X`$ is the number of rounds needed for stabilization, then as soon as time accuracy holds, no clock increases beyond $`𝒯𝒟`$ until the core system has stabilized. Requirement 4 implies stabilization within $`𝒯7𝒟`$ rounds, which ensures that the core system stabilizes before there is the possibility of a second double-reset. To complete the proof we address the period between the first double-reset and before time accuracy holds. This is at most $`𝒟`$ rounds, and it is easy to show that no clock increases from zero to beyond $`𝒯𝒟`$ within $`𝒟`$ rounds, so a second double-reset does not occur in the period before time accuracy holds.
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# 1 The predicted branching ratio for the 𝐵⁻ Δ𝑆=2 two - body decays calculated using the factorization approach within Standard Model (the first column), Minimal Supersymmetric Standard Model (the second column), Minimal Supersymmetric Standard Model extended by ℛ parity breaking (the third column), and Two Higgs Doublet Model (the fourth column). The values in columns two, three, four are upper limits, as determined from present knowledge of upper limits for couplings involved.
The intensive search for physics beyond the Standard Model (SM) is performed nowdays in various areas of particle physics. Among these, rare B meson decays are suggested to give good opportunities for discovering new physics beyond SM . Recently, it has been suggested to investigate effects of new physics possibly arising from $`bss\overline{d}`$ or $`bdd\overline{s}`$ decays. As shown in Ref. , the $`bss\overline{d}`$ transition is mediated in the standard model by the box-diagram and its calculation results in a branching ratio of nearly $`10^{11}`$, the exact value depending on the relative unknown phase between t, c contributions in the box. The $`bdd\overline{s}`$ branching ratio is even smaller by a factor of $`10^2`$, due to the relative $`|V_{td}/V_{ts}|`$ factor in the amplitudes. In Ref. different scenarios were used in the analysis of the $`bdd\overline{s}`$ decay, which might be important in $`B^\pm K\pi `$ decays. The authors of Refs. have calculated the $`bss\overline{d}`$ transition in various extensions of the SM. It appears that for certain plausible values of the parameters, this decay may proceed with a branching ratio of $`10^810^7`$ in the minimal supersymmetric standard model (MSSM) and in two Higgs doublet models .
Thus, decays related to the $`bss\overline{d}`$ transition which was calculated to be very rare in the Standard Model, provide a good opportunity for investigating beyond the Standard Model physics. In Ref. it was suggested that the most suitable channels to see effects of the $`bss\overline{d}`$ transition are the $`B^{}K^{}K^{}\pi ^+`$ or $`\overline{B}^0K^{}K^{}\pi ^+\pi ^+`$ decays. Moreover, when one considers supersymmetric models with $``$-parity violating couplings, it turned out that the existing bounds on the involved couplings of the superpotential did not provide any constraint on the $`bss\overline{d}`$ mode . Recently, the OPAL collaboration has set bounds on these couplings from the establishment of un upper limit for the $`B^{}K^{}K^{}\pi ^+`$ decay $`BR(B^{}K^{}K^{}\pi ^+)1.3\times 10^4`$. The long distance effects in $`B^{}K^{}K^{}\pi ^+`$ decay have also been estimated recently and they have been found to be of the order $`10^{12}`$, comparable in size with the short - distance SM contribution, thus leaving this decay ”free” for the search of new physics. Although it appears that $`B^{}K^{}K^{}\pi ^+`$ or $`\overline{B}^0K^{}K^{}\pi ^+\pi ^+`$ are very good candidates to search for the $`\mathrm{\Delta }S=2`$ transitions, we investigate here another possibilty for the observation of the $`bss\overline{d}`$ transition: the two body decays of $`B^{}`$.
We consider the $`VV`$, $`VP`$, $`PP`$ states. Although in principle two body decays would appear to be simpler to analyze, there is the complication of $`K^0\overline{K}^0`$ mixing. Hence one needs also a good estimate for the $`bs\overline{s}d`$ transitions as well. Nevertheless, not all the two-body states involve neutral $`K^{}s`$ and we shall return to this point in our summary. First, we proceed to describe the framework used in our analysis in which we concentrate on MSSM, with and without $``$ parity and two Higgs doublet models as possible alternatives to the SM.
The minimal supersymmetric extension of the Standard Model leads to the following effective Hamiltonian describing the $`bss\overline{d}`$ transition
$``$ $`=`$ $`\stackrel{~}{C}_{MSSM}(\overline{s}\gamma ^\mu d_L)(\overline{s}\gamma _\mu b_L),`$ (1)
where we have denoted
$`\stackrel{~}{C}_{MSSM}`$ $`=`$ $`{\displaystyle \frac{\alpha _s^2\delta _{12}^d\delta _{23}^d}{216m_{\stackrel{~}{d}}^2}}[24xf_6(x)+66\stackrel{~}{f}_6(x)]`$ (2)
with $`x=m_{\stackrel{~}{g}}^2/m_{\stackrel{~}{d}}^2`$, and the functions $`f_6(x)`$ and $`\stackrel{~}{f}_6(x)`$ are given in . The couplings $`\delta _{ij}^d`$ parametrize the mixing between the down-type left-handed squarks. At the scale of $`b`$ quark mass and by taking the existing upper limits on $`\delta _{ij}^d`$ from and the coupling $`\stackrel{~}{C}_{MSSM}`$ is estimated to be $`|\stackrel{~}{C}_{MSSM}|1.2\times 10^9`$ $`GeV^2`$ for an average squark mass $`m_{\stackrel{~}{d}}=500`$ $`GeV`$ and $`x=8`$, which leads to an inclusive branching ratio for $`bss\overline{d}`$ of $`2\times 10^7`$ . The corresponding factor calculated in SM is found to be
$`C_{SM}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{G_F^2}{2\pi ^2}}m_W^2V_{tb}V_{ts}^{}[V_{td}V_{ts}^{}f({\displaystyle \frac{m_W^2}{m_t^2}})+V_{cd}V_{cs}^{}{\displaystyle \frac{m_c^2}{m_W^2}}g({\displaystyle \frac{m_W^2}{m_t^2}},{\displaystyle \frac{m_c^2}{m_W^2}})]]`$ (3)
with $`f(x)`$ and $`g(x,y)`$ given in . Taking numerical valus from , neglecting the CKM phases, one estimates $`|C_{SM}|4\times 10^{12}`$ $`GeV^2`$.
The authors of have also investigated beyond MSSM cases by including $`R`$\- parity violating interactions. The part of the superpotential which is relevant here is $`W=\lambda _{ijk}^{}L_iQ_jd_k`$, where $`i,j,k`$ are indices for the families and $`L,Q,d`$ are superfields for the lepton doublet, the quark doublet, and the down-type quark singlet, respectively. Following notations of and the tree level effective Hamiltonian is
$``$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{f_{QCD}}{m_{\stackrel{~}{\nu }_n}^2}}[\lambda _{n32}^{}\lambda _{n21}^{}(\overline{s}_Rb_L)(\overline{s}_Ld_R)+\lambda _{n21}^{}\lambda _{n32}^{}(\overline{s}_Rd_L)(\overline{s}_Lb_R)].`$ (4)
The QCD corrections were found to be important for this transition . For our purpose it suffices to follow retaining the leading order QCD result $`f_{QCD}2`$, for $`m_{\stackrel{~}{\nu }}=100`$ $`GeV`$.
Most recently an upper bound on the specific combination of couplings entering (4) has been obtained by OPAL from a search for the $`B^{}K^{}K^{}\pi ^+`$ decay
$`{\displaystyle \underset{n}{}}\sqrt{|\lambda _{n32}^{}\lambda _{n21}^{}|^2+|\lambda _{n21}^{}\lambda _{n32}^{}|^2}<10^4.`$ (5)
Here we take the order of magnitude, while the OPAL result is $`5.9\times 10^4`$ based on a rough estimate $`\mathrm{\Gamma }(B^{}K^{}K^{}\pi ^+)1/4`$ $`\mathrm{\Gamma }(bss\overline{d})`$.
The decay $`bss\overline{d}`$ has been investigated using two Higgs doublet models (THDM) as well . These authors found that the charged Higgs box contribution in MSSM is negligible. On the other hand, THDM involving several neutral Higgses could have a more sizable contribution to these modes. The part of the effective Hamiltonian relevant in our case is the tree diagram exchanging the neutral Higgs bosons $`h`$ (scalar) and $`A`$ (pseudoscalar)
$`_{TH}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\xi _{sb}\xi _{sd}({\displaystyle \frac{1}{m_h^2}}(\overline{s}d)(\overline{s}b){\displaystyle \frac{1}{m_A^2}}(\overline{s}\gamma _5d)(\overline{s}\gamma _5b)),`$ (6)
with the coupling $`\xi _{ij}`$ defined in as a Yukawa coupling of the FCNC transitions $`d_id_j`$. In our estimation we use the bound $`|\xi _{sb}\xi _{sd}|/m_H^2>10^{10}`$ $`GeV^2`$, $`H=h,A`$, which was obtained in by using the $`\mathrm{\Delta }m_K`$ limit on $`\xi _{bd}/m_H`$ and assuming $`|\xi _{sb}/m_H|>10^3`$.
We proceed now to study the effect of Hamiltonians (1), (4), (6) on the various two body $`\mathrm{\Delta }S=2`$ decays of charged B - mesons. In order to calculate the matrix elements of the operators appearing in the effective Hamiltonian, we use the factorization approximation , which requires the knowledge of the matrix elements of the current operators or the density operators. Here we use the standard form factor representation of the following matrix elements:
$`P^{}(p^{})|\overline{q}_j\gamma ^\mu q_i|P(p)`$ $`=`$ $`F_1(q^2)(p^\mu +p^\mu {\displaystyle \frac{m_P^2m_P^{}^2}{q^2}}(p^\mu p^\mu ))`$ (7)
$`+`$ $`F_0(q^2){\displaystyle \frac{m_P^2m_P^{}^2}{q^2}}(p^\mu p^\mu ),`$
where $`F_1`$ and $`F_0`$ contain the contribution of vector and scalar states respectively and $`q^2=(pp^{})^2`$. Also, $`F_1(0)=F_0(0)`$ . For these form factors, one usually assumes pole dominance
$`F_1(q^2)`$ $`=`$ $`{\displaystyle \frac{F_1(0)}{1\frac{q^2}{m_V^2}}};F_0(q^2)={\displaystyle \frac{F_0(0)}{1\frac{q^2}{m_S^2}}}`$ (8)
and in order to simplify, we shall take $`m_V=m_S`$. The matrix element between pseudoscalar and vector meson is usually decomposed as
$`V(q,ϵ_V)|\overline{q}_j\gamma ^\mu (1\gamma _5)q_i|P(p)=`$
$`=`$ $`{\displaystyle \frac{2V(Q^2)}{m_P+m_V}}ϵ^{\mu \nu \alpha \beta }ϵ_{V\nu }^{}p_\alpha q_\beta +iϵ_V^{}Q{\displaystyle \frac{2m_V}{Q^2}}Q^\mu \left(A_3(Q^2)A_0(Q^2)\right)`$
$`+`$ $`i(m_P+m_V)\left[ϵ_V^\mu A_1(Q^2){\displaystyle \frac{ϵ_V^{}Q}{(m_P+m_V)^2}}(p+q)^\mu A_2(Q^2)\right],`$
where $`Q=pq`$.
$$A_3(Q^2)\frac{m_H+m_V}{2m_V}A_1(Q^2)+\frac{m_Hm_V}{2m_V}A_2(Q^2)=0,$$
(10)
and $`A_3(0)=A_0(0)`$. For the vector and axial vector form factor we use again pole dominance , and relevant parametrs are taken from $`F_0^{BK}(0)=0.38`$, $`A_0^{BK}(0)=0.32`$. For the calculations of the density operators we use the relations
$`^\alpha (\overline{s}\gamma _\alpha b)`$ $`=`$ $`i(m_bm_s)\overline{s}b`$ (11)
and
$`^\alpha (\overline{s}\gamma _\alpha \gamma _5b)`$ $`=`$ $`i(m_b+m_s)\overline{s}\gamma _5b`$ (12)
We will use also the following decay constants:
$$V(ϵ_V,q)|\overline{q}_j\gamma ^\mu q_i|0=ϵ_\mu ^{}(q)g_V(q^2),$$
(13)
and
$$P(q)|\overline{q}_j\gamma ^\mu \gamma _5q_i|0=if_Pq_\mu $$
(14)
with $`f_K=0.162`$ $`GeV`$, $`g_K^{}=0.196`$ $`GeV^2`$ . Now we turn to the analysis of the specific modes.
a) $`B^{}K^{}\overline{K}^0`$ decay
For the analysis of pseudoscalar meson decay to two vector mesons it is convenient to use helicity formalism (see e.g. ). We denote $`𝒪=(\overline{s}\gamma ^\mu (1\gamma _5)d)`$ $`(\overline{s}\gamma _\mu (1\gamma _5)b)`$, and then we use $`=C𝒪`$ with $`C`$ being $`1/4\stackrel{~}{C}_{MSSM}`$, $`1/4C_{SM}`$. Using factorization and the definitions given above, one finds the following helicity amplitudes
$`H_{00}(B^{}K^{}\overline{K}^0)=Cg_K^{}(m_B+m_K^{})[\alpha A_1^{BK^{}}(m_K^{}^2)\beta A_2(m_K^{}^2)]`$ (15)
$`H_{\pm \pm }(B^{}K^{}\overline{K}^0)=Cg_K^{}(m_B+m_K^{})[\alpha A_1^{BK^{}}(m_K^{}^2)\gamma V^{BK^{}}(m_K^{}^2)]`$ (16)
where
$`\alpha `$ $`=`$ $`{\displaystyle \frac{12r^2}{2r^2}},\beta ={\displaystyle \frac{k^2}{2r^2(1+r)^2}},\gamma =(14r^2)`$ (17)
with $`r=m_K^{}/m_B`$, $`k^2=1+r^4+t^42r^2`$ $`2t^22r^2t^2`$. The decay width is then
$`\mathrm{\Gamma }(B^{}K^{}\overline{K}^0)`$ $`=`$ $`{\displaystyle \frac{|\stackrel{}{p}|}{8\pi m_B^2}}[|H_{00}|^2+|H_{++}|^2+|H_{}|^2].`$ (18)
Within MSSM model the branching ratio becomes $`6.2\times 10^9`$, while SM gives this rate to be $`6.8\times 10^{14}`$. The $``$ \- parity term described by the effective Hamiltonian (4) cannot be seen in this decay mode when factorization approach is used, since the density operator matrix element $`\overline{K}^0|(\overline{s}d)|0`$ vanishes. The two Higgs doublet model also cannot be tested in this mode due to the same reason.
b) $`B^{}K^{}\overline{K}^0`$ decay
The matrix element of the operator $`𝒪`$ is calculated to be
$`\overline{K}^0(k_0)K^{}(k_{},ϵ)|𝒪|B^{}(p_B)`$ $`=`$ $`2m_K^{}f_KA_0^{BK^{}}(m_K^{}^2)ϵ^{}k_0`$ (19)
Denoting the decay amplitude by $`𝒜`$, one finds
$`{\displaystyle \underset{pol}{}}|𝒜|^2`$ $`=`$ $`|C|^2f_K^2|A_0(m_K^2)|^2\lambda (m_B^2,m_K^2,m_K^2)`$ (20)
with the $`\lambda (a,b,c)=a^2+b^2+c^22(ab+bc+ac)`$. The branching ratio is straightforwardly found to be $`BR(B^{}K^{}\overline{K}^0)_{MSSM}1.6\times 10^9`$, which is comparable to the SM prediction of Ref. for the $`\mathrm{\Delta }S=0`$ $`B^{}K^{}K^0`$ decay, given as $`BR(B^{}K^{}K^0)=1\times 10^9`$, $`5\times 10^9`$, $`2\times 10^9`$ obtained for the number of colours $`N_c=2`$, $`N_c=3`$, $`N_c=\mathrm{}`$, respectively.
The SM calculation for the $`\mathrm{\Delta }S=2`$ transition leads to $`BR(B^{}K^{}\overline{K}^0)_{SM}=1.7\times 10^{14}`$. The MSSM which includes $``$ parity breaking terms can occur in this decay. The matrix element of the operator $`𝒪_{}=(\overline{s}(1+\gamma _5)d)`$ $`(\overline{s}(1\gamma _5)b)`$ can be found to be
$`\overline{K}^0(k_0)K^{}(k_{},ϵ)|𝒪_{}|B^{}(p_B)`$ $`=`$
$`{\displaystyle \frac{m_K^2f_K}{(m_s+m_d)(m_s+m_b)}}(2m_Kϵ^{}k_0)A_0^{BK^{}}(m_K^2).`$ (21)
Taking the values of the quark masses as in $`m_b=4.88`$ $`GeV`$, $`m_s=122`$ $`MeV`$, $`m_d=7.6`$ $`MeV`$ and using the bound given in Eq. (5) we obtain the estimation of the upper limit of the branching ratio $`BR(B^{}K^{}\overline{K}^0)_{}`$ to be $`4.4\times 10^8`$. This limit can be raised to $`1.5\times 10^6`$ for the upper bound on the couplings of $`5.9\times 10^4`$ given in .
The two Higgs doublet model (6) gives for the amplitude of this decay
$`𝒜_{THDM}(B^{}(p_B)\overline{K}^0(k_0)K^{}(k_{},ϵ)`$ $`=`$
$`{\displaystyle \frac{i}{2}}{\displaystyle \frac{\xi _{sb}\xi _{sd}}{m_A^2}}[2m_K^{}f_KA_0^{BK}(m_K^2)ϵ^{}k_0]{\displaystyle \frac{m_K^2f_K}{(m_s+m_d)(m_s+m_b)}},`$ (22)
which gives for the limit $`|\xi _{sb}\xi _{sd}|/m_H^2>10^{10}`$ $`GeV^2`$, a branching ratio of the order $`10^{11}`$. Due to specific combination of the products of the scalar (pseudoscalar) densities this is the only decay which has nonvanishing amplitude within the factorization assumption.
c) $`B^{}K^{}\overline{K}^0`$ decay
For this decay mode the matrix element of the operator $`𝒪`$ is determined to be
$`\overline{K}^0(k_0,ϵ))K^{}(k_{})|𝒪|B^{}(p_B)`$ $`=`$ $`2g_K^{}f_KF_1^{BK}(m_K^2)ϵ^{}k_{}`$ (23)
giving the branching ratio in MSSM with an upper limit
$$BR(B^{}K^{}\overline{K}^0)_{MSSM}=5.9\times 10^9$$
(24)
in comparison with SM result $`6.5\times 10^{14}`$.
The amplitude calculated in MSSM including $``$ breaking and THDM vanishes, due to vanishing of the matrix element of the density operator for $`\overline{K}^0`$ state.
d) $`B^{}K^{}\overline{K}^0`$ decay
The matrix element of the operator $`𝒪`$ becomes in this case
$`\overline{K}^0(k_0)K^{}(k_{})|𝒪|B^{}(p_B)`$ $`=`$ $`if_KF_0^{BK}(m_K^2)(m_B^2m_K^2).`$ (25)
The multiplication with the corresponding $`1/4\stackrel{~}{C}_{MSSM}`$ gives the required amplitude $`\stackrel{~}{𝒜}`$. The branching ratio is then
$`\mathrm{\Gamma }(B^{}K^{}\overline{K}^0)`$ $`=`$ $`{\displaystyle \frac{1}{16\pi m_B^2}}\sqrt{m_B^24m_K^2}|\stackrel{~}{𝒜}|^2,`$ (26)
The branching ratio for MSSM is found to be $`BR(B^{}K^{}\overline{K}^0)_{MSSM}2.3\times 10^9`$, in comparison with the $`2.5\times 10^{14}`$ found in the SM. The matrix element of the $`R`$ parity breaking MSSM operator $`𝒪^{(1)}=(\overline{s}\gamma _5d)`$$`(\overline{s}b)`$ is found to be
$`K^{}\overline{K}^0|𝒪^{(1)}|B^{}`$ $`=`$ $`\overline{K}^0|\overline{s}\gamma _5d|0K^{}|\overline{s}b|B^{}`$
$`=`$ $`i{\displaystyle \frac{m_K^2}{(m_s+m_d)(m_b+m_s)}}f_KF_0^{BK}(m_K^2)(m_B^2m_K^2)`$
while the operator $`(\overline{s}\gamma _5b)(\overline{s}d)`$ gives the same result with the opposite sign. The decay width is then
$`\mathrm{\Gamma }(B^{}K^{}\overline{K}^0)_{}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi m_B^2}}\sqrt{m_B^24m_K^2}|K^{}\overline{K}^0|𝒪^{(1)}|B^{}/4|^2`$ (28)
$`\times `$ $`{\displaystyle \frac{f_{QCD}^2}{m_{\stackrel{~}{\nu }}^4}}({\displaystyle \underset{i=n}{}}|\lambda _{n32}^{}\lambda _{n21}^{}|^2+|\lambda _{n21}^{}\lambda _{n32}^{}|^2)`$
The constraint in (5) gives the bound $`9.4\times 10^8`$, while for the bound of $`5.9\times 10^4`$ for the coupling constants (6) the rate $`BR(B^{}K^{}\overline{K}^0)_{}`$ can reach $`3.3\times 10^6`$.
The long distance effects are usually suppressed in the $`B`$ meson decays. One might wonder if they are important in decays we consider here. We have estimated the tree level contribution of the $`D(D^{})`$ which then goes into $`K(K^{})`$ via weak annihilation. We found that these contributions give a branching ratio of the order $`10^{18}`$ and therefore they can be safely neglected. One might think that the exchange of two intermediate states $`D(D^{})`$, $`K(K^{})`$ can introduce certain long distance contributions. In decay $`B\mathrm{"}D\mathrm{"}\mathrm{"}K\mathrm{"}\mathrm{"}K\mathrm{"}\mathrm{"}K\mathrm{"}`$ the first weak vertex arises from the decay $`B\mathrm{"}D\mathrm{"}\mathrm{"}K\mathrm{"}`$ and the second weak vertex (see e.g. ) can be generally obtained from the three body decays of $`DKKK`$. In Ref. it was found that such contributions are also very small. Therefore, we are quite confident to suggest that the long distance effects are not important in the two body $`\mathrm{\Delta }S=2`$ $`B`$ decays.
Let us turn now to the possibility of detecting these decay modes. The $`B^{}K^{}\overline{K}^0`$ and $`B^{}K^{}\overline{K}^0`$ modes have clean signatures of a $`\mathrm{\Delta }S=2`$ transition and therefore these are the channels we recommend to look for. The other two modes we discussed, $`b)`$ and $`d)`$ have a $`\overline{K}^0`$ in the final states which complicates the possibilty of a detection because of $`K^0`$ $`\overline{K}^0`$ mixing. Separating the desired amplitude requires the measurement of the decays of both $`K_S`$ and $`K_L`$, since one can express
$`{\displaystyle \frac{\mathrm{\Gamma }(B^{}K^{}K_S)\mathrm{\Gamma }(B^{}K^{}K_L)}{\mathrm{\Gamma }(B^{}K^{}K_S)+\mathrm{\Gamma }(B^{}K^{}K_L)}}`$ $`=`$ $`Re\eta (B^{}KK^{}),`$ (29)
where
$`Re\eta (B^{}KK^{})`$ $`=`$ $`{\displaystyle \frac{A(B\overline{K}^0K^{})}{A(BK^0K^{})}}.`$ (30)
We summarize our results in the Table 1. The MSSM gives rates of the order $`10^910^8`$, while the $``$ parity breaking terms in the MSSM can be seen only in the $`B^{}K^{}\overline{K}^0`$ and $`B^{}K^{}\overline{K}^0`$ decay. These are the modes which as we mentioned are more difficult on the experimental side. The THDM model can give nonvanishing contribution only in the case of $`B^{}K^{}\overline{K}^0`$ decay, with a rate too small to be seen. Thus, we conclude by stressing the possibilty of detecting physics beyond SM mainly in the $`K^{}\overline{K}^0`$, $`K^{}\overline{K}^0`$ decays.
We thank Y. Rozen, S. Tarem and D. Zavrtanik for stimulating discussions on experimental aspects of this investigation.
This work has been supported in part by the Ministry of Science of the Republic of Slovenia (SF) and by the Fund for Promotion of Research at the Technion (PS).
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# Neutrinos produced by ultrahigh-energy photons at high red shift
\[
## Abstract
Some of the proposed explanations for the origin of ultrahigh-energy cosmic rays invoke new sources of energetic photons (e.g., topological defects, relic particles, etc.). At high red shift, when the cosmic microwave background has a higher temperature but the radio background is low, the ultrahigh-energy photons can generate neutrinos through pair-production of muons and pions. Slowly evolving sources produce a detectable diffuse background of $`10^{17}`$eV neutrinos. Rapidly evolving sources of photons can be ruled out based on the existing upper limit on the neutrino flux.
preprint: UCLA/00/TEP/22; BNL-HET-00/26
\]
Discovery of cosmic rays with energies beyond the Greisen-Zatsepin-Kuzmin (GZK) cutoff presents an outstanding puzzle in astrophysics and cosmology . Many proposed explanations invoke a new source, such as superheavy relic particles or topological defects , that can generate photons at both low and high red shifts. In understanding the origin of the ultrahigh-energy cosmic rays (UHECR), it is crucial to distinguish such sources from more conventional astrophysical ones . In this letter we show that a diffuse background of neutrinos with energies $`10^{17}`$eV can be generated by ultrahigh-energy photons at high red shift.
Generation of ultrahigh-energy neutrinos has been studied for various sources at small red shift, for which muon pair-production can be neglected. However, a substantial flux of neutrinos could be produced at earlier times, when the propagation of photons was different from that in the present universe because the intergalactic magnetic field was weaker, the density of radio background was lower, and the cosmic microwave background (CMB) density and temperature were higher.
At red shift $`z`$ the cosmic microwave background radiation (CMBR) has temperature $`T_{_{CMB}}(z)=2.7(1+z)\mathrm{K}`$. Because of this, at high red shift photon-photon and electron-photon interactions can produce pairs of muons and charged pions, whose decays generate neutrinos. This is in sharp contrast with the $`z\stackrel{<}{_{}}1`$ case, where the photons do not produce neutrinos as they lose energy mainly by scattering off the radio background through electron-positron pair production and subsequent electromagnetic cascade . The ratio of the CMBR density to that of universal radio background (RB) increases at higher $`z`$, and the processes $`\gamma \gamma _{_{CMB}}\mu ^+\mu ^{}`$ and $`e\gamma _{_{CMB}}e\mu ^+\mu ^{}`$ can produce muons, which decay into neutrinos: $`\mu e\nu _e\nu _\mu `$. The threshold for these interactions is $`\sqrt{s}>2m_\mu =0.21`$GeV, or
$$E_{\gamma ,e}>E_{\mathrm{th}}(z)=\frac{10^{20}\mathrm{eV}}{1+z}$$
(1)
Our discussion applies to any source of photons active at high red shift. The latter requirement excludes some astrophysical sources . Topological defects and decaying relic particles , however, could operate even at $`z1`$. These sources are expected to produce photons with energies as high as $`10^{20}`$eV. We will describe a neutrino signature of this class of sources.
At $`z<1`$ the main source of energy loss for photons is electromagnetic cascade that involves $`e^+e^{}`$ pair production (PP) on the radio background photons. The radio background is generated by normal and radio galaxies. Its present density is higher than that of CMB photons in the same energy range. The radio background determines the mean interaction length for the $`e^+e^{}`$ pair production. At red shift $`z`$, however, the density of CMB photons is higher by a factor $`(1+z)^3`$, while the density of radio background is either constant or, more likely, lower. Some models of cosmological evolution of radio sources predict a sharp drop in the density of radio background at red shift $`z\stackrel{>}{_{}}2`$. More recent observations indicate that the decrease of radio background at $`z>2`$ is slow. However, one expects the CMB to become a more important source of energy losses for photons at higher $`z`$ because of the $`(1+z)^3`$ increase in the density of CMB photons. Let $`z__R`$ be the value of red shift at which the scattering of high-energy photons off CMBR dominates over their scattering off RB. Based on the analyses of Refs. , we take $`z__R5`$. Another source of energy losses in the electromagnetic cascade is the synchrotron radiation by the electrons in the intergalactic magnetic field (IGMF). This is an important effect for red shift $`z<z__M`$, where $`z__M5`$ corresponds to the time when the synchrotron losses are not as significant as the interactions with the CMB radiation. We will use the value $`z_{\mathrm{min}}=\mathrm{max}(z__R,z__M)5`$ in what follows. As discussed below, a higher value of $`z_{\mathrm{min}}`$, even as high as 10, would not make a big difference in the flux of the signature neutrinos.
Let us now consider the propagation of UHE photons at $`z>z_{\mathrm{min}}`$. In particular, we are interested in neutrino-generating processes, that is, reactions that produce muons and pions. For photon energies above the threshold for muon pair production (1), the reactions $`\gamma \gamma _{_{CMB}}e^+e^{}`$, $`\gamma \gamma _{_{CMB}}e^+e^{}e^+e^{}`$ and $`\gamma \gamma _{_{CMB}}\mu ^+\mu ^{}`$ are possible. For $`\sqrt{s}>2m_{\pi ^\pm }=0.28`$GeV the charged pion production may also occur. Among the processes listed above, the electron pair production (PP) has the highest cross section for photon energies $`E_\gamma \stackrel{<}{_{}}5\times 10^{20}\mathrm{eV}/(1+z)`$. Since the energies of the two interacting photons are vastly different, either the electron or the positron from PP has energy close to that of the initial photon. At higher photon energies, double pair production (DPP) becomes more important . Four electrons, each carrying about $`1/4`$ of the initial photon energy, are produced in this reaction. Thus, after an initial $`\gamma \gamma _{_{CMB}}`$ interaction one ends up with one or more UHE electrons.
These electrons continue to scatter off CMBR. At lower energies, inverse Compton scattering (ICS), $`e\gamma _{_{CMB}}e\gamma `$, converts high-energy electrons into high-energy photons . However, at energies above the muon threshold, higher order processes, such as triplet production (TPP) $`e\gamma _{_{CMB}}ee^+e^{}`$ and muon electron-pair production (MPP) $`e\gamma _{_{CMB}}e\mu ^+\mu ^{}`$, dominate. For center of mass energies $`sm_e^2`$, the inelasticity $`\eta `$ for TPP is very small: $`\eta 1.768(s/m_e^2)^{3/4}<10^3`$ . One of the electrons produced through TPP, carries almost all ($`1\eta `$) of the incoming electron’s energy. It can interact once again with the CMBR. As a result, the leading electron can scatter many times before losing a considerable amount of energy. Hence, the energy attenuation length $`\lambda _{\mathrm{eff}}`$ is much greater than the TPP interaction length: $`\lambda _{_{\mathrm{TPP}}}\eta \lambda _{\mathrm{eff}}`$.
To see if neutrinos are produced, one must compare this energy attenuation length with the interaction length for muon pair production in processes like $`e\gamma _{_{CMB}}e\mu ^+\mu ^{}`$. Above the pion threshold, pion production is yet another channel that drains the energy out of the electromagnetic cascade and into neutrinos. We note that even a single neutrino-producing channel is enough for UHE photons to produce neutrinos at high red shift. The fact that there are several such channels makes little difference.
Let us compare the TPP energy attenuation length $`\lambda _{\mathrm{eff}}`$ with the interaction length for muon pair production. The interaction length is given by $`\lambda ^1n_{_{CMB}}v\sigma `$, and thus the ratio is $`R=\lambda _{\mathrm{eff}}/\lambda _{_{\mathrm{MPP}}}\sigma _{_{\mathrm{MPP}}}/(\eta \sigma _{_{\mathrm{TPP}}})`$. For $`sm_e^2`$ the cross section for TPP is
$$\sigma _{_{TPP}}\frac{3\alpha }{8\pi }\sigma __T\left(\frac{28}{9}\mathrm{ln}\frac{s}{m_e^2}\frac{218}{27}\right),$$
(2)
where $`\sigma __\mathrm{T}`$ is the Thompson cross section. The MPP cross section in the energy range just above the threshold $`5m_\mu ^2<s<20m_\mu ^2`$ is of the order of $`0.11\mathrm{m}\mathrm{b}`$, and the ratio $`R100`$.
Since $`\lambda _{\mathrm{eff}}\lambda _{_{\mathrm{MPP}}}`$, in the absence of dense radio background and intergalactic magnetic fields, all electrons with $`E>E_{\mathrm{th}}`$ pair-produce muons before their energy is reduced by the cascade. For muon production close to the threshold, each muon carries on average $`1/4`$ of the incoming electron’s energy . Muons decay before they can interact with the photon background. Each energetic muon produces two neutrinos and an electron. The electron produced alongside the muon pair gets half or more of the incoming electron’s energy; it can interact again with the CMBR to produce muons. This process can repeat until the energy of the regenerated electron decreases below the threshold for muon pair production. Higher energy electrons with energies $`E>2\times 10^{20}\mathrm{eV}/(1+z)`$eV can also produce pions through the reaction $`e\gamma _{_{CMB}}e\pi \pi `$. Charged pions decay into neutrinos, while neutral pions reproduce photons. As explained above, it makes little difference through which channel the neutrinos are produced – as long as there is at least one reaction with a shorter mean free path than the energy attenuation length.
One can parameterize the rate of photon production as $`\dot{n}__X=\dot{n}_{{}_{\gamma }{}^{},0}(t/t_0)^m`$, with $`m=0`$ for decaying relic particles, $`m=3`$ for ordinary string and necklaces, and $`m4`$ for superconducting strings . Let $`z_{\mathrm{max}}`$ be the red shift at which the universe becomes opaque to ultrahigh-energy neutrinos. Its value is determined by the neutrino interactions with the relic neutrino background. The absorption red shift for neutrinos with energy $`10^{17}`$eV is $`z_{\mathrm{max}}3\times 10^3`$ . All neutrinos coming from red shift $`z_{\mathrm{min}}<z<z_{\mathrm{max}}`$ contribute to the present flux. The neutrino flux is
$`n_\nu `$ $`=`$ $`\xi {\displaystyle _{z_{\mathrm{min}}}^{z_{\mathrm{max}}}}𝑑t\dot{n}_\gamma (z)(1+z)^4`$ (3)
$`=`$ $`\xi {\displaystyle \frac{3}{2a}}\dot{n}_{\gamma ,0}t_0[(1+z_{\mathrm{min}})^a(1+z_{\mathrm{max}})^a],`$ (4)
where $`\xi `$ is the number of neutrinos produced per UHE photon, and $`a=(3m11)/2`$. We take $`\xi 4`$ because one photon produces one UHE electron, which generates a pair of muons, which decay into four neutrinos. This is probably an underestimate because the remaining electron may have enough energy for a second round of muon pair-production. Also, pion decays produce three neutrinos each.
For $`m<11/3`$, $`a<0`$, and, according to eq. (4), most of neutrinos come from red shift $`zz_{\mathrm{min}}5`$. All these neutrinos are produced by photons with energies $`E_\gamma >E_{\mathrm{min}}=10^{20}\mathrm{eV}/(1+z_{\mathrm{min}})2\times 10^{19}`$eV. If decaying TD’s or relic particles are the origin of the UHECR today, one can use the observed UHECR flux to fix the overall normalization constant $`\dot{n}_{\gamma ,0}(E>E_{\mathrm{min}})`$. We will use the photon fluxes calculated in . For various sources, $`\dot{n}_{\gamma ,0}(E>E_{\mathrm{min}})L^1_{E_{\mathrm{min}}}dEJ(E)`$, where is $`J(E)`$ the differential photon flux, and $`L`$ the length scale from which the photons are collected. Because the photon flux is a sharply falling function of energy, $`\dot{n}_{\gamma ,0}`$ is dominated by photons with energies $`EE_{\mathrm{min}}`$. Relic particles $`(m=0)`$ and monopolonia $`(m=3)`$ cluster in galaxies and have $`LL_{\mathrm{gal}}100\mathrm{k}\mathrm{p}\mathrm{c}`$, the size of our galaxy. Their over-density in the galaxy is $`2\times 10^5`$ . Necklaces $`(m=3)`$ on the other hand are distributed uniformly throughout the universe; for them $`L=L_\gamma 5`$Mpc, the photon absorption length at these energies.
Using the photon flux from Refs. one obtains the neutrino flux:
$$\varphi _\nu \{\begin{array}{ccc}10^{21}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1,\hfill & \mathrm{relic}\mathrm{particles}(m=0),\hfill & \\ 10^{18}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1,\hfill & \mathrm{monopolonium}(m=3),\hfill & \\ 10^{16}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1,\hfill & \mathrm{necklaces}(m=3).\hfill & \end{array}$$
(5)
Taking $`z_{\mathrm{min}}=10`$ reduces the flux of signature neutrinos from $`m=3`$ sources only by a factor $`2`$.
The energy of these neutrinos at red shift $`z`$ is $`E_\nu (z)E_\mu /3`$. It is then further red shifted by a factor $`(1+z)^1`$. Assuming a falling photon spectrum, we expect most of neutrinos to come from photons near the threshold, eq. (1). We estimate the energy of these neutrinos after the red shift $`E_\nu 10^{17}`$eV.
If the source in question is a slowly decaying relic particle or some other source with $`m=0`$, the neutrinos produced at high red shift are probably not detectable. Of course, if the relic particles or topological defects produce UHECR through $`Z`$-bursts , neutrinos from $`Z`$ decays can be detected.
Can other sources produce a comparable flux of neutrinos at $`10^{17}10^{18}`$eV? The neutrino flux $`\varphi _\nu 10^{16}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1`$ at $`E_\nu 10^{17}`$eV exceeds the background flux from the atmosphere and from pion photoproduction on CMBR at this energy , as well as the fluxes predicted by a number of models . TD can produce a large flux of primary neutrinos. However, the primary flux peaks at $`E_\nu 10^{20}`$eV, while the secondary flux peaks at $`E_\nu 10^{17}`$eV and creates a distinctive “bump” in the spectrum. Models of active galactic nuclei (AGN) have predicted a similar flux of neutrinos at these energies . The predictions of these models have been a subject of debate . However, every one agrees that AGN cannot produce neutrinos with energies of $`10^{20}`$eV . So, an observation of $`10^{17}`$eV neutrinos accompanied by a comparable flux of $`10^{20}`$eV neutrinos would be a signature of a TD rather than an AGN.
There is yet another interesting possibility. TD with $`m4`$, e.g., superconducting strings, cannot give a large enough flux of UHECR because of the EGRET bound on the flux of $`\gamma `$-photons . However, this does not mean they did not exist in the early universe. Neutrinos with energy $`10^{17}`$eV are probably the only observable signature of some rapidly evolving sources that could be active at high red shift but would have “burned out” by now.
To summarize, we have shown that sources of ultrahigh-energy photons that operate at red shift $`z\stackrel{>}{_{}}5`$ produce neutrinos with energy $`E_\nu 10^{17}`$eV. The flux depends on the evolution index $`m`$ of the source. A distinctive characteristic of this type of neutrino background is a cutoff below $`10^{17}`$eV due to the universal radio background at $`z<z_{\mathrm{min}}`$. Detection of these neutrinos can help understand the origin of ultrahigh-energy cosmic rays.
We thank J. Alvarez-Muniz, P. Biermann, V. Berezinsky, F. Halzen, and G. Sigl for helpful comments. This work was supported in part by the US Department of Energy grant DE-FG03-91ER40662, Task C, as well as by an Assistant Professor Initiative grant from UCLA Council on Research. A.K. thanks CERN Theory Division for hospitality during his stay at CERN when part of this work was performed.
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# Chandra constraints on the thermal conduction in the intracluster plasma of A2142
## 1 INTRODUCTION
The heat stored in the intracluster plasma is conducted down any temperature gradient present in the gas in a way that can be described through the following equations (Spitzer 1956, Sarazin 1988):
$$q=\kappa \frac{d(kT_\mathrm{e})}{dr},$$
(1)
where $`q`$ is the heat flux, $`T_\mathrm{e}`$ is the electron temperature, and $`\kappa `$ is the thermal conductivity that can be expressed in term of the density, $`n_\mathrm{e}`$, the electron mass, $`m_\mathrm{e}`$, and the electron mean free path, $`\lambda _\mathrm{e}`$, as (Cowie and McKee 1977)
$$\kappa =1.31n_\mathrm{e}\lambda _\mathrm{e}\left(\frac{kT_\mathrm{e}}{m_\mathrm{e}}\right)^{1/2}.$$
(2)
In a fully ionized gas of (mostly) hydrogen, the electron mean free path is function of the gas temperature, density and impact parameters of the Coulomb collisions, $`\mathrm{\Lambda }`$:
$$\lambda _\mathrm{e}=30.2\mathrm{T}^2n^1\left(\frac{\mathrm{ln}\mathrm{\Lambda }}{37.9+\mathrm{ln}\left(\mathrm{T}/n^{1/2}\right)}\right)^1\text{kpc},$$
(3)
where we have adopted the following dimensionless quantities:
$$\mathrm{T}=\left(\frac{kT_\mathrm{e}}{\text{10 keV}}\right),n=\left(\frac{n_\mathrm{e}}{\text{10}\text{-3}\text{ cm}\text{-3}}\right)$$
(4)
Using this expression and the adopted typical values for cluster plasma, we can then write
$$\kappa =8.2\times 10^{20}\mathrm{T}^{5/2}\text{erg s}\text{-1}\text{ cm}\text{-1}\text{ keV}\text{-1}.$$
(5)
If the mean electron free path is comparable to the scale length $`\delta r`$ of the temperature gradient, the heat flux tends to saturate to the limiting value which may be carried by the electrons (Cowie and McKee 1977):
$`q_{\mathrm{sat}}`$ $`=`$ $`0.42\left({\displaystyle \frac{2kT_\mathrm{e}}{\pi m_\mathrm{e}}}\right)^{1/2}n_\mathrm{e}kT_\mathrm{e}`$ (6)
$`=`$ $`0.023\mathrm{T}^{3/2}n\text{erg s}\text{-1}\text{ cm}\text{-2},`$
where the factor 0.42 comes from the reduction effect on the heat conducted by the electrons and that is produced from the secondary electric field that maintain the total electric current along the temperature gradient at zero (Spitzer 1956).
In this Letter, we apply these equations to estimate the efficiency of thermal conductivity in the intracluster medium of A2142, a cluster of galaxies observed by the X-ray telescope Chandra during its calibration phase in August 1999.
## 2 THERMAL CONDUCTIVITY IN A2142
Markevitch et al. (2000) have analyzed the Chandra observation of the merging cluster of galaxies A2142 and made an important discovery. The X-ray image reveals of sharp edges to the surface brightness of the central elliptical-shaped region. The edges are located about 3 arcmin to the northwest and 1 arcmin to the south with respect to the X–ray centre. Markevitch et al (2000) show that the bright and fainter regions either side of an edge are in pressure equilibrium with each other, but with a dramatic electron temperature decrease on the inside.
Considering the values of the intracluster gas properties in A2142 from their Figure 4, together with the equations presented in our Introduction, we can estimate whether thermal conductivity is efficient in erasing the observed temperature gradient.
The electron temperature (panel $`b`$ in their Fig. 4) varies from 5.8 to 10.6 keV on either side of the boundary of the southern edge of the central bright patch in A2142, and from 7.5 to 13.8 keV at the northern edge. The relative uncertainties on these values are about 20 per cent at the 90 per cent confidence level. The electron density (panel $`d`$ in their Fig. 4) at the edges is $``$ 1.2 $`\times `$ 10<sup>-2</sup> cm<sup>-3</sup>, 3.0 $`\times `$ 10<sup>-3</sup> cm<sup>-3</sup> to the South and North, respectively.
The scale length $`\delta r`$ on which this temperature gradient is observed is spatially unresolved in the temperature profile and appears enclosed between 0 and 35 kpc ($`H_0=`$ 70 km s<sup>-1</sup> Mpc<sup>-1</sup>) for the Southern edge, 0–70 kpc for the Northern edge. However, the surface brightness profiles (panel $`c`$) show a radially discontinuous derivative at the positions of the sharp edges, on scales of about 10–15 kpc. We adopt hereafter these values as representative of $`\delta r`$, using the larger values only for upper limit purposes.
The case for a saturated flux is reached when these scales are comparable with the electron mean free path of about 2 and 12 kpc for the Southern and Northern edge, respectively, calculated using the temperature and density estimates in eqn. 3.
Therefore, we will consider hereafter the two extreme cases where (i) $`\delta r\lambda _\mathrm{e}`$ and the heat flux is un-saturated, (ii) $`\delta r\lambda _\mathrm{e}`$ and the heat flux is saturated and represented by eqn. 6.
The maximum heat flux in a plasma is given by
$$q=\frac{3}{2}n_\mathrm{e}kT_\mathrm{e}\overline{v},$$
(7)
where $`\overline{v}=dr/d\tau `$ is a characteristic velocity that we are now able to constrain equalizing the latter equation to eqn. 1.
In particular, given the observed values (and relative errors) of density and temperature across the two edges and $`\delta r`$ of (i: non-saturated flux) 10 and 20 kpc (upper limits: 35 and 70 kpc for the edges to South and North, respectively) and (ii: saturated flux) $`\lambda _\mathrm{e}`$, the characteristic time, $`\delta \tau `$, required to erase the electron temperature gradient and due to the action of the thermal conduction alone would be:
$$\delta \tau =\frac{\delta r}{\overline{v}}=\{\begin{array}{c}3.6(<80)\times 10^6\text{yrs},\delta r2\text{kpc}\hfill \\ 0.3(<0.4)\times 10^6\text{yrs},\delta r2\text{kpc}\hfill \end{array}$$
(8)
for the Southern edge, and
$$\delta \tau =\frac{\delta r}{\overline{v}}=\{\begin{array}{c}2.4(<52)\times 10^6\text{yrs},\delta r12\text{kpc}\hfill \\ 1.9(<2.4)\times 10^6\text{yrs},\delta r12\text{kpc}\hfill \end{array}$$
(9)
for the Northern edge. The upper limits are obtained propagating the uncertainties on the temperature and, for the “$`\delta r\lambda _\mathrm{e}`$” condition only, assuming the spatial resolution of the temperature profile as indicative of the length of the gradient. Here we note that the limit on the timescale for saturated flux is the minimum value given the condition of the gas. Any value of $`\delta \tau `$ estimated on scales considerably larger than the electron mean free path has to be longer than the limit for saturated flux (also significantly, given that the timescale is proportional to $`\delta r^2`$).
When this time interval is compared with the core crossing time of the interacting clumps of about 10<sup>9</sup> yrs, we conclude that thermal conduction needs to be suppressed by a factor larger than 10 and with a minimum characteristic value enclosed between 250 and 2500. On the other side, Markevitch et al. suggest a dynamical model in which the dense cores of the two interacting clumps are moving through the host, less dense, intracluster medium at a subsonic velocity of less than 1000 km s<sup>-1</sup> and 400 km s<sup>-1</sup>, for the Northern and Southern edge, respectively, leading to a timescale of about $`2\times 10^7`$ yrs or larger for the cool and hot phases to be in contact. In this scenario, our results implies that the conduction is suppressed at least by a factor of 2–200 in the South and 2–32 in the North. However, we note that it is unlikely that the cooler gas can have arrived at its present arrangement and settled in the hotter enviroment on a timescale much shorter than a core crossing time. The frequency of the occurance of similar structures in other cluster cores will be important in establishing the timescale for their formation and duration and, hence, improve the constraint on the thermal conduction in the intracluster plasma.
## 3 CONCLUSIONS
In this letter, we calculate the thermal conductivity in the intracluster medium of A2142, a interacting cluster of galaxies observed by Chandra during the calibration phase.
We have shown that the time interval in which the action of thermal conduction should propagate heat to neighbouring regions is shorter by a factor of about 250–2500 than the likely estimated age of the structure. We note that the observed sharp temperature boundary also means that mixing and diffusion are minimal.
The results presented here are a direct measurement of a physical process in the intracluster plasma and imply that thermal conduction is particularly inefficient within 280 $`h_{70}^1`$ kpc of the central core. The gas in the central regions of many clusters has a cooling time lower than the overall age of the system, so that a slow flow of hotter plasma moves here from the outer parts to maintain hydrostatic equilibrium. In such a cooling flow (Fabian, Nulsen, Canizares 1991; Fabian 1994), several phases of the gas (i.e. with different temperatures and densities) are in equilibrium and would thermalize if the conduction time were short. The large suppression of plasma conductivity in the cluster core allows an inhomogeneous, multi–phase cooling flow to form and be maintained, as is found from spatial and spectral X-ray analyses of many clusters (e.g. Allen et al. 2000).
How the conduction is reduced by a so large factor is still unclear. Binney & Cowie (1981) explain the reservoir for heat observed in the region of M87 as requiring an rms field strength considerably larger than the component of the field parallel to the direction along which conduction occurs. (The transport processes are reduced in the direction perpendicular to magnetic field lines). This would imply either highly tangled magnetic fields or large–scale fields perpendicular to the lines connecting the hotter to the cooler zones. Such fields could become dynamically important. Chandran, Cowley & Albright (1999) use asymptotic analysis and Monte Carlo particle simulations to show that tangled field lines and, with larger uncertainties, magnetic mirrors reduce the Spitzer conductivity by a large factor. Via a phenomenological approach Tribble (1989) argued that a multiphase intracluster medium is an inevitable consequence of the effect of a tangled magnetic field on the flow of the heat through the cluster plasma. Electromagnetic instabilities driven by the temperature gradient (e.g. Pistinner, Levinson & Eichler 1996) can represent another possible explanation for the suppression of thermal conductivity. Finally, we speculate that cooler gas dumped in the cluster core by a merger (see e.g. Fabian & Daines 1991) would be part of a different magnetic structure to the hotter gas and so thermally isolated.
## ACKNOWLEDGEMENTS
We acknowledge the support of the Royal Society.
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# 1 Introduction
## 1 Introduction
If in the future $`e^{}e^+`$ Linear Colliders (LC) , the option to develop high energy $`\gamma \gamma `$ collisions will also be available, then many new opportunities for new physics (NP) searches should arise. Employing back-scattering of laser photons, this option transforms an<sup>1</sup><sup>1</sup>1In this case it would be best to run LC in its $`e^{}e^{}`$ mode, . LC to essentially a $`\gamma \gamma `$ Collider ($`\mathrm{LC}_{\gamma \gamma }`$) with about $`80\%`$ of the initial energy and a comparable luminosity . The importance of $`\mathrm{LC}_{\gamma \gamma }`$ stems from the fact that the cross sections for gauge boson and top production in $`\gamma \gamma `$ collisions at sufficiently high energies, are often considerably larger than the corresponding quantities in the $`e^{}e^+`$ case .
To some extent, such an enhancement should arise for Higgs production also. For the neutral Higgs particles in particular, an $`\mathrm{LC}_{\gamma \gamma }`$ may act as a Higgs factory which can be used to study their detail properties, including possible Higgs anomalous couplings . Since the anomalous gauge boson, top and Higgs couplings are interconnected and constitute an important possible source of new physics, an $`\mathrm{LC}_{\gamma \gamma }`$ should be very helpful for its identification. In case the NP scale is very high, such forms of NP may be described by the complete list of $`dim=6`$ operators involving gauge bosons and/or quarks of the third family presented in .
Alternatively, it may turn out that the NP scale is nearby, as it would be expected in the usual SUSY scenario . In such a case many neutral spinless particles of Higgs and sneutrino type may exist, and an $`\mathrm{LC}_{\gamma \gamma }`$ may be used for an s-channel production of the CP-even light and heavy neutral Higgs bosons $`h^0`$ and $`H^0`$ respectively, as well as the CP-odd $`A^0`$. The study of the various branching ratios, and the polarization of the incoming photons, could then be very helpful to establish and disentangle the nature of these Higgs particles .
Once any of these spinless bosons is discovered, its properties should be carefully looked at, in order to be sure that they fulfill the SUSY expectations. Motivated by this, we study in this paper the process $`\gamma \gamma A^0A^0`$ in the context of a minimal SUSY model, where no new sources of CP violation, apart from those already known in the Standard Model (SM) Yukawa potential, are assumed to exist. Thus, the various new SUSY couplings are taken to be real, but no specific assumption on their relative magnitudes or signs is made . As we will see below, in such a case, there are only two independent helicity amplitudes for $`\gamma \gamma A^0A^0`$ , denoted below as $`F_{++}`$ and $`F_+`$, where the indices describe the helicities of the incoming photons.
It is also interesting to study the phases of these amplitudes. The motivation for this stems from the recent observation in , that at c.m energies $`250GeV`$, out of the many independent helicity amplitudes for the processes $`(\gamma \gamma \gamma \gamma ,\gamma Z,ZZ)`$, only the two helicity conserving amplitudes $`F_{++++}`$ and $`F_{++}`$ are important, which moreover turn out to be almost purely imaginary<sup>2</sup><sup>2</sup>2For $`\gamma \gamma ZZ`$ in SM the further assumption is made that the standard Higgs particle is light; e.g. below $`200GeV`$.. The physical reason for this result is not very clear . Therefore, it seems worthwhile to investigate what happens in other processes, like e.g. the neutral Higgs boson production, which, as the neutral gauge bosn production, also vanish at tree order and they first appear at the 1-loop level.
Below, in Section 2 we give an overall view of the $`\gamma \gamma A^0A^0`$ helicity amplitudes in SUSY. The needed SUSY vertices appear in Appendix A, while the corresponding contributions to the amplitudes are given in Appendix B. The results are expressed in terms of $`C_0`$ and $`D_0`$ Passarino-Veltman functions only , using expressions analogous to those encountered in the $`\gamma \gamma ZZ`$ calculation . Finally in Section 3, we give our Conclusions.
Coming now to the related studies already existing in the literature, we first remark that $`\gamma \gamma h^0h^0`$ has been studied in SM by Jikia . In the non-linear gauge defined in (A.1) and used here, the only contributing diagrams involve $`W`$ or top-loops, similar to those appearing in Figs.1,3. We have repeated the calculations of and agree with the results, apart from the overall sign of the<sup>3</sup><sup>3</sup>3For the gauge boson polarization vectors, here and in , we use the same conventions as in . The only difference is that we use the JW convention , which introduces an additional minus to the polarization vector of a longitudinal ”Number 2” Z and affects $`\gamma \gamma ZZ,\gamma Z`$. $`F_{++}`$ amplitude. For the top contributions, our results are fully consistent with those of . The relevant amplitudes are presented and compared to those of $`\gamma \gamma A^0A^0`$ at the end of Section 2.
In a calculation of $`\gamma \gamma h^0h^0`$ in a general SUSY model has been presented in terms of the general $`C_j`$ and $`D_j`$ of Passarino and Veltman functions. The production of two neutral Higgs pairs in SUSY models at hadronic Colliders has also been studied in ; where of course the complications from loops involving $`W`$-bosons, or ”single” and ”mixed” charginos, are avoided. Finally the processes $`\gamma \gamma H^0H^0`$ and $`\gamma \gamma A^0A^0`$ have also appeared in a non-Supersymmetric gauge model involving a two Higgs doublet scalar sector .
## 2 An overall view of the $`\gamma \gamma A^0A^0`$ amplitudes.
The invariant helicity amplitudes for<sup>4</sup><sup>4</sup>4We use the same conventions as in . $`\gamma \gamma A^0A^0`$ are denoted as $`F_{\lambda _1,\lambda _2}(\widehat{s},\widehat{t},\widehat{u})`$, where $`\lambda _j`$ describe the helicities of the incoming photons, and the kinematics are defined in Appendix B. Assuming that the SUSY Higgs potential is CP-invariant we get (see (B.2))
$$F_{\lambda _1,\lambda _2}(\widehat{s},\widehat{t},\widehat{u})=F_{\lambda _1,\lambda _2}(\widehat{s},\widehat{t},\widehat{u}),$$
(1)
which implies that there are only two independent helicity amplitudes, $`F_{++}(\widehat{s},\widehat{t},\widehat{u})`$ and $`F_+(\widehat{s},\widehat{t},\widehat{u})`$.
As in , we employ the non-linear gauge of , which implies the gauge fixing and FP-ghost interactions of (A.1, A.2), leading to the conclusion that there are no $`\gamma W^\pm G^{}`$, $`ZW^\pm G^{}`$ vertices. The diagrams contributing to $`A^0`$-pair production are then given in Figs.1-4.
The contribution to the $`F_{++}`$ and $`F_+`$ amplitudes from the diagrams in Fig.1 consists of two types. The first is induced by the two diagrams in the first line in Fig.1 and describes the $`(h^0,H^0)`$-pole contributions appearing in (B.19, B.20). The diagrams in the second to last line of Fig.1 involve loops in which $`W^\pm `$ and/or $`H^\pm `$ are running along their internal lines. These induce the second type of contributions contained in (B.22, B.23), and expressed in terms of the $`(C_0,D_0)`$-functions explained in (B.8-B.14); as well as the functions $`\stackrel{~}{F}^{WH^\pm }`$, $`\stackrel{~}{F}^{H^\pm W}`$, $`E_1^{WH^\pm }`$ defined in (B.15, B.16). The contributions (B.22, B.23) give the largest effect to the $`\gamma \gamma A^0A^0`$ amplitudes, for the numerical applications considered below.
The chargino loop contribution is described by the diagrams in Fig.2. It consists also of an $`(h^0,H^0)`$-pole contribution given in (B.24); the box contributions involving a ”single chargino”-loop giving (B.26, B.27); and the ”mixed chargino” contribution (B.28, B.29), arising when both charginos are running along the loop. Analytically, the later is the most complicated one. Nevertheless, it is simple enough to be possible to write it. Numerically, it has to be taken into account only when both charginos are relatively light.
The $`t`$ and $`b`$ quark contributions are described by the diagrams in Fig.3. They are given in (B.31) for the $`(h^0,H^0)`$-pole contribution, and in (B.33, B.34) for the box diagrams.
As an example of a sfermion contribution, we only considered the one arising from the $`(\stackrel{~}{t}_1,\stackrel{~}{t}_2)`$-loop, described by the diagrams in Fig.4. Their contributions are given by (B.35-B.38).
For the numerical applications we use the three CP-invariant mSUGRA set of parameters introduced in and presented in Table 1. For the electromagnetic coupling we take $`\alpha =1/127.8`$. The results are shown in Figs.5-7.
The real and imaginary parts of the helicity amplitudes $`F_{++}(\gamma \gamma A^0A^0)`$ and $`F_+`$ are presented in Figs. 5-7 . As indicated there, the most important contributions to the amplitudes arise from the $`(W,H^\pm )`$-loop diagrams presented in the 2nd to last line of Fig.1 and appearing in (B.22, B.23). At sufficiently high energies, these contributions are mainly imaginary. But the predominance of the imaginary parts of the amplitudes is not so strong, as the one observed in the gauge boson production cases .
As indicated in Figs. 5-7, the chargino contribution is generally quite important; while the $`t,b`$-quark contribution is somewhat smaller; and the stop contribution is negligible for the above cases.
For the $`t,b`$-quark contribution we also remark that in the mSUGRA(1) and mSUGRA(3) cases, where $`\mathrm{tan}\beta `$ is small, the $`b`$-contribution is negligible compared to the top one. On the contrary, for the mSUGRA(2) case of $`\mathrm{tan}\beta =30`$, the $`b`$-quark contribution may be more important than the $`t`$-quark one.
For comparison, we have also looked at the $`F_{++}`$ and $`F_+`$ amplitudes for $`\gamma \gamma h^0h^0`$ in the Standard Model. The results for $`m_{h^0}=120\mathrm{G}\mathrm{e}\mathrm{V}`$ are given in Fig.8. For the $`F_{++}`$ we find that the $`top`$-loop contributions is comparable the $`W`$-one, and the amplitude is never particularly imaginary. It is only for $`F_+`$, for which there is no Higgs-pole contribution; that at energies $`600\mathrm{G}\mathrm{e}\mathrm{V}`$, the $`W`$-loop is more important than the top-one, and the imaginary part of the amplitude becomes predominant.
The $`\gamma \gamma A^0A^0`$ unpolarized cross section for the sets of parameters in Table 1, are given in Fig.9. It lies in the range of $`(0.10.2)\mathrm{fb}`$, which is similar but somewhat smaller, than the result expected for $`\sigma (\gamma \gamma h^0h^0)`$ in SM for $`m_{h^0}120\mathrm{G}\mathrm{e}\mathrm{V}`$ . This result does not seem particularly sensitive to SUSY parameters like e.g. $`\mathrm{tan}\beta `$; but mainly depends on the $`A^0`$-mass. It should also be compared to the situation for a single $`A^0`$ or $`H^0`$ production studied in . We also remark that a cross section at the $`(0.10.2)\mathrm{fb}`$-level may be observable, if a luminosity of e.g. $`_{\gamma \gamma }250\mathrm{f}\mathrm{b}^1/\mathrm{year}`$ is realized in TESLA .
## 3 Conclusions
The Higgs sector, which is responsible for giving masses to almost all particles immediately after our Universe started, is definitely the most fascinating part of the present elementary particle theory. Motivated by this and assuming that the SUSY option is chosen by nature, we have studied here the process $`\gamma \gamma A^0A^0`$.
In the non-linear gauge used here, the types of contributing diagrams may be divided into two categories constructed on the basis of whether an s-channel neutral Higgs-pole is involved or not. Each category may then be further divided into three classes, on the basis of whether their loops involve the $`(W,H^\pm )`$-pair, charginos or sfermions. General formulae have been presented which allow the description of the process in any SUSY model, minimal or non-minimal.
For the numerical applications we only considered three SUGRA examples presented in Table 1, leading to an $`A^0`$ heavier than $`250\mathrm{G}\mathrm{e}\mathrm{V}`$. Excluding the forward and backward regions, the $`\sigma (\gamma \gamma A^0A^0)`$ cross-section is found in the $`(0.10.2)\mathrm{fb}`$ region.
At sufficiently high energies, both amplitudes $`F_{++}`$ and $`F_+`$ are found to be to largely imaginary; an effect reminiscent, but not so predominant, as the one noticed in neutral gauge boson production . On the contrary, nothing like this appears for the $`F_{++}`$ amplitude in $`\gamma \gamma h^0h^0`$, in the Standard Model. It seems that the predominance of the imaginary part of a loop amplitude at high energies, is somehow associated with the predominance of a $`W`$-involving loop. The understanding of such properties may be useful for new physics searches; since e.g. for $`\gamma \gamma \gamma \gamma `$ they determine the way the interference between the ”old” and possible forms of ”new” physics may appear .
Thus, after the discovery of $`A^0`$ and the study of the single production process $`\gamma \gamma A^0`$ , the study of the double $`A^0`$ production through $`\gamma \gamma A^0A^0`$, should certainly be useful for verifying the Higgs identification.
Appendix A: The MSSM vertices for $`\gamma \gamma A^0A^0`$.
In order to reduce the number of diagrams contributing to $`\gamma \gamma A^0A^0`$, we use the nonlinear gauge defined by the gauge fixing term
$`_{GF}`$ $`=`$ $`{\displaystyle \frac{1}{\xi _W}}F^+F^{}{\displaystyle \frac{1}{2\xi _Z}}(F_Z)^2{\displaystyle \frac{1}{2\xi _\gamma }}(F_\gamma )^2,`$
$`F^\pm `$ $`=`$ $`^\mu W_\mu ^\pm \pm i\xi _Wm_WG^\pm \pm ig^{}B^\mu W_\mu ^\pm ,`$
$`F_Z`$ $`=`$ $`^\mu Z_\mu +\xi _Zm_ZG^0,`$
$`F_\gamma `$ $`=`$ $`^\mu A_\mu ,`$ (A.1)
which is free from $`\gamma W^\pm G^{}`$ and $`ZW^\pm G^{}`$ vertices . The implied ghost-photon and ghost-scalar field interactions then are
$`_{FP}`$ $`=`$ $`ieA_\mu (^\mu \overline{\eta }^{}\eta ^+^\mu \overline{\eta }^+\eta ^{}+\overline{\eta }^+^\mu \eta ^{}\overline{\eta }^{}^\mu \eta ^+)+e^2A_\mu A^\mu (\overline{\eta }^+\eta ^{}+\overline{\eta }^{}\eta ^+)`$ (A.2)
$``$ $`{\displaystyle \frac{1}{2}}\xi _Wgm_W(\overline{\eta }^+\eta ^{}+\overline{\eta }^{}\eta ^+)[\mathrm{cos}(\alpha \beta )H^0+\mathrm{sin}(\beta \alpha )h^0].`$
The complete list of diagrams contributing to the process $`\gamma \gamma A^0A^0`$ in the present gauge, appear in Figs.1-4. Below we give the interaction Lagrangian describing the vertices for these sets of diagrams.
The diagrams in Fig.1 describe the $`(H^\pm ,W^\pm )`$ loop contribution (together of course with the accompanying ghost and Goldstone ones). The relevant vertices involve, in addition to the gauge boson self-interactions present in SM, also the triple and quartic vertices
$`_{VH}`$ $`=`$ $`{\displaystyle \frac{g}{2}}[(A^0\stackrel{\mu }{\stackrel{}{}}H^{})W_\mu ^++(A^0\stackrel{\mu }{\stackrel{}{}}H^+)W_\mu ^{}]+i{\displaystyle \frac{gm_W}{2}}(A^0H^{}G^+A^0H^+G^{})`$ (A.3)
$`+`$ $`gm_W[\mathrm{cos}(\beta \alpha )H^0+\mathrm{sin}(\beta \alpha )h^0]W_\mu ^+W^\mu ie(H^{}\stackrel{\mu }{\stackrel{}{}}H^+)A_\mu `$
$`+`$ $`{\displaystyle \frac{gm_W}{2c_W^2}}\mathrm{cos}2\beta [\mathrm{sin}(\alpha +\beta )h^0\mathrm{cos}(\alpha +\beta )H^0]G^+G^{}`$
$`+`$ $`{\displaystyle \frac{gm_W}{4c_W^2}}\mathrm{cos}2\beta [\mathrm{cos}(\alpha +\beta )H^0\mathrm{sin}(\alpha +\beta )h^0]A^0A^0`$
$`+`$ $`gm_W\left[\mathrm{sin}(\alpha \beta ){\displaystyle \frac{\mathrm{cos}2\beta }{2c_W^2}}\mathrm{sin}(\alpha +\beta )\right]h^0H^+H^{}`$
$``$ $`gm_W\left[\mathrm{cos}(\alpha \beta ){\displaystyle \frac{\mathrm{cos}2\beta }{2c_W^2}}\mathrm{cos}(\alpha +\beta )\right]H^0H^+H^{}`$
$`+`$ $`{\displaystyle \frac{g^2}{4}}\left[W_\mu ^+W^\mu \left(1{\displaystyle \frac{\mathrm{cos}^22\beta }{2c_W^2}}\right)G^+G^{}{\displaystyle \frac{\mathrm{cos}^22\beta }{2c_W^2}}H^+H^{}\right]A^0A^0`$
$`+`$ $`i{\displaystyle \frac{ge}{2}}A^\mu A^0[W_\mu ^+H^{}W_\mu ^{}H^+]+e^2H^+H^{}A_\mu A^\mu .`$
On the basis of this we define the $`h^0`$-couplings <sup>5</sup><sup>5</sup>5For the definition of the scalar sector mixing angles we follow the standard notation of e.g. .
$`g_{h\overline{\eta }\eta }{\displaystyle \frac{1}{2}}\xi _Wgm_W\mathrm{sin}(\beta \alpha )`$ $`,`$ $`g_{hWW}gm_W\mathrm{sin}(\beta \alpha ),`$
$`g_{hGG}{\displaystyle \frac{gm_W}{2c_W^2}}\mathrm{cos}2\beta \mathrm{sin}(\alpha +\beta )`$ $`,`$ $`g_{hAA}{\displaystyle \frac{gm_W}{2c_W^2}}\mathrm{cos}2\beta \mathrm{sin}(\alpha +\beta ),`$
$$g_{hH^+H^{}}=gm_W\left[\mathrm{sin}(\alpha \beta )\frac{\mathrm{cos}2\beta }{2c_W^2}\mathrm{sin}(\alpha +\beta )\right],$$
(A.4)
and the $`H^0`$-couplings
$`g_{H^0\overline{\eta }\eta }{\displaystyle \frac{1}{2}}\xi _Wgm_W\mathrm{cos}(\beta \alpha )`$ $`,`$ $`g_{H^0WW}gm_W\mathrm{cos}(\beta \alpha ),`$
$`g_{H^0GG}{\displaystyle \frac{gm_W}{2c_W^2}}\mathrm{cos}2\beta \mathrm{cos}(\alpha +\beta )`$ $`,`$ $`g_{H^0AA}{\displaystyle \frac{gm_W}{2c_W^2}}\mathrm{cos}2\beta \mathrm{cos}(\alpha +\beta ),`$
$$g_{H^0H^+H^{}}=gm_W\left[\mathrm{cos}(\alpha \beta )\frac{\mathrm{cos}2\beta }{2c_W^2}\mathrm{cos}(\alpha +\beta )\right],$$
(A.5)
which are used in Appendix B for expressing the (Higgs-$`W`$) loop contribution of the diagrams in Fig.1, as well as the s-channel $`(h^0,H^0)`$-pole diagrams contained in Figs.2-4.
The chargino loop contribution is described by the diagrams in Fig.2. To define them we first list the chargino mass matrix term as<sup>6</sup><sup>6</sup>6The gaugino fields are defined so that they satisfy $`C\overline{\stackrel{~}{W}}^{+\tau }=\stackrel{~}{W}^{}`$. In such a case there is no $`i`$ in front of $`\stackrel{~}{W}^\pm `$ in (A.6).
$$_{M_\chi }=\left(\begin{array}{c}\stackrel{~}{W}^\tau ,\stackrel{~}{H}_1^\tau \end{array}\right)_LC\left(\begin{array}{cc}M_2& \sqrt{2}m_W\mathrm{sin}\beta \\ \sqrt{2}m_W\mathrm{cos}\beta & +\mu \end{array}\right)\left(\begin{array}{c}\stackrel{~}{W}^+\\ \stackrel{~}{H}_2^+\end{array}\right)_L+\mathrm{h}.\mathrm{c}..$$
(A.6)
Assuming that in MSSM there no new sources of CP-violation, apart from those already known in the Yukawa part of SM; we take the quantities $`(M_2,\mu )`$ as real, but of arbitrary sign. $`C`$ is the usual charge conjugation matrix, and the $`\tau `$ index indicates transposition of the spinorial field. In terms of
$$\stackrel{~}{D}\left[(M_2^2+\mu ^2+2m_W^2)^24(M_2\mu m_W^2\mathrm{sin}(2\beta ))^2\right]^{1/2},$$
(A.7)
the physical chargino masses are expressed as
$$m_{\stackrel{~}{\chi }_1,\stackrel{~}{\chi }_2}=\frac{1}{\sqrt{2}}[M_2^2+\mu ^2+2m_W^2\stackrel{~}{D}]^{1/2}.$$
(A.8)
The mixing-angles $`\varphi _R,\varphi _L`$ in the $`(\stackrel{~}{W}^+,\stackrel{~}{H}_2^+)_L`$ and $`(\stackrel{~}{W}^{},\stackrel{~}{H}_1^{})_L`$ sectors respectively, are defined so that they always lie in the second quarter
$$\frac{\pi }{2}\varphi _L<\pi ,\frac{\pi }{2}\varphi _R<\pi .$$
(A.9)
They are written as
$`\mathrm{cos}\varphi _L`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\stackrel{~}{D}}}}[\stackrel{~}{D}M_2^2+\mu ^2+2m_W^2\mathrm{cos}2\beta ]^{1/2},`$
$`\mathrm{cos}\varphi _R`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\stackrel{~}{D}}}}[\stackrel{~}{D}M_2^2+\mu ^22m_W^2\mathrm{cos}2\beta ]^{1/2}.`$ (A.10)
We always describe the chargino field so that it absorbs a positive chargino particle; i.e. $`\stackrel{~}{\chi }_j\stackrel{~}{\chi }_{j}^{}{}_{}{}^{+}`$ $`(j=1,2)`$. Using this and the sign-quantities
$`\stackrel{~}{}_L`$ $`=`$ $`\mathrm{Sign}(\mu \mathrm{sin}\beta +M_2\mathrm{cos}\beta ),`$
$`\stackrel{~}{}_R`$ $`=`$ $`\mathrm{Sign}(\mu \mathrm{cos}\beta +M_2\mathrm{sin}\beta ),`$
$`\stackrel{~}{\mathrm{\Delta }}_1`$ $`=`$ $`\mathrm{Sign}(M_2[\stackrel{~}{D}M_2^2+\mu ^22m_W^2]2m_W^2\mu \mathrm{sin}2\beta ),`$
$`\stackrel{~}{\mathrm{\Delta }}_2`$ $`=`$ $`\mathrm{Sign}(\mu [\stackrel{~}{D}M_2^2+\mu ^2+2m_W^2]+2m_W^2M_2\mathrm{sin}2\beta ),`$
$`\stackrel{~}{}_{LR}`$ $``$ $`\mathrm{Sign}(M_2\mu +{\displaystyle \frac{\mu ^2+M_2^2}{2}}\mathrm{sin}2\beta )=\stackrel{~}{}_L\stackrel{~}{}_R,`$
$`\stackrel{~}{\mathrm{\Delta }}_{12}`$ $``$ $`\mathrm{Sign}(M_2\mu m_W^2\mathrm{sin}2\beta )=\stackrel{~}{\mathrm{\Delta }}_1\stackrel{~}{\mathrm{\Delta }}_2,`$ (A.11)
the neutral gauge boson-chargino couplings are written as
$``$ $`=`$ $`eA^\mu \overline{\stackrel{~}{\chi }}_j\gamma _\mu \stackrel{~}{\chi }_j{\displaystyle \frac{e}{2s_Wc_W}}Z^\mu \overline{\stackrel{~}{\chi }}_j\left(\gamma _\mu g_{vj}\gamma _\mu \gamma _5g_{aj}\right)\stackrel{~}{\chi }_j`$ (A.12)
$`{\displaystyle \frac{e}{2s_Wc_W}}Z^\mu \left[\overline{\stackrel{~}{\chi }}_1\left(\gamma _\mu g_{v12}\gamma _\mu \gamma _5g_{a12}\right)\stackrel{~}{\chi }_2+\text{h.c.}\right],`$
with
$`g_{v1}={\displaystyle \frac{3}{2}}2s_W^2+{\displaystyle \frac{1}{4}}[\mathrm{cos}2\varphi _L+\mathrm{cos}2\varphi _R]`$ $`,`$ $`g_{a1}={\displaystyle \frac{1}{4}}[\mathrm{cos}2\varphi _L\mathrm{cos}2\varphi _R],`$ (A.13)
$`g_{v2}={\displaystyle \frac{3}{2}}2s_W^2{\displaystyle \frac{1}{4}}[\mathrm{cos}2\varphi _L+\mathrm{cos}2\varphi _R]`$ $`,`$ $`g_{a2}={\displaystyle \frac{1}{4}}[\mathrm{cos}2\varphi _L\mathrm{cos}2\varphi _R],`$ (A.14)
$`g_{v12}={\displaystyle \frac{\mathrm{Sign}(M_2)}{4}}[\stackrel{~}{}_R\stackrel{~}{\mathrm{\Delta }}_{12}\mathrm{sin}2\varphi _R+\stackrel{~}{}_L\mathrm{sin}2\varphi _L],`$
$`g_{a12}={\displaystyle \frac{\mathrm{Sign}(M_2)}{4}}[\stackrel{~}{}_R\stackrel{~}{\mathrm{\Delta }}_{12}\mathrm{sin}2\varphi _R\stackrel{~}{}_L\mathrm{sin}2\varphi _L],`$ (A.15)
where the sign-factors $`(\stackrel{~}{}_L,\stackrel{~}{}_R,\stackrel{~}{\mathrm{\Delta }}_{12})`$ are given in<sup>7</sup><sup>7</sup>7 These expressions are equivalent to those given e.g. in , where a more common definition of the $`\varphi _{L,R}`$-angles is employed. (A.11).
The corresponding chargino-neutral Higgs vertices are
$`_{A^0}`$ $`=`$ $`iA^0\left[g_{A1}\overline{\stackrel{~}{\chi }}_1\gamma _5\stackrel{~}{\chi }_1+g_{A2}\overline{\stackrel{~}{\chi }}_2\gamma _5\stackrel{~}{\chi }_2+\overline{\stackrel{~}{\chi }}_1\left(g_{As12}+\gamma _5g_{Ap12}\right)\stackrel{~}{\chi }_2\overline{\stackrel{~}{\chi }}_2\left(g_{As12}\gamma _5g_{Ap12}\right)\stackrel{~}{\chi }_1\right]`$ (A.16)
$`+`$ $`(g_{h1}h^0+g_{H^01}H^0)\overline{\stackrel{~}{\chi }}_1\stackrel{~}{\chi }_1+(g_{h2}h^0+g_{H^02}H^0)\overline{\stackrel{~}{\chi }}_2\stackrel{~}{\chi }_2,`$
where
$`g_{h1}`$ $`=`$ $`{\displaystyle \frac{g\stackrel{~}{\mathrm{\Delta }}_1}{\sqrt{2}}}[\stackrel{~}{}_L\mathrm{sin}\alpha \mathrm{cos}\varphi _R\mathrm{sin}\varphi _L+\stackrel{~}{}_R\mathrm{cos}\alpha \mathrm{sin}\varphi _R\mathrm{cos}\varphi _L],`$
$`g_{H^01}`$ $`=`$ $`{\displaystyle \frac{g\stackrel{~}{\mathrm{\Delta }}_1}{\sqrt{2}}}[\stackrel{~}{}_L\mathrm{cos}\alpha \mathrm{cos}\varphi _R\mathrm{sin}\varphi _L+\stackrel{~}{}_R\mathrm{sin}\alpha \mathrm{sin}\varphi _R\mathrm{cos}\varphi _L],`$
$`g_{h2}`$ $`=`$ $`{\displaystyle \frac{g\stackrel{~}{\mathrm{\Delta }}_2}{\sqrt{2}}}[\stackrel{~}{}_R\mathrm{sin}\alpha \mathrm{sin}\varphi _R\mathrm{cos}\varphi _L+\stackrel{~}{}_L\mathrm{cos}\alpha \mathrm{sin}\varphi _L\mathrm{cos}\varphi _R],`$
$`g_{H^02}`$ $`=`$ $`{\displaystyle \frac{g\stackrel{~}{\mathrm{\Delta }}_2}{\sqrt{2}}}[\stackrel{~}{}_R\mathrm{cos}\alpha \mathrm{sin}\varphi _R\mathrm{cos}\varphi _L+\stackrel{~}{}_L\mathrm{sin}\alpha \mathrm{cos}\varphi _R\mathrm{sin}\varphi _L],`$
$`g_{A1}`$ $`=`$ $`{\displaystyle \frac{g\stackrel{~}{\mathrm{\Delta }}_1}{\sqrt{2}}}[\stackrel{~}{}_L\mathrm{sin}\beta \mathrm{cos}\varphi _R\mathrm{sin}\varphi _L+\stackrel{~}{}_R\mathrm{cos}\beta \mathrm{sin}\varphi _R\mathrm{cos}\varphi _L],`$
$`g_{A2}`$ $`=`$ $`{\displaystyle \frac{g\stackrel{~}{\mathrm{\Delta }}_2}{\sqrt{2}}}[\stackrel{~}{}_R\mathrm{sin}\beta \mathrm{sin}\varphi _R\mathrm{cos}\varphi _L+\stackrel{~}{}_L\mathrm{cos}\beta \mathrm{sin}\varphi _L\mathrm{cos}\varphi _R],`$
$`g_{As12}`$ $`=`$ $`{\displaystyle \frac{g\mathrm{Sign}(M_2)}{2\sqrt{2}}}[\stackrel{~}{}_{LR}(\stackrel{~}{\mathrm{\Delta }}_1\mathrm{cos}\beta \stackrel{~}{\mathrm{\Delta }}_2\mathrm{sin}\beta )\mathrm{sin}\varphi _R\mathrm{sin}\varphi _L`$
$``$ $`(\stackrel{~}{\mathrm{\Delta }}_1\mathrm{sin}\beta \stackrel{~}{\mathrm{\Delta }}_2\mathrm{cos}\beta )\mathrm{cos}\varphi _L\mathrm{cos}\varphi _R],`$
$`g_{Ap12}`$ $`=`$ $`{\displaystyle \frac{g\mathrm{Sign}(M_2)}{2\sqrt{2}}}[\stackrel{~}{}_{LR}(\stackrel{~}{\mathrm{\Delta }}_1\mathrm{cos}\beta +\stackrel{~}{\mathrm{\Delta }}_2\mathrm{sin}\beta )\mathrm{sin}\varphi _R\mathrm{sin}\varphi _L`$ (A.17)
$``$ $`(\stackrel{~}{\mathrm{\Delta }}_1\mathrm{sin}\beta +\stackrel{~}{\mathrm{\Delta }}_2\mathrm{cos}\beta )\mathrm{cos}\varphi _L\mathrm{cos}\varphi _R].`$
The appearance in (A.15, A.17) of the sign-factors defined in (A.11), guarantees that the physical charginos always have positive masses; irrespective of the signs of the arbitrary real parameters $`M_2`$ and $`\mu `$. These signs are of course intimately related to the definition of the chargino mixing angles employed in (A.9, A.10).
We next turn to $`t`$ and $`b`$ quark loop contribution. The relevant diagrams for the $`t`$-quark case are shown in Fig.3. The necessary vertices are determined by
$`_t`$ $`=`$ $`eA_\mu [Q_t\overline{t}\gamma ^\mu t+Q_b\overline{b}\gamma ^\mu b]+i{\displaystyle \frac{g}{2m_W}}A^0[m_t\mathrm{cot}\beta \overline{t}\gamma _5t+m_b\mathrm{tan}\beta \overline{b}\gamma _5b]`$ (A.18)
$`{\displaystyle \frac{gm_t}{2m_W\mathrm{sin}\beta }}[h^0\mathrm{cos}\alpha +H^0\mathrm{sin}\alpha ]\overline{t}t`$
$`{\displaystyle \frac{gm_b}{2m_W\mathrm{cos}\beta }}[H^0\mathrm{cos}\alpha h^0\mathrm{sin}\alpha ]\overline{b}b,`$
where $`Q_t,Q_b`$ are the $`t`$ and $`b`$ quark charges. The implied $`t`$ quark couplings are
$`g_{h^0tt}={\displaystyle \frac{gm_t}{2m_W\mathrm{sin}\beta }}\mathrm{cos}\alpha ,g_{H^0tt}={\displaystyle \frac{gm_t}{2m_W\mathrm{sin}\beta }}\mathrm{sin}\alpha ,`$
$`g_{Att}={\displaystyle \frac{gm_t}{2m_W}}\mathrm{cot}\beta ,`$ (A.19)
and correspondingly fro the $`b`$-couplings.
Finally, for the stop loop contribution, the relevant interaction Lagrangian is
$`_{\stackrel{~}{t}}`$ $`=`$ $`ieQ_tA_\mu \left[(\stackrel{~}{t}_1^{}\stackrel{\mu }{\stackrel{}{}}\stackrel{~}{t}_1)+(\stackrel{~}{t}_2^{}\stackrel{\mu }{\stackrel{}{}}\stackrel{~}{t}_2)\right]+e^2Q_t^2A_\mu A^\mu (\stackrel{~}{t}_1^{}\stackrel{~}{t}_1+\stackrel{~}{t}_2^{}\stackrel{~}{t}_2)`$ (A.20)
$`+`$ $`i{\displaystyle \frac{gm_t}{2m_W}}(A_t\mathrm{cot}\beta +\mu )A^0[\stackrel{~}{t}_L^{}\stackrel{~}{t}_R\stackrel{~}{t}_R^{}\stackrel{~}{t}_L]`$
$``$ $`\left[{\displaystyle \frac{gm_t^2\mathrm{cos}\alpha }{m_W\mathrm{sin}\beta }}g_Zm_Z\mathrm{sin}(\alpha +\beta )\left(T_t^{(3)}Q_ts_W^2\right)\right]h^0\stackrel{~}{t}_L^{}\stackrel{~}{t}_L`$
$``$ $`[{\displaystyle \frac{gm_t^2\mathrm{cos}\alpha }{m_W\mathrm{sin}\beta }}Q_tg_Zm_Z\mathrm{sin}(\alpha +\beta )s_W^2)]h^0\stackrel{~}{t}_R^{}\stackrel{~}{t}_R`$
$``$ $`{\displaystyle \frac{gm_t}{2m_W\mathrm{sin}\beta }}(\mu \mathrm{sin}\alpha +A_t\mathrm{cos}\alpha )h^0(\stackrel{~}{t}_R^{}\stackrel{~}{t}_L+\stackrel{~}{t}_L^{}\stackrel{~}{t}_R)`$
$``$ $`\left[{\displaystyle \frac{gm_t^2\mathrm{sin}\alpha }{m_W\mathrm{sin}\beta }}+g_Zm_Z\mathrm{cos}(\alpha +\beta )\left(T_t^{(3)}Q_ts_W^2\right)\right]H^0\stackrel{~}{t}_L^{}\stackrel{~}{t}_L`$
$``$ $`[{\displaystyle \frac{gm_t^2\mathrm{sin}\alpha }{m_W\mathrm{sin}\beta }}+Q_tg_Zm_Z\mathrm{cos}(\alpha +\beta )s_W^2)]H^0\stackrel{~}{t}_R^{}\stackrel{~}{t}_R`$
$`+`$ $`{\displaystyle \frac{gm_t}{2m_W\mathrm{sin}\beta }}(\mu \mathrm{cos}\alpha A_t\mathrm{sin}\alpha )H^0(\stackrel{~}{t}_R^{}\stackrel{~}{t}_L+\stackrel{~}{t}_L^{}\stackrel{~}{t}_R)`$
$``$ $`\left[{\displaystyle \frac{g^2m_t^2}{4m_W^2}}\mathrm{cot}^2\beta {\displaystyle \frac{1}{4}}g_Z^2(T_t^{(3)}s_W^2Q_t)\mathrm{cos}2\beta \right]A^0A^0\stackrel{~}{t}_L^{}\stackrel{~}{t}_L`$
$``$ $`\left[{\displaystyle \frac{g^2m_t^2}{4m_W^2}}\mathrm{cot}^2\beta {\displaystyle \frac{1}{4}}g_Z^2s_W^2Q_t\mathrm{cos}2\beta \right]A^0A^0\stackrel{~}{t}_R^{}\stackrel{~}{t}_R,`$
where, as usual, $`g_Z=g/c_W`$. The various neutral Higgs-$`\stackrel{~}{t}_{L,R}`$ couplings are determined from the coefficients of the various terms in (A.20). For two examples, we note<sup>8</sup><sup>8</sup>8As usual, in the definition of the $`A^0A^0\stackrel{~}{t}_L^{}\stackrel{~}{t}_L`$ and $`A^0A^0\stackrel{~}{t}_R^{}\stackrel{~}{t}_R`$ couplings from the last two terms in (A.20), the relevant coefficient is multiplied by a 2, due to the identity of the two $`A^0`$-fields.
$$g_{A\stackrel{~}{t}_L\stackrel{~}{t}_R}=\frac{gm_t}{2m_W}(A_t\mathrm{cot}\beta +\mu ),$$
$$g_{AA\stackrel{~}{t}_R\stackrel{~}{t}_R}=2\left[\frac{g^2m_t^2}{4m_W^2}\mathrm{cot}^2\beta \frac{1}{4}g_Z^2s_W^2Q_t\mathrm{cos}2\beta \right].$$
For determining the corresponding couplings to the physical $`\stackrel{~}{t}_{1,2}`$ we write
$$(\begin{array}{c}\stackrel{~}{t}_L\\ \stackrel{~}{t}_R\end{array})=(\begin{array}{cc}\mathrm{cos}\theta _t& \mathrm{sin}\theta _t\mathrm{Sign}(A_t\mu \mathrm{cot}\beta )\\ \mathrm{sin}\theta _t\mathrm{Sign}(A_t\mu \mathrm{cot}\beta )& \mathrm{cos}\theta _t\end{array})(\begin{array}{c}\stackrel{~}{t}_1\\ \stackrel{~}{t}_2\end{array}),$$
(A.21)
where $`\theta _t`$ is fully determined by<sup>9</sup><sup>9</sup>9 The quantities $`m_{\stackrel{~}{t}_L}^2`$, $`m_{\stackrel{~}{t}_R}^2`$ in (A.22) are the usual soft SUSY breaking parameters in which the small D-contributions have also been included.
$$\mathrm{sin}(2\theta _t)=\frac{2m_t|A_t\mu \mathrm{cot}\beta |}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2},\mathrm{cos}(2\theta _t)=\frac{m_{\stackrel{~}{t}_L}^2m_{\stackrel{~}{t}_R}^2}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2},$$
(A.22)
while $`A_t`$ is also real. We observe from (A.22) that
$$\frac{\pi }{2}<\theta _t<\pi ,$$
since $`m_{\stackrel{~}{t}_1}<m_{\stackrel{~}{t}_2}`$, by definition. We have checked that this stop-mixing-formalism is equivalent to the usual one found e.g. in .
Appendix B: The MSSM contributions to $`\gamma \gamma A^0A^0`$.
The invariant helicity amplitudes for the process
$$\gamma (p_1,\lambda _1)\gamma (p_2,\lambda _2)A^0(p_3)A^0(p_4),$$
(B.1)
are denoted as<sup>10</sup><sup>10</sup>10Their sign is related to the sign of the $`S`$-matrix through $`S_{\lambda _1\lambda _2}=1+i(2\pi )^4\delta (p_fp_i)F_{\lambda _1\lambda _2}`$. $`F_{\lambda _1\lambda _2}(\widehat{s},\widehat{t},\widehat{u})`$, where the particle-momenta and helicities of the incoming photons, are indicated in parentheses. Assuming no new (beyond SM) source of CP violation, these invariant helicity amplitudes satisfy
$$F_{\lambda _1,\lambda _2}(\widehat{s},\widehat{t},\widehat{u})=F_{\lambda _1,\lambda _2}(\widehat{s},\widehat{t},\widehat{u}),$$
(B.2)
which implies that there are only two independent helicity amplitudes; namely $`F_{++}`$ and $`F_+`$. As in we make the definitions
$`\widehat{s}=(p_1+p_2)^2={\displaystyle \frac{4m_A^2}{1\beta _A^2}},\widehat{t}=(p_1p_3)^2,\widehat{u}=(p_1p_4)^2`$ , (B.3)
$`\widehat{s}_4=\widehat{s}4m_A^2,\widehat{s}_2=\widehat{s}2m_A^2,\widehat{t}_1=\widehat{t}m_A^2,\widehat{u}_1=\widehat{u}m_A^2`$ , (B.4)
$`\widehat{t}=m_A^2{\displaystyle \frac{\widehat{s}}{2}}(1\beta _A\mathrm{cos}\vartheta ^{}),\widehat{u}=m_A^2{\displaystyle \frac{\widehat{s}}{2}}(1+\beta _A\mathrm{cos}\vartheta ^{})`$ $`,`$ (B.5)
$`Y=\widehat{t}\widehat{u}m_A^4={\displaystyle \frac{\widehat{s}^2\beta _A^2}{4}}\mathrm{sin}^2\vartheta ^{}`$ . (B.6)
where $`\beta _A`$ is the $`A^0`$-velocity in the $`A^0A^0`$-c.m. frame, and $`\vartheta ^{}`$ the c.m. scattering angle. Moreover, the combinations
$$m_{ab}^2=m_A^2+m_a^2m_b^2,\widehat{s}_{ab}=\widehat{s}m_{ab}^2,$$
(B.7)
often appear below for the charged particle pairs $`(a,b)=(H^\pm ,W^{})`$, $`(W^{},H^\pm )`$, and $`(\stackrel{~}{\chi }_1,\stackrel{~}{\chi }_2)`$.
All 1-loop results are expressed in terms of the $`C_0`$ and $`D_0`$ Passarino-Veltman functions , for which we follow the notation of . Similarly to , we also introduce the short hand writing<sup>11</sup><sup>11</sup>11In the middle terms of (B.8-B.13) $`k_1=p_1`$, $`k_2=p_2`$ denote the momenta of the photons, while $`k_3=p_3`$, $`k_4=p_4`$ those of the $`A^0`$, always taken as incoming; compare (B.1).
$`C_0^{abc}(\widehat{s})C_0(k_1,k_2)`$ $`=`$ $`C_0(0,0,\widehat{s};m_a,m_b,m_c),`$ (B.8)
$`C_A^{abc}(\widehat{t})C_0(k_3,k_1)`$ $`=`$ $`C_0(m_A^2,0,\widehat{t};m_a,m_b,m_c),`$ (B.9)
$`C_{AA}^{abc}(\widehat{s})C_0(k_3,k_4)`$ $`=`$ $`C_0(m_A^2,m_A^2,\widehat{s};m_a,m_b,m_c),`$ (B.10)
$`D_{AA}^{abcd}(\widehat{s},\widehat{t})D_0(k_4,k_3,k_1)`$ $`=`$ $`D_0(m_A^2,m_A^2,0,0,\widehat{s},\widehat{t};m_a,m_b,m_c,m_d),`$ (B.11)
$`D_{AA}^{abcd}(\widehat{s},\widehat{u})D_0(k_3,k_4,k_1)`$ $`=`$ $`D_0(m_A^2,m_A^2,0,0,\widehat{s},\widehat{u};m_a,m_b,m_c,m_d),`$ (B.12)
$`D_{AA}^{abcd}(\widehat{t},\widehat{u})D_0(k_3,k_1,k_4)`$ $`=`$ $`D_0(m_A^2,0,m_A^2,0,\widehat{t},\widehat{u};m_a,m_b,m_c,m_d),`$ (B.13)
$`D_{AA}^{abcd}(\widehat{u},\widehat{t})D_0(k_4,k_1,k_3)`$ $`=`$ $`D_0(m_A^2,0,m_A^2,0,\widehat{u},\widehat{t};m_a,m_b,m_c,m_d),`$ (B.14)
in which we have also emphasized the fact that the masses running along the various sides of the loop may be different.
The fact that the masses along the loops in Figs.1, 2, 4 may be different, considerably complicates the formulae. Nevertheless, expressions analogous to those encountered for the SM contributions to $`\gamma \gamma ZZ`$ may be defined, which allows writing the amplitudes in a compact way. We thus define
$`\stackrel{~}{F}^{ab}(\widehat{s},\widehat{t},\widehat{u})`$ $`=`$ $`D_{AA}^{abba}(\widehat{t},\widehat{u})+D_{AA}^{abaa}(\widehat{s},\widehat{t})+D_{AA}^{abaa}(\widehat{s},\widehat{u}),`$ (B.15)
$`E_1^{ab}(\widehat{s},\widehat{t})`$ $`=`$ $`\widehat{t}_1\left[C_A^{abb}(\widehat{t})+C_A^{baa}(\widehat{t})\right]\widehat{s}\widehat{t}D_{AA}^{abaa}(\widehat{s},\widehat{t}),`$ (B.16)
$`E_2^{ab}(\widehat{t},\widehat{u})`$ $`=`$ $`\widehat{t}_1\left[C_A^{abb}(\widehat{t})+C_A^{baa}(\widehat{t})\right]+\widehat{u}_1\left[C_A^{abb}(\widehat{u})+C_A^{baa}(\widehat{u})\right]`$ (B.17)
$``$ $`YD_{AA}^{abba}(\widehat{t},\widehat{u}),`$
which are closely related to the definitions in Eqs.(A.22-A.24) in . We also note that
$`D_{AA}^{abba}(\widehat{t},\widehat{u})=D_{AA}^{abba}(\widehat{u},\widehat{t})=D_{AA}^{baab}(\widehat{t},\widehat{u})=D_{AA}^{baab}(\widehat{u},\widehat{t}),`$ (B.18)
$`\stackrel{~}{F}^{ab}(\widehat{s},\widehat{t},\widehat{u})=\stackrel{~}{F}^{ab}(\widehat{s},\widehat{u},\widehat{t}),E_2^{ab}(\widehat{t},\widehat{u})=E_2^{ab}(\widehat{u},\widehat{t})=E_2^{ba}(\widehat{t},\widehat{u}).`$
The $`(W^\pm ,H^\pm )`$-loop diagrams.
There two kinds of contributions to the invariant amplitudes $`F_{\lambda _1\lambda _2}(\gamma \gamma A^0A^0)`$ from the diagrams of Fig.1. The first arises from the two diagrams in the first row of Fig.1 and contains $`\widehat{s}`$-pole contributions generated by exchanging the CP-even neutral Higgs particles $`h^0`$ and $`H^0`$. For the $`\widehat{s}`$-channel $`h^0`$-case, this is given by
$`F_{++}^{WH^\pm (h^0\text{pole})}={\displaystyle \frac{e^2g_{hAA}}{8\pi ^2(\widehat{s}m_h^2)}}\{g_{hH^+H^{}}[1+2m_{H^\pm }^2C_0^{H^+H^+H^+}(\widehat{s})]+g_{hGG}2g_{h\eta \eta }`$
$`4g_{hWW}+2[(g_{hGG}2g_{h\eta \eta }4g_{hWW})m_W^2+2g_{hWW}\widehat{s}]C_0^{WWW}(\widehat{s})\},`$ (B.19)
while for the $`H^0`$-case we get
$$F_{++}^{WH^\pm (H^0\text{pole})}=F_{++}^{WH^\pm (h^0\text{pole})}(h^0H^0),$$
(B.20)
where the needed $`h^0`$ and $`H^0`$ couplings are given in (A.4, A.5). For such contributions we obviously also have
$$F_+^{WH^\pm (H^0\text{pole})}=F_+^{WH^\pm (h^0\text{pole})}=0.$$
(B.21)
The second comes from the 3rd to last diagrams in Fig.1, and it is written as
$`F_{++}^{WH^\pm }={\displaystyle \frac{e^2g^2}{16\pi ^2}}\{4+[6{\displaystyle \frac{\mathrm{cos}^2(2\beta )}{c_W^2}}]m_W^2C_0^{WWW}(\widehat{s})+[2+{\displaystyle \frac{\mathrm{cos}^2(2\beta )}{c_W^2}}]m_{H^\pm }^2C_0^{H^\pm H^\pm H^\pm }(\widehat{s})`$
$`+2{\displaystyle \frac{m_{HW}^2}{\widehat{s}}}E_2^{H^\pm W}(\widehat{t},\widehat{u})+2m_{H^\pm }^2\widehat{s}_{HW}\stackrel{~}{F}^{H^\pm W}(\widehat{s},\widehat{t},\widehat{u})`$
$`+2(\widehat{s}m_{H^\pm }^2m_{HW}^2m_W^2)\stackrel{~}{F}^{WH^\pm }(\widehat{s},\widehat{t},\widehat{u})\},`$ (B.22)
$`F_+^{WH^\pm }={\displaystyle \frac{e^2g^2}{16\pi ^2Y}}\{\widehat{s}[2(m_{HW}^2m_{WH}^2Y)+\widehat{s}(m_H^2m_W^2)+\widehat{t}\widehat{t}_1+\widehat{u}\widehat{u}_1]C_0^{WWW}(\widehat{s})`$
$`+\widehat{s}\widehat{s}_{HW}(\widehat{s}_{HW}m_{HW}^2)C_0^{H^\pm H^\pm H^\pm }(\widehat{s})+\widehat{s}_{HW}(\widehat{t}^2+\widehat{u}^22m_A^4)\left[C_{AA}^{H^\pm WH^\pm }(\widehat{s})+C_{AA}^{WH^\pm W}(\widehat{s})\right]`$
$`+2\left[\widehat{s}\widehat{s}_{HW}(m_{H^\pm }^2m_W^2)^2+Y\left[2\widehat{s}m_W^2+Ym_{HW}^2(m_W^2+m_{H^\pm }^2)\right]\right]\stackrel{~}{F}^{WH^\pm }(\widehat{s},\widehat{t},\widehat{u})`$
$`+\{\widehat{s}_{HW}[\widehat{s}(m_{H^\pm }^2m_W^2)^2+\widehat{s}\widehat{t}(\widehat{t}2m_W^2)2m_{H^\pm }^2\widehat{t}_1^2][D_{AA}^{H^\pm WH^\pm H^\pm }(\widehat{s},\widehat{t})D_{AA}^{WH^\pm WW}(\widehat{s},\widehat{t})]`$
$`+2(m_A^4+\widehat{t}^2\widehat{t}m_{WH}^2)E_1^{WH^\pm }(\widehat{s},\widehat{t})+(\widehat{t}\widehat{u})\}\},`$ (B.23)
where all needed quantities have been defined in (B.3-B.17).
The chargino loop diagrams
The relevant diagrams are presented in Fig.2. The first diagram in Fig.2 contains an $`\widehat{s}`$-channel pole due to $`(h^0,H^0)`$ exchanges, and is characterized by a single chargino $`\stackrel{~}{\chi }_i`$ running along the loop. It gives
$`F_{++}^{\stackrel{~}{\chi }_i(\text{pole})}={\displaystyle \frac{e^2m_{\stackrel{~}{\chi }_i}}{4\pi ^2}}\left({\displaystyle \frac{g_{hi}g_{hAA}}{\widehat{s}m_h^2}}+{\displaystyle \frac{g_{Hi}g_{H^0AA}}{\widehat{s}m_H^2}}\right)\left[2+(4m_{\stackrel{~}{\chi }_i}^2\widehat{s})C_0^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{s})\right],`$ (B.24)
$`F_+^{\stackrel{~}{\chi }_i(\text{pole})}=0.`$ (B.25)
The other diagrams in Fig.2 involve contributions containing either a single chargino running along the loop, or mixed contributions where both charginos run. The single chargino contribution is
$`F_{++}^{\stackrel{~}{\chi }_i}`$ $`=`$ $`{\displaystyle \frac{e^2g_{Ai}^2}{4\pi ^2}}\{2+4m_{\stackrel{~}{\chi }_i}^2C_0^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{s})m_{\stackrel{~}{\chi }_i}^2(\widehat{t}+\widehat{u})\stackrel{~}{F}^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{s},\widehat{t},\widehat{u})`$ (B.26)
$`+`$ $`{\displaystyle \frac{m_A^2}{\widehat{s}}}E_2^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{t},\widehat{u})\},`$
$`F_+^{\stackrel{~}{\chi }_i}`$ $`=`$ $`{\displaystyle \frac{e^2g_{Ai}^2}{8\pi ^2Y}}\{\widehat{s}(\widehat{s}_2^22Y)C_0^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{s})+\widehat{s}_2(\widehat{t}^2+\widehat{u}^22m_A^4)C_{AA}^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{s})`$ (B.27)
$`+`$ $`2m_{\stackrel{~}{\chi }_i}^2\widehat{s}_2Y\stackrel{~}{F}^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{s},\widehat{t},\widehat{u})+(\widehat{t}^2+m_A^4)E_1^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{s},\widehat{t})`$
$`+`$ $`(\widehat{u}^2+m_A^4)E_1^{\stackrel{~}{\chi }_i\stackrel{~}{\chi }_i}(\widehat{s},\widehat{u})\}.`$
Of course, in calculating the total ”single” chargino contribution, the results in (B.24-B.27) should be summed for both the $`\stackrel{~}{\chi }_1`$ and $`\stackrel{~}{\chi }_2`$ charginos. The necessary couplings are given in (A.17).
The considerably more complicated mixed chargino contribution, arising from the 3rd and 4th diagram in Fig.2, is
$`F_{++}^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2}`$ $`=`$ $`{\displaystyle \frac{e^2}{4\pi ^2}}(g_{As12}^2+g_{Ap12}^2)\{2+4m_{\stackrel{~}{\chi }_1}^2C_0^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_1\stackrel{~}{\chi }_1}(\widehat{s})+{\displaystyle \frac{1}{2\widehat{s}}}(\widehat{s}X)E_2^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2}(\widehat{t},\widehat{u})`$ (B.28)
$`m_{\stackrel{~}{\chi }_1}\left[\widehat{s}\left(m_{\stackrel{~}{\chi }_1}+{\displaystyle \frac{(g_{As12}^2g_{Ap12}^2)}{(g_{As12}^2+g_{Ap12}^2)}}m_{\stackrel{~}{\chi }_2}\right)m_{\stackrel{~}{\chi }_1}X\right]\stackrel{~}{F}^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2}(\widehat{s},\widehat{t},\widehat{u})`$
$`+`$ $`(\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2)\},`$
$`F_+^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2}`$ $`=`$ $`{\displaystyle \frac{e^2(g_{As12}^2+g_{Ap12}^2)}{8\pi ^2Y}}\{\widehat{s}[X(\widehat{s}_{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2}m_{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2})2Y]C_0^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_1\stackrel{~}{\chi }_1}(\widehat{s})`$ (B.29)
$`+`$ $`X\left[(m_{\stackrel{~}{\chi }_1}^2+m_{\stackrel{~}{\chi }_2}^2)Y+\widehat{s}(m_{\stackrel{~}{\chi }_1}^2m_{\stackrel{~}{\chi }_2}^2)^2\right]\stackrel{~}{F}^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2}(\widehat{s},\widehat{t},\widehat{u})`$
$`+`$ $`X(\widehat{t}^2+\widehat{u}^22m_A^4)C_{AA}^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2\stackrel{~}{\chi }_1}(\widehat{s})(\widehat{t}X+Y)E_1^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2}(\widehat{s},\widehat{t})`$
$``$ $`(\widehat{u}X+Y)E_1^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2}(\widehat{s},\widehat{u})(m_{\stackrel{~}{\chi }_1}^2m_{\stackrel{~}{\chi }_2}^2)\left[2\widehat{t}_1^2X+Y(X\widehat{s})\right]D_{AA}^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2\stackrel{~}{\chi }_1\stackrel{~}{\chi }_1}(\widehat{s},\widehat{t})`$
$``$ $`(m_{\stackrel{~}{\chi }_1}^2m_{\stackrel{~}{\chi }_2}^2)[2\widehat{u}_1^2X+Y(X\widehat{s})]D_{AA}^{\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2\stackrel{~}{\chi }_1\stackrel{~}{\chi }_1}(\widehat{s},\widehat{u})+(\stackrel{~}{\chi }_1\stackrel{~}{\chi }_2)\},`$
which where
$$X=\widehat{s}_2+2m_{\stackrel{~}{\chi }_1}^2+2m_{\stackrel{~}{\chi }_2}^2+4m_{\stackrel{~}{\chi }_1}m_{\stackrel{~}{\chi }_2}\frac{(g_{As12}^2g_{Ap12}^2)}{(g_{As12}^2+g_{Ap12}^2)},$$
(B.30)
and the necessary couplings are given in (A.17).
The $`t`$ and $`b`$-quark loop diagrams.
The top-loop contribution arises from the diagrams in Fig.3. The first of them contains the $`(h^0,H^0)`$-pole contribution
$`F_{++}^{t(\text{pole})}={\displaystyle \frac{3e^2Q_t^2m_t}{4\pi ^2}}\left({\displaystyle \frac{g_{htt}g_{hAA}}{\widehat{s}m_h^2}}+{\displaystyle \frac{g_{H^0tt}g_{H^0AA}}{\widehat{s}m_H^2}}\right)\left[2+(4m_t^2\widehat{s})C_0^{ttt}(\widehat{s})\right],`$ (B.31)
$`F_+^{t(\text{pole})}=0,`$ (B.32)
while the second gives
$`F_{++}^t`$ $`=`$ $`{\displaystyle \frac{3e^2Q_t^2g_{Att}^2}{4\pi ^2}}[2+4m_t^2C_0^{ttt}(\widehat{s})m_t^2(\widehat{t}+\widehat{u})\stackrel{~}{F}^{tt}(\widehat{s},\widehat{t},\widehat{u})`$ (B.33)
$`+`$ $`{\displaystyle \frac{m_A^2}{\widehat{s}}}E_2^{tt}(\widehat{t},\widehat{u})],`$
$`F_+^t`$ $`=`$ $`{\displaystyle \frac{3e^2Q_t^2g_{Att}^2}{8\pi ^2Y}}\{\widehat{s}(\widehat{s}_2^22Y)C_0^{ttt}(\widehat{s})+\widehat{s}_2(\widehat{t}^2+\widehat{u}^22m_A^4)C_{AA}^{ttt}(\widehat{s})`$ (B.34)
$`+`$ $`2m_t^2\widehat{s}_2Y\stackrel{~}{F}^{tt}(\widehat{s},\widehat{t},\widehat{u})+(\widehat{t}^2+m_A^4)E_1^{tt}(\widehat{s},\widehat{t})+(\widehat{u}^2+m_A^4)E_1^{tt}(\widehat{s},\widehat{u})\}.`$
All needed couplings are given in (A.4, A.5, A.19). In (B.31-(B.34) a factor three for colour has already been introduced. The corresponding $`b`$-quark contribution is analogously obtained through (A.18) and the use of $`Q_b`$ instead of $`Q_t`$.
$`\stackrel{~}{t}`$-loop diagrams
These diagrams are shown in Fig.4 and will be relevant in case one or two stop sqarks turn out to be not too heavy. The first two of these diagrams describe the $`(h^0,H^0)`$ $`\widehat{s}`$-channel pole contributions and have just one kind of $`\stackrel{~}{t}_i`$ running along the loop. For each such $`\stackrel{~}{t}_i`$, the pole contribution is
$`F_{++}^{\stackrel{~}{t}_i(\text{pole})}={\displaystyle \frac{3e^2Q_t^2}{8\pi ^2}}\left({\displaystyle \frac{g_{hAA}g_{h\stackrel{~}{t}_i\stackrel{~}{t}_i}}{\widehat{s}m_h^2}}+{\displaystyle \frac{g_{HAA}g_{H\stackrel{~}{t}_i\stackrel{~}{t}_i}}{\widehat{s}m_H^2}}\right)\left[1+2m_{\stackrel{~}{t}_i}^2C_0^{\stackrel{~}{t}_i\stackrel{~}{t}_i\stackrel{~}{t}_i}(\widehat{s})\right],`$ (B.35)
$$F_+^{\stackrel{~}{t}_i(\text{pole})}=0.$$
(B.36)
In addition, we have the loop contribution from the no-pole last five diagrams of Fig.4
$`F_{++}^{\stackrel{~}{t}_1\stackrel{~}{t}_2}`$ $`=`$ $`{\displaystyle \frac{3e^2Q_t^2}{8\pi ^2}}\{g_{AA\stackrel{~}{t}_1\stackrel{~}{t}_1}[1+2m_{\stackrel{~}{t}_1}^2C_0^{\stackrel{~}{t}_1\stackrel{~}{t}_1\stackrel{~}{t}_1}(\widehat{s})]+{\displaystyle \frac{g_{A\stackrel{~}{t}_1\stackrel{~}{t}_2}^2}{\widehat{s}}}E_2^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(\widehat{t},\widehat{u})`$ (B.37)
$``$ $`2g_{A\stackrel{~}{t}_1\stackrel{~}{t}_2}^2m_{\stackrel{~}{t}_1}^2\stackrel{~}{F}^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(\widehat{s},\widehat{t},\widehat{u})+(\stackrel{~}{t}_1\stackrel{~}{t}_2)\},`$
$`F_+^{\stackrel{~}{t}_1\stackrel{~}{t}_2}`$ $`=`$ $`{\displaystyle \frac{3e^2Q_t^2g_{A\stackrel{~}{t}_1\stackrel{~}{t}_2}^2}{8\pi ^2Y}}\{\widehat{s}(\widehat{s}2m_{\stackrel{~}{t}_1\stackrel{~}{t}_2}^2)C_0^{\stackrel{~}{t}_1\stackrel{~}{t}_1\stackrel{~}{t}_1}(\widehat{s})+(\widehat{t}^2+\widehat{u}^22m_A^4)C_{AA}^{\stackrel{~}{t}_1\stackrel{~}{t}_2\stackrel{~}{t}_1}(\widehat{s})`$ (B.38)
$`+`$ $`Y(m_{\stackrel{~}{t}_1}^2+m_{\stackrel{~}{t}_2}^2)D_{AA}^{\stackrel{~}{t}_1\stackrel{~}{t}_2\stackrel{~}{t}_2\stackrel{~}{t}_1}(\widehat{t},\widehat{u})+\widehat{s}(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2)^2\stackrel{~}{F}^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(\widehat{s},\widehat{t},\widehat{u})`$
$``$ $`\widehat{t}E_1^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(\widehat{s},\widehat{t})\widehat{u}E_1^{\stackrel{~}{t}_1\stackrel{~}{t}_2}(\widehat{s},\widehat{u})2(\widehat{s}\widehat{t}m_{\stackrel{~}{t}_2}^2+\widehat{t}_1^2m_{\stackrel{~}{t}_1}^2)D_{AA}^{\stackrel{~}{t}_1\stackrel{~}{t}_2\stackrel{~}{t}_1\stackrel{~}{t}_1}(\widehat{s},\widehat{t})`$
$``$ $`2(\widehat{s}\widehat{u}m_{\stackrel{~}{t}_2}^2+\widehat{u}_1^2m_{\stackrel{~}{t}_1}^2)D_{AA}^{\stackrel{~}{t}_1\stackrel{~}{t}_2\stackrel{~}{t}_1\stackrel{~}{t}_1}(\widehat{s},\widehat{u})+(\stackrel{~}{t}_1\stackrel{~}{t}_2)\},`$
which involve contributions either from a single $`\stackrel{~}{t}_j`$ running along the loop, or mixed contributions involving both both $`\stackrel{~}{t}_1,\stackrel{~}{t}_2`$. In (B.35-(B.38) a factor three for colour has already been introduced, while the necessary couplings are given by combining (A.20, A.21).
If other kinds of sfermions turn out also to be light, then their contribution can readily be derived from (B.35-B.38) by changing the appropriate couplings.
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# Untitled Document
hep-th/0007175
Overview Of $`K`$-Theory
Applied To Strings
Edward Witten
School of Natural Sciences, Institute for Advanced Study
Olden Lane, Princeton, NJ 08540, USA
and
Department of Physics, Caltech, Pasadena CA 91125
and
CIT-USC Center For Theoretical Physics, Los Angeles CA
$`K`$-theory provides a framework for classifying Ramond-Ramond (RR) charges and fields. $`K`$-theory of manifolds has a natural extension to $`K`$-theory of noncommutative algebras, such as the algebras considered in noncommutative Yang-Mills theory or in open string field theory. In a number of concrete problems, the $`K`$-theory analysis proceeds most naturally if one starts out with an infinite set of $`D`$-branes, reduced by tachyon condensation to a finite set. This suggests that string field theory should be reconsidered for $`N=\mathrm{}`$.
August, 2000
1. Introduction And Definition Of $`K(X)`$
A $`D`$-brane wrapped on a submanifold $`S`$ of spacetime may carry a nonzero Ramond-Ramond (RR) charge. RR fields are $`p`$-forms, and superficially it seems that the conserved charge should be measured by the cohomology class of the RR form (or of the cycle $`S`$ itself). However, $`D`$-branes carry gauge fields; and gauge fields are not natural in (co)homology theory. They are natural in “$`K`$-theory.” $`K`$-theory has been used to answer some questions about RR charges and fields; the aim of the present article<sup>1</sup> The article is based on my lecture at Strings 2000, Univ. of Michigan, July 10, 2000 as well as earlier lectures at the CIT-USC Center for Theoretical Physics. I thank both audiences for questions and comments. is to give an overview of this, along with some speculations (for which I have only very modest evidence) about how one might want to rethink open string field theory in the large $`N`$ limit. These subjects fill sections 1-3. Some mathematical details have been postponed to section 4.
If $`X`$ is spacetime and $`𝒜(X)`$ is the commutative, associative algebra of continuous complex-valued functions on $`X`$, then the $`K`$-theory of $`X`$ can be defined in terms of representations of $`𝒜(X)`$. A representation of a ring is usually called a module. Here are some examples of $`𝒜(X)`$-modules.
The most obvious example of an $`𝒜(X)`$-module is $`𝒜(X)`$ itself. For $`f𝒜(X)`$ (regarded as a ring) and $`g𝒜(X)`$ (regarded as a module), we define
$$f(g)=fg,$$
where on the right hand side the multiplication occurs in $`𝒜(X)`$. This obviously obeys the defining condition of a module, which is that $`(f_1f_2)(g)=f_1(f_2(g))`$.
More generally, consider a $`Dp`$-brane (or a collection of $`N`$ $`Dp`$-branes for some $`N>0`$) wrapped on a submanifold $`S`$ of $`X`$, with any Chan-Paton gauge bundle $`W`$ on the $`D`$-brane. Let $`M(S)`$ be the space of sections of $`W`$, that is, the space of one-particle states for a charged scalar coupled to the bundle $`W`$. Then $`M(S)`$ is an $`𝒜(X)`$-module; for $`f𝒜(X)`$, $`gM(S)`$, we simply set again $`f(g)=fg`$. On the right hand side, the multiplication is defined by restricting $`f`$ (which is a function on $`X`$) to $`S`$ and then multiplying $`f`$ and $`g`$.
So in, say, Type IIB superstring theory, a collection of $`D9`$-branes defines a representation or module $`E`$ of $`𝒜(X)`$. A collection of $`\overline{D9}`$-branes defines another module $`F`$. So any configuration of $`D9`$ and $`\overline{D9}`$-branes determines a pair $`(E,F)`$.
To classify $`D`$-brane charge, we want to classify pairs $`(E,F)`$ modulo physical processes. An important process is brane-antibrane creation and annihilation – the creation or annihilation of a set of $`D9`$’s and $`\overline{D9}`$’s each bearing the same gauge bundle $`G`$. This amounts to
$$(E,F)(EG,FG).$$
The equivalence classes make up a group called $`K(X)`$ (or $`K(𝒜(X))`$ if we want to make the interpretation in terms of $`𝒜(X)`$-modules more explicit). The addition law in this group is just
$$(E,F)+(E^{}F^{})=(EE^{},FF^{}).$$
The inverse of $`(E,F)`$ is $`(F,E)`$; note that $`(E,F)(F,E)=(EF,EF)`$, and using the equivalence relation (1.1), this is equivalent to zero. $`D`$-branes of Type IIB carry conserved charges that take values in $`K(X)`$.
In the above definition of $`K(X)`$, we used only ninebranes, even though, as we explained earlier, an $`𝒜(X)`$ module can be constructed using $`Dp`$-branes (or antibranes) for any $`p`$. In fact, we can classify $`D`$-brane charge just using the ninebranes, and then build the $`Dp`$-branes of $`p<9`$ via pairs $`(E,F)`$ with a suitable tachyon condensate. This construction (due to Atiyah, Bott, and Shapiro) was reviewed in , section 4, in the $`D`$-brane context. A more systematic explanation of the definition of $`K(X)`$ in terms of modules, clarifying the role of ninebranes, can be found in section 4.
Advantages Of $`K`$-Theory Description
What do we gain by knowing that $`D`$-brane charge is classified by $`K`$-theory?
First of all, it is the right answer. Wherever one looks closely at topological properties of RR charges (or fields), one sees effects that reflect the $`K`$-theory structure. For example, there are stable $`D`$-brane states (like the nonsupersymmetric $`D0`$-branes of Type I) that would not exist if $`D`$-brane charge were classified by cohomology instead of $`K`$-theory. Conversely, it is possible to have a $`D`$-brane state that would be stable if $`D`$-brane charge were measured by cohomology, but which is in fact unstable (via a process that involves nucleation of $`D9`$-$`\overline{D9}`$ pairs in an intermediate state). This occurs in Type II superstring theory, in which a $`D`$-brane wrapped on a homologically nontrivial cycle in spacetime is in fact in certain cases unstable. Finally, the $`K`$-theory interpretation of $`D`$-branes is needed to make sense of a certain global worldvolume anomaly.
But I think that there is a deeper reason that it is good to know about the $`K`$-theory interpretation of $`D`$-branes: it may be naturally adapted for stringy generalizations. In fact (though some mathematical details have been postponed to section 4), we defined $`K(X)`$ in terms of representations of the algebra $`𝒜(X)`$ of functions on spacetime. We can similarly define $`K(𝒜)`$ for any noncommutative algebra $`𝒜`$, in terms of pairs $`(E,F)`$ of $`𝒜`$-modules. By contrast, we would not have an equally useful and convenient notion of “cohomology” if the algebra of functions on spacetime is replaced by a noncommutative ring.
For example, turning on a Neveu-Schwarz $`B`$-field, we can make $`𝒜(X)`$ noncommutative; the associated $`K(𝒜)`$ was used by Connes, Douglas, and Schwarz in the original paper on noncommutative Yang-Mills theory applied to string theory . This is an interesting example, even though it involves only the zero modes of the strings. One would much like to have a fully stringy version involving a noncommutative algebra constructed using all of the modes of the string, not just the zero modes.
What is the right noncommutative algebra that uses all of the modes? We do not know, of course. One concrete candidate is the $``$-algebra of open string field theory, defined in terms of gluing strings together. If I call this algebra $`𝒜_{st}`$, it seems plausible that $`D`$-brane charge is naturally labeled by $`K(𝒜_{st})`$. (For a manifold of very large volume compared to the string scale, I would conjecture that $`K(𝒜_{st})`$ is the same as the ordinary $`K(X)`$ of topological $`K`$-theory.) I will come back to $`K(𝒜_{st})`$ in section 3.3. First, I want to finish our survey of known applications of $`K`$-theory in string physics.
2. $`K`$-Theory And RR Fields
$`K`$-theory is relevant to understanding RR fields as well as charges \[6--8,,3\].
Naively speaking, an RR $`p`$-form field $`G_p`$ obeys a Dirac quantization law according to which, for any $`p`$-cycle $`U`$ in spacetime,
$$_U\frac{G_p}{2\pi }=\mathrm{integer}.$$
If that were the right condition, then RR fields would be classified by cohomology.
But that is not the right answer, because the actual quantization condition on RR periods is much more subtle than (2.1). There are a variety of corrections to (2.1) that involve spacetime curvature and the gauge fields on the brane, as well as self-duality and global anomalies \[9--13,,7,,8,,3\].
The answer, for Type IIB superstrings, turns out to be that RR fields are classified by $`K^1(X)`$. For our purposes, $`K^1(X)`$ can be defined as the group of components of the group of continuous maps from $`X`$ to $`U(N)`$, for any sufficiently large $`N`$.
This statement means that topological classes of RR fields on $`X`$ are classified by a map $`U:XU(N)`$ for some large $`N`$. The relation of $`G_p`$ to $`U`$ is roughly $`G_p\mathrm{Tr}(U^1dU)^p`$; here I have ignored corrections due to spacetime curvature and subtleties associated with self-duality of RR fields.
The physical meaning of $`U`$ is not clear. For Type IIA, the analog is that RR fields are classified by a $`U(N)`$ gauge bundle (for some large $`N`$) with connection $`A`$ and curvature $`F_A`$, the relation being $`G_p\mathrm{Tr}F_A^{p/2}`$. The analog for $`M`$-theory involves $`E_8`$ gauge bundles with connection \[13,,14,,3\]. Again, the physical meaning of the $`U(N)`$ or $`E_8`$ gauge fields is not clear.
The value of using $`K^1`$ to classify RR fields of Type IIB is that this gives a concise way to summarize the otherwise rather complicated quantization conditions obeyed by the RR fields. In addition, this framework is useful in describing subtle phase factors that enter in the RR partition function. In hindsight, once it is known that RR charges are classified by $`K`$-theory, one should have suspected a similar classification for RR fields. After all, RR charges produce RR fields! So the math used to classify RR charges must be similar to the math used to classify RR fields.
Just like $`K(X)`$, $`K^1(X)`$ has an analog for any noncommutative algebra $`𝒜`$. Given $`𝒜`$, we let $`𝒜_N`$ denote the group of invertible $`N\times N`$ matrices whose matrix elements are elements of $`𝒜`$. Then $`K^1(𝒜)`$ is the group of components of $`𝒜_N`$, for large $`N`$.
For example, for $`𝒜=𝒜(X)`$ the ring of complex-valued continuous functions on $`X`$, $`𝒜_N`$ is the group of maps of $`X`$ to $`GL(N,𝐂)`$. This is contractible to the group of maps of $`X`$ to $`U(N)`$, so for large $`N`$ the group of components of $`𝒜_N`$ is the same as $`K^1(X)`$, as we defined it initially.
The existence of a generalization of $`K^1(X)`$ for noncommutative rings means that the description of Type IIB RR fields by $`K^1(X)`$ in the long distance limit may be a useful starting point for stringy generalizations.
3. $`N\mathrm{}`$
In the last section, $`N`$ was a sufficiently large but finite integer. Our next task will be to describe some things that depend on setting $`N`$ equal to infinity.
Before doing so, let us recall the role of the $`N\mathrm{}`$ limit in physics. It is important in the conjectured link of gauge theory with strings; in the old matrix models that are used to give soluble examples of string theory; in the matrix model of $`M`$-theory; and in the correspondence between gravity in an asymptotically AdS spacetime and conformal field theory on the boundary.
My theme here will be to suggest that we should somehow study $`D`$-branes with $`N=\mathrm{}`$, with tachyon condensation to annihilate most of the branes and reduce to something more manageable. To motivate this, I will consider two concrete questions that seem to require taking the number of $`D`$-branes to be infinite. One question is the relation of Type IIA superstrings to $`K`$-theory; the other is the inclusion of a topologically non-trivial NS three-form field $`H`$ in the $`K`$-theory classification of RR charges and fields.
After giving the talk, I became aware of a standard problem that has already been interpreted in terms of an infinite number of $`D9`$ and $`\overline{D9}`$-branes with tachyon condensation to something manageable (see , sections 2.2, 4, and 5). This is a problem involving a $`D5`$-brane probe of a $`D5`$-$`D9`$ system.
3.1. Type IIA
For Type IIB superstrings, we used $`K(X)`$ to classify RR charges, and $`K^1(X)`$ to classify RR fields.
The $`T`$-dual statement is that for Type IIA, $`K^1(X)`$ should classify RR charges, and $`K(X)`$ should classify RR fields. (By Bott periodicity, $`K^{i+2}(X)=K^i(X)`$, so the only $`K`$-groups of $`X`$ are $`K^0(X)`$, which we have called simply $`K(X)`$, and $`K^1(X)`$.)
The most concrete and natural attempt to explain in general why $`K^1(X)`$ classifies RR charges for Type IIA is that of Horava . The starting point here is to consider a system of $`N`$ unstable $`D9`$-branes of Type IIA. The branes support a $`U(N)`$ gauge field and a tachyon field $`T`$ in the adjoint representation of $`U(N)`$. There is a symmetry $`TT`$.
The effective potential for the tachyon field is believed to have the general form
$$V(T)=\frac{1}{g_{st}}\mathrm{Tr}F(T),$$
where the function $`F(T)`$ is non-negative and, after scaling $`T`$ correctly, vanishes precisely if $`T=\pm 1`$. Hence $`V(T)`$ is minimized if and only if every eigenvalue of $`T`$ is $`\pm 1`$.
It was argued in that, in flat $`𝐑^{10}`$, one can make supersymmetric $`Dp`$-branes (for even $`p`$) as solitons of $`T`$. For example, to make a $`D6`$-brane, we set $`N=2`$. Let $`\stackrel{}{x}`$ be the three coordinates normal to the $`D6`$-brane, and set
$$T=\frac{\stackrel{}{\sigma }\stackrel{}{x}}{|x|}f(|x|),$$
where $`f(r)r`$ for small $`r`$, and $`f(r)\mathrm{}`$ for $`r\mathrm{}`$. So for $`|x|\mathrm{}`$, the eigenvalues of $`T`$ are everywhere $`\pm 1`$. Near $`x=0`$, there is a topological knot that we interpret as the $`D6`$-brane.
In flat $`𝐑^{10}`$, one can similarly make $`Dp`$-branes for other even $`p`$. But on a general spacetime, this does not work for arbitrary $`Dp`$-branes unless we set $`N=\mathrm{}`$. The problem is most obvious if $`X`$, or at least its spatial part, is compact. The tachyon field $`T`$, being adjoint-valued, maps $`X`$ to the Lie algebra of $`U(N)`$; since the Lie algebra is contractible, $`T`$ carries no topology. So a map from $`X`$ to the Lie algebra does not represent an element of $`K^1(X)`$; indeed, it does not carry topological information at all.
To define an element of $`K^1(X)`$, we need the group, not the Lie algebra; a map $`U:XU(N)`$ does the job, as I have stated before.
Amazingly, as Atiyah and Singer showed long ago, we get back the right topology from the Lie algebra if we set $`N=\mathrm{}`$! We have to interpret $`U(\mathrm{})`$ to be the unitary group $`𝒰`$ of a Hilbert space $``$ of countably infinite dimension. (Such a Hilbert space is also called a separable Hilbert space.) We interpret the $`N=\mathrm{}`$ analog of the space of hermitian $`N\times N`$ matrices to be the space of bounded self-adjoint operators $`T`$ on $``$ whose spectrum is as follows: there are infinitely many positive eigenvalues and infinitely many negative ones, and zero is not an accumulation point of the spectrum. (The last condition makes $`T`$ a “Fredholm operator.”) Physically, $`T`$ should be required to obey these conditions, since they are needed to make the energy and the $`D8`$-brane charge finite. In fact, to make the energy finite, almost all the eigenvalues of $`T`$ are very close to $`\pm 1`$. Anyway, with these conditions imposed on $`T`$, it turns out that the space of $`T`$’s has the same topology as that of $`U(N)`$ for large $`N`$.
So we can use tachyon condensation on a system of $`D9`$-branes to describe RR charges for Type IIA. But we have to start with infinitely many $`D9`$-branes, which then undergo tachyon condensation down to a configuration of finite energy.
Obstruction To Finite $`N`$
Let us now explain in more concrete terms the obstruction to making $`D`$-branes in this way for finite $`N`$, and how it vanishes for $`N=\mathrm{}`$.
Let us go back to the example of a $`D6`$-brane constructed with 2 $`D9`$-branes. We took the transverse directions to be a copy of $`𝐑^3`$, and the tachyon field to be
$$T=\frac{\stackrel{}{\sigma }\stackrel{}{x}}{|x|}f(|x|).$$
If we try to compactify the transverse directions to $`𝐒^3`$, we run into trouble because $`T`$ is not constant at infinity. The conjugacy class of $`T`$ is constant at infinity – the eigenvalues of $`T`$ are everywhere $`1`$ and $`1`$ – but $`T`$ itself is not constant.
Moreover, $`T`$ is not homotopic to a constant at infinity. If $`T`$ were homotopic to a constant near infinity, we would deform it to be constant and then extend it over $`𝐒^3`$. But it is not homotopic to a constant.
The basic obstruction to making $`T`$ constant at infinity is the “magnetic charge.” Let $`𝐒^2`$ be a sphere at infinity in $`𝐑^3`$. Over $`𝐒^2`$, we can define a line bundle $`_+`$ whose fiber is the $`+1`$ eigenspace of $`T`$, and a line bundle $`_{}`$ whose fiber is the $`1`$ eigenspace of $`T`$. The line bundles $`_+`$ and $`_{}`$ are topologically nontrivial – their first Chern classes are respectively 1 and $`1`$. As long as we try to deform $`T`$ preserving the fact that its eigenvalues are 1 and $`1`$, the line bundles $`_\pm `$ are well-defined, and their first Chern classes are invariant. So the nontriviality of $`_+`$ (or $`_{}`$) prevents us from making a homotopy to constant $`T`$.
Let us add some additional “spectator” $`D9`$-branes, and see if anything changes. Suppose there are $`M=2k`$ additional branes, so that the total number of branes is $`N=2+M=2+2k`$. Let the tachyon field be $`T^{}=TU`$, where $`T`$ is as above and $`U`$ is the sum of $`k`$ copies of the matrix
$$\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
acting on the $`2k`$ additional branes.
Thus $`T^{}`$ has near infinity $`k+1`$ eigenvalues $`+1`$ and $`k+1`$ eigenvalues $`1`$. The $`+1`$ eigenspace of $`T^{}`$ is a vector bundle $`V_+`$ of first Chern class $`1`$ (since it is constructed by adding a trivial bundle to $`_+`$), and the $`1`$ eigenspace of $`T^{}`$ is similarly a vector bundle $`V_{}`$ of first Chern class $`1`$. In particular, $`V_+`$ and $`V_{}`$ are nontrivial, so we have not gained anything by adding the spectator branes: $`T^{}`$ is not homotopic to a constant, and cannot be extended over infinity. The nontriviality of $`V_+`$ is controlled by $`\pi _1(U(k+1))=𝐙`$, which is associated with the existence of a first Chern class.
Instead, what happens if we set $`k=\mathrm{}`$? To be more precise, we take the number of spectator $`D9`$-branes to be countably infinite, and assume $`T^{}=TU`$, where $`U`$ is the direct sum of countably many copies of the matrix in (3.1). We can still define the bundles $`V_+`$ and $`V_{}`$; their fibers are separable Hilbert spaces (that is, Hilbert spaces of countably infinite dimension). $`U(k+1)`$ is replaced by $`𝒰`$, the unitary group of a separable Hilbert space. Now we run into the fundamental fact (Kuiper’s theorem) that $`𝒰`$ is contractible; its homotopy groups are all zero. Thus, any bundle of separable Hilbert spaces is trivial. In particular, $`V_+`$ and $`V_{}`$ are trivial, so $`T^{}`$ is homotopic to a constant and can be extended over infinity.
So if the total number of unstable $`D9`$-branes is $`N=\mathrm{}`$, we can make a $`D6`$-brane localized at a point in $`𝐒^3`$. More generally, in view of the result of Atiyah and Singer, we can starting at $`N=\mathrm{}`$ build an arbitrary class in $`K^1(X)`$ via tachyon condensation.
In terms of applying this result to physics, there are a few issues that we should worry about. One question is simply whether it is physically sensible to start with infinitely many branes and rely on tachyon condensation to get us down to something of finite energy. We will have to leave this question to the future (but see for an example).
Quite a different question is whether the answer that we have obtained by setting $`N=\mathrm{}`$ is the right one for physics. In the field of a $`D6`$-brane that is localized at a point on $`𝐒^3`$, the equation for the RR two-form field $`G_2`$ (of which the $`D6`$-brane is a magnetic source) has no solution, since “the flux has nowhere to go.”
It seems that the situation is that $`N=\mathrm{}`$ corresponds to the correct answer in classical open string theory, where the effective action comes from worldsheets with the topology of a disc. The RR fields enter as a correction of relative order $`g_{st}`$ (the closed string coupling constant) coming from worldsheets with cylinder topology, and should be ignored in the classical approximation.
The classification of $`D`$-branes by brane creation and annihilation holds at $`g_{st}=0`$, and (as we recalled in section 1) leads for Type IIB to a classification of $`D`$-brane charge by $`K(X)`$. To get the analogous answer – namely $`K^1(X)`$ – for Type IIA via unstable $`D9`$-branes and tachyon condensation, we need to start at $`N=\mathrm{}`$.
Intuitively, in the absence of tachyon condensation, $`N=\mathrm{}`$ should correspond to $`g_{st}=0`$, since the effective expanstion parameter for open strings is $`g_{st}N`$. If $`N`$ is infinite, then prior to tachyon condensation, $`g_{st}`$ must be zero, or the quantum corrections diverge. If we want $`g_{st}`$ to be nonzero, we need tachyon condensation to reduce to an effective finite value of $`N`$.
3.2. Turning On An $`H`$-Field
A somewhat analogous problem is to consider $`D`$-branes when the Neveu-Schwarz three-form field $`H`$ is topologically nontrivial. We will carry out this discussion in Type IIB (for Type IIA, we would have to combine what follows with what we said above in the absence of the $`H`$-field).
Just as at $`H=0`$, we would like to classify $`D`$-brane states by pairs $`(E,F)`$ (where $`E`$ is a $`D9`$ state and $`F`$ is a $`\overline{D9}`$ state) subject to the usual sort of equivalence relation. But there is a problem in having a $`D9`$ state in the presence of an $`H`$-field.
In fact, when $`H`$ is topologically non-trivial, one cannot have a single $`D9`$-brane. On the $`D9`$-brane, there is a $`U(1)`$ gauge field with field strength $`F`$. The relation $`dF=H`$ shows, at the level of de Rham cohomology, that $`H`$ must be topologically trivial if a single $`D9`$-brane is present. This conclusion actually holds precisely, not just in de Rham cohomology.
There is a special case in which there is a comparatively elementary cure for this difficulty (see , section 5.3 and ). If $`H`$ is torsion, that is if there is an integer $`M>0`$ such that $`MH`$ is topologically trivial, then it is possible to have a set of $`M`$ $`D9`$-branes whose “gauge bundle” actually has structure group $`U(M)/𝐙_M`$, rather than $`U(M)`$. (The obstruction to lifting the $`U(M)/𝐙_M`$ bundle to a $`U(M)`$ bundle is determined by $`H`$.) We will call such a gauge bundle a twisted bundle. More generally, for any positive integer $`m`$, we can have $`N=mM`$ $`D9`$-branes with the structure group of the bundle being $`U(mM)/𝐙_M`$. In such a situation, $`D`$-brane charge is classified, as one would guess, by pairs $`(E,F)`$ of twisted bundles (or $`D9`$ and $`\overline{D9}`$ states) subject to the usual equivalence relation. The equivalence classes make a group $`K_H(X)`$.
If one wishes to interpret $`K_H(X)`$ as the $`K`$-theory of representations of an algebra, one must pick a particular twisted bundle $`W`$ and consider a $`D`$-brane state with boundary conditions determined by $`W`$. The $`W`$-$`W`$ open strings transform in the adjoint representation, so the gauge parameters of the zero mode sector of the open strings are sections of $`W\overline{W}`$. Notice that although $`W`$ is a twisted bundle (with structure group $`U(M)/𝐙_M`$ rather than $`U(M)`$), $`W\overline{W}`$ is an ordinary bundle, since the center acts trivially in the adjoint representation.
The sections of $`W\overline{W}`$ form an algebra, defined as follows: if $`s^i_j`$ and $`t^k_l`$ are sections of $`W\overline{W}`$ – where the upper and lower indices are respectively $`W`$\- and $`\overline{W}`$-valued – then their product is $`(st)^i{}_{l}{}^{}=_ks^i{}_{k}{}^{}t_{}^{k}_l`$. This is the algebra $`𝒜_W(X)`$ of all endomorphisms or linear transformations of the bundle $`W`$. The algebra of open string field theory, for $`W`$-$`W`$ open strings, reduces to $`𝒜_W`$ if one looks only at the zero modes of the strings. This is a sensible approximation at low energies in a limit in which $`X`$ is very large compared to the string scale.
If $`H`$ is zero and $`W`$ is a trivial rank one complex bundle, then $`𝒜_W(X)`$ is our friend $`𝒜(X)`$. If $`H`$ is zero and $`W`$ is a trivial rank $`N`$ complex bundle, then including $`W`$ means simply that there are $`N\times N`$ Chan-Paton matrices everywhere. So in this case, $`𝒜_W(X)=𝒜(X)M_N`$, where $`M_N`$ is the algebra of $`N\times N`$ complex-valued matrices. In general, whatever $`H`$ is, $`W`$ is always trivial locally, so locally $`𝒜_W(X)`$ is isomorphic to $`𝒜(X)M_N`$.
A twisted bundle is equivalent to an $`𝒜_W`$-module, and the group $`K_H(X)`$ of pairs $`(E,F)`$ of twisted bundles (modulo the usual equivalence) coincides with $`K(𝒜_W)`$, the $`K`$-group of $`𝒜_W`$-modules. This will be explained in section 4. This assertion leads to an immediate puzzle; $`K_H(X)`$ as defined in terms of pairs $`(E,F)`$ of twisted bundles is manifestly independent of $`W`$ while $`K(𝒜_W)`$ appears to depend on $`W`$. Indeed, given any two distinct twisted bundles $`W`$ and $`W^{}`$, the corresponding algebras $`𝒜_W`$ and $`𝒜_W^{}`$ are distinct, but are “Morita-equivalent.” This concept is explained in section 4, where we also show that the Morita equivalence implies that $`K(𝒜_W)=K(𝒜_W^{})`$.
So far, we have only considered the case that $`H`$ is torsion. A typical example, important in the AdS/CFT correspondence, is the spacetime $`X=\mathrm{AdS}_5\times \mathrm{RP}^5`$, where a torsion $`H`$-field on $`\mathrm{RP}^5`$ is used to describe $`Sp(n)`$ rather than $`SO(2n)`$ gauge theory in the boundary CFT.
However, in most physical applications, $`H`$ is not torsion. In that case, we must somehow take a large $`M`$ limit of what has been said above. The right way to do this has been shown by Bouwknegt and Mathai and Atiyah and Segal . The suitable large $`M`$ limit of $`U(M)/𝐙_M`$ is $`PU()=U()/U(1)`$. In other words, for $`M=\mathrm{}`$, one replaces $`U(M)`$ by the unitary group $`U()`$ of a separable Hilbert space $``$; and one replaces $`𝐙_M`$ by $`U(1)`$. This means, in particular, that when $`H`$ is not torsion, one cannot have a finite set of $`D9`$\- or $`\overline{D9}`$-branes, but one can have an infinite set, with a suitable infinite rank twisted gauge bundle $`E`$ or $`F`$. Then $`D`$-brane charge is classified by the group $`K_H`$ of pairs $`(E,F)`$ modulo the usual equivalence relation. A detailed explanation can be found in . Kuiper’s theorem – the contractibility of $`𝒰=U()`$ – plays an important role, as it did in section 3.1.
This construction, in the $`M=\mathrm{}`$ limit, has the beautiful property, explained in , that the noncommutative algebra whose $`K`$-group is $`K_H`$ is unique, independent of any arbitrary choice of twisted bundle $`W`$ or $`W^{}`$. This really depends on the number of $`D9`$ and $`\overline{D9}`$ branes being infinite.
3.3. Stringy Generalization?
At the risk of going out on a somewhat shaky limb, I will now try to propose a stringy generalization of some of this.
We start with a closed string background and try to form an open string algebra. To do so, we must pick an open string boundary condition – call it $`\alpha `$. Then, in open string field theory, the $`\alpha `$-$`\alpha `$ open strings form an algebra $`𝒜_\alpha `$.
This algebra is not unique; we could instead pick another boundary condition $`\beta `$ and define another algebra $`𝒜_\beta `$.
The example of strings in a background $`H`$-field that is torsion, and the associated algebras $`𝒜_W`$, $`𝒜_W^{}`$, suggests the nature of the relation between $`𝒜_\alpha `$ and $`𝒜_\beta `$: they are different algebras, but are Morita-equivalent and hence have the same $`K`$-theory.
The assertion that $`𝒜_\alpha 𝒜_\beta `$, if true, is quite troublesome. It is a sharp statement of the lack of manifest background independence of open string field theory. (Background independence in this context means independence of the open string background; of course, classical open string field theory depends on a closed string background in which the open strings propagate.) It means that the formalism depends on which open string background one uses in setting up the theory.
The example of open strings in an $`H`$-field suggests a cure: take the number of $`D9`$\- and $`\overline{D9}`$-branes to be infinite.
Let $`m\alpha `$ and $`m\beta `$ denote $`m`$ copies of $`\alpha `$ or $`\beta `$ (that is, $`\alpha `$ or $`\beta `$ supplemented with $`U(m)`$-valued Chan-Paton factors). Then the conjecture is that $`𝒜_{m\alpha }𝒜_{m\beta }`$ for any finite $`m`$, but that they are equal for $`m=\mathrm{}`$. Actually, the statement about what happens for $`m=\mathrm{}`$ can be formulated much more precisely by analogy with statements in . Let $`𝒦`$ be the algebra of compact operators in a separable Hilbert space $``$. Then the conjecture is that for any $`\alpha `$ and $`\beta `$,
$$𝒜_\alpha 𝒦=𝒜_\beta 𝒦.$$
The main evidence for the conjecture is that, according to , the corresponding statement ($`𝒜_W𝒜_W^{}`$ but $`𝒜_W𝒦=𝒜_W^{}𝒦`$) holds for the zero mode algebra in the presence of a torsion $`H`$-field. In addition, if it is necessary to tensor with $`𝒦`$ before the algebras become isomorphic, this helps explain why background independence in open string field theory is so hard to understand; in other words it helps explain why the theory constructed with classical solution $`\alpha `$ and algebra $`𝒜_\alpha `$ looks different from the theory constructed with classical solution $`\beta `$ and algebra $`𝒜_\beta `$. However, the evidence for the conjecture is quite limited; it may be that whenever $`\alpha `$ is continuously connected to $`\beta `$, $`𝒜_\alpha =𝒜_\beta `$, and that the difficulty in understanding background independence in open string field theory is “just” a technical problem.
If the conjecture is right, one would suspect that to get a greater degree of background independence in open string field theory, we should start with infinitely many $`D9`$-branes and rely on tachyon condensation to get us down to something reasonable.
For this to be useful, we would need a description of $`𝒜_\alpha 𝒦`$ much simpler and more incisive than any description we have today for $`𝒜_\alpha `$ or $`𝒜_\beta `$. In such a hypothetical new description, the BRST operator $`Q`$ might be harder to describe. If that happened, we would have to take our lumps!
4. Mathematical Details
This concluding section will be devoted to explaining a few of the mathematical points that we have skimmed over so far.
Let $`𝒜`$ be a ring – in fact, let us momentarily assume that $`𝒜`$ is commutative and associative. An $`𝒜`$-module $`M`$ is called “free” if it is a direct sum of copies of $`𝒜`$: $`M=𝒜𝒜\mathrm{}𝒜`$. An $`𝒜`$-module $`M`$ is called “projective” if there is another $`𝒜`$-module $`M^{}`$ such that $`MM^{}`$ is free.
Let us consider what these definitions mean in case $`𝒜=𝒜(X)`$ is the ring of continuous complex-valued functions on a manifold $`X`$. We can think of $`𝒜`$ as the space of sections of a trivial complex line bundle $`𝒪`$. Given a trivial rank $`n`$ complex vector bundle $`V=𝒪𝒪\mathrm{}𝒪`$, the space of sections of $`V`$ is a free module $`M`$ which is the sum of $`n`$ copies of $`𝒜`$.
Now suppose we are given any complex vector bundle $`E`$ over $`X`$. The space of sections of $`E`$ is an $`𝒜`$-module $`M(E)`$ (given a function $`a𝒜(X)`$ and a section $`m`$ of $`E`$, we simply define $`a(m)`$ to be the product $`am`$).<sup>2</sup> In the informal spirit of our discussion, I have generally used the same name $`E`$ for a bundle and the corresponding module. In this paragraph only, I distinguish $`E`$ from $`M(E)`$ in the notation, to facilitate the statement of the Serre-Swan theorem. $`M(E)`$ is a free module if and only if $`E`$ is a trivial vector bundle. But $`M(E)`$ is always projective. Indeed, there is always a “complementary” vector bundle $`F`$ such that $`EF`$ is trivial, and hence $`M(E)M(F)=M(EF)`$ is free. Conversely (by a theorem of Serre and Swan), every projective $`𝒜(X)`$-module is $`M(E)`$ for some complex vector bundle $`E`$ over $`X`$.
For any ring $`𝒜`$, the projective modules form a semigroup: if $`E`$ and $`E^{}`$ are projective modules, so is $`EE^{}`$. Given any semigroup (with an addition operation that we will write as $``$), there is a canonical way to form a group. This is done the same way that one builds the group of integers, starting with the semigroup of positive integers. An element of the group is a pair $`(E,F)`$, with $`E`$ and $`F`$ elements of the semigroup (so in our situation, $`E`$ and $`F`$ are projective modules) and subject to the equivalence relation $`(E,F)(EG,FG)`$ for any element $`G`$ in the semigroup. Pairs are added by $`(E,F)+(E^{},F^{})=(EE^{},FF^{})`$. The equivalence classes form a group (the zero element or additive identity is $`0=(G,G)`$ for any $`G`$, and the additive inverse of $`(E,F)`$ is $`(F,E)`$, since $`(E,F)+(F,E)=(EF,EF)=0`$). In case one starts with the semigroup of projective modules for a ring $`𝒜`$, the group formed this way is called $`K(𝒜)`$.
For $`𝒜=𝒜(X)`$, we have in section 1 defined $`K(𝒜)`$ just in terms of vector bundles or in other words $`D9`$\- and $`\overline{D9}`$-brane configurations. Now we see how to describe this restriction in a way that has broader validity: the ninebranes correspond to projective modules, and the definition of $`K(X)`$ in terms of ninebranes is a special case of the definition of $`K(𝒜)`$ for any ring $`𝒜`$ in terms of projective modules.
Now let us drop the assumption that $`𝒜`$ is commutative (but keep associativity). If $`𝒜`$ is a noncommutative ring, one can canonically associate with it a second ring $`𝒜^{op}`$ (called the opposite ring) whose elements are in one to one correspondence with elements of $`𝒜`$, but which are multiplied in the opposite order. Thus, for every element $`a𝒜`$, there is a corresponding $`a^{op}𝒜^{op}`$, with the multiplication law of $`𝒜^{op}`$ being $`a^{op}b^{op}=(ba)^{op}`$.
For a noncommutative ring, we must distinguish two types of modules, left-modules and right-modules. $`𝒜`$ acts on the left on a left-module $`M`$, the basic axiom being that for $`a,a^{}𝒜`$ and $`mM`$, one has $`(aa^{})(m)=a(a^{}(m))`$. $`𝒜`$ acts on the right on a right-module $`M`$, the basic axiom being that $`m(aa^{})=(ma)a^{}`$. A left-module of $`𝒜`$ is the same as a right-module of $`𝒜^{op}`$, and vice-versa.
The definition of projective module is the same as in the commutative case: a left or right $`𝒜`$-module $`M`$ is projective if there is another left or right $`𝒜`$-module $`N`$ such that $`MN`$ is a free left or right $`𝒜`$-module. From the semigroup of projective left $`𝒜`$-modules, we form a group called $`K(𝒜)`$. We can also form a semigroup from the projective right $`𝒜`$-modules; it is the same as the semigroup of projective left $`𝒜^{op}`$-modules, so the associated abelian group is $`K(𝒜^{op})`$. For the rings we are interested in, there is generally an operation of complex (or hermitian) conjugation that maps $`𝒜`$ to $`𝒜^{op}`$. Consequently, the $`K`$-groups formed using left or right $`𝒜`$-modules are equivalent to each other under complex conjugation.
Given two rings $`𝒜`$ and $``$, a bimodule (or $`(𝒜,)`$-bimodule) is a group $`M`$ which is simultaneously a left $`𝒜`$-module and a right $``$-module, with the left action of $`𝒜`$ and the right action of $``$ commuting. For $`a𝒜`$, $`b`$, and $`mM`$, the left action of $`𝒜`$ is denoted $`a(m)=am`$, the right action of $``$ is denoted $`b(m)=mb`$, and since we assume that $`𝒜`$ and $``$ commute, we have $`(am)b=a(mb)`$; we abbreviate these expressions as $`amb`$.
We can now define a very important relationship between rings, called Morita equivalence. It was first exploited in connection with $`D`$-branes in ; for further developments see . Let $`M`$ be an $`(𝒜,)`$-bimodule. The choice of $`M`$ enables us to define a map from left $``$-modules to left $`𝒜`$-modules: given a left $``$-module $`W`$, we map it to the left $`𝒜`$-module $`W^{}=M_{}W`$. Similarly, if we are given a $`(,𝒜)`$-bimodule $`N`$, we can map a left $`𝒜`$-module $`W^{}`$ to a left $``$-module $`W^{\prime \prime }`$ by $`W^{\prime \prime }=N_𝒜W^{}`$. If these two operations are inverse to each other, in other words if one has $`W^{\prime \prime }=W`$ for all $`W`$, we say that $`𝒜`$ and $``$ are Morita-equivalent. In this case, the semigroup of left $`𝒜`$-modules is isomorphic to the semigroup of left $``$-modules, and $`K(𝒜)=K()`$.
The meaning of this relationship in open string field theory was explained in , section 6.4. For any open string boundary condition $`\alpha `$, one defines a ring $`𝒜_\alpha `$ consisting of the $`\alpha `$-$`\alpha `$ open strings. For any other boundary condition $`\gamma `$, the $`\alpha `$-$`\gamma `$ open strings form a left $`𝒜_\alpha `$-module $`M_{\alpha \gamma }`$, and the $`\gamma `$-$`\alpha `$ open strings form a right $`𝒜_\alpha `$-module $`M_{\gamma \alpha }`$. Presumably these are in a suitable sense a basis of left and right $`𝒜_\alpha `$-modules. Now consider any other open string boundary condition $`\beta `$ and the associated algebra $`𝒜_\beta `$. One forms left and right $`𝒜_\beta `$ modules $`M_{\beta \gamma }`$ and $`M_{\gamma \beta }`$ in the same way, using $`\beta `$-$`\gamma `$ and $`\gamma `$-$`\beta `$ open strings. Obviously, there is a natural association between left $`𝒜_\alpha `$ modules and left $`𝒜_\beta `$ modules, namely
$$M_{\alpha \gamma }M_{\beta \gamma }.$$
So $`K(𝒜_\alpha )=K(𝒜_\beta )`$ for all $`\alpha `$ and $`\beta `$. We established this last fact directly without talking about Morita equivalence, but in fact, it seems that the relation between $`𝒜_\alpha `$ and $`𝒜_\beta `$ is that they are Morita-equivalent. Indeed, one has natural $`(𝒜_\alpha ,𝒜_\beta )`$ and $`(𝒜_\beta ,𝒜_\alpha )`$ bimodules $`M_{\alpha \beta }`$ and $`M_{\beta \alpha }`$ constructed from the $`\alpha `$-$`\beta `$ and $`\beta `$-$`\alpha `$ open strings; it seems very likely that these bimodules have the right properties to establish a Morita equivalence between $`𝒜_\alpha `$ and $`𝒜_\beta `$.
Now let us consider a more elementary version of this, with the goal of adding some details to the discussion in section 3.2. Let $`M_n`$ be the algebra of $`n\times n`$ complex matrices. The most obvious left $`M_n`$-module is the space $`W`$ of $`n`$-component column vectors. This is a projective module because, as a left module, $`M_n`$ consists of $`n`$ copies of $`W`$ ($`M_n`$ consists of $`n\times n`$ matrices; in the left action of $`M_n`$ on itself, $`M_n`$ acts separately on each of the $`n`$ columns making up the $`n\times n`$ matrix, so these comprise $`n`$ copies of $`W`$). So if $`W^{}`$ is the sum of $`n1`$ copies of $`W`$, then $`WW^{}`$ is free, and thus $`W`$ is projective. It can be shown that every projective left $`M_n`$ module is the sum of copies of $`W`$. So the semigroup of projective left $`M_n`$ modules is isomorphic to the semigroup of non-negative integers, and $`K(M_n)=𝐙`$.
Since $`K(M_n)`$ is thus independent of $`n`$, one might wonder if the $`M_n`$’s of different $`n`$ are Morita-equivalent. This is indeed so. For any positive integers $`n`$ and $`k`$, the most obvious $`(M_n,M_k)`$ bimodule is the space $`C_{n,k}`$ of $`n\times k`$ matrices, with $`M_n`$ acting on the left and $`M_k`$ on the right. Likewise, the most obvious $`(M_k,M_n)`$ bimodule is $`C_{k,n}`$. It can be shown that these bimodules establish the Morita equivalence between $`M_n`$ and $`M_k`$.
Now, we reconsider the torsion $`H`$-field that was discussed in section 3.2. For any twisted bundle $`W`$, there is an algebra $`𝒜_W(X)`$ generated by the ground states of the $`W`$-$`W`$ open strings. An element of $`𝒜_W(X)`$ is a section of $`W\overline{W}`$. Given any two different twisted bundles $`W`$ and $`W^{}`$, we can try to make a Morita equivalence between $`𝒜_W(X)`$ and $`𝒜_W^{}(X)`$ by specializing to this situation the more abstract discussion of general open string algebras $`𝒜_\alpha `$ and $`𝒜_\beta `$ that was given above. This means that we should find $`(𝒜_W(X),𝒜_W^{}(X))`$ and $`(𝒜_W^{}(X),𝒜_W(X))`$ bimodules by looking at the ground states of the $`W`$-$`W^{}`$ and $`W^{}`$-$`W`$ open strings. The ground states of the $`W`$-$`W^{}`$ open strings are sections of $`W\overline{W^{}}`$; the space of these sections is a $`(𝒜_W(X),𝒜_W^{}(X))`$ bimodule $`M`$. Likewise, the space of ground states of the $`W^{}`$-$`W`$ open strings is a $`(𝒜_W^{}(X),𝒜_W(X))`$ bimodule $`N`$ consisting of the sections of $`W^{}\overline{W}`$. The bimodules $`M`$ and $`N`$ do indeed establish a Morita equivalence between $`𝒜_W(X)`$ and $`𝒜_W^{}(X)`$. Indeed, in proving this one can work locally on $`X`$. Locally, $`W`$ and $`W^{}`$ are trivial, and (if $`n`$ and $`k`$ are the ranks of $`W`$ and $`W^{}`$) we have locally as noted in section 3.2, $`𝒜_W(X)=𝒜(X)M_n`$, $`𝒜_W^{}(X)=𝒜(X)M_k`$. The same computation that gives a Morita equivalence of $`M_n`$ with $`M_k`$ shows that $`𝒜_W(X)`$ is Morita equivalent to $`𝒜_W^{}(X)`$.
A projective right $`𝒜_W(X)`$ module is the space of sections of $`E\overline{W}`$, for any twisted bundle $`E`$. (If $`s^\alpha _j`$ is such a section – where $`\alpha `$ is an $`E`$ index and $`j`$ a $`\overline{W}`$ index – and $`w^k_l`$ is a section of $`W\overline{W}`$, then the product of $`s`$ and $`w`$ is defined by $`(sw)^\alpha {}_{j}{}^{}=_ks^\alpha {}_{k}{}^{}w_{}^{k}_j`$.) So projective right $`𝒜_W(X)`$ modules are in natural one-to-one correspondence with twisted bundles $`E`$, and the $`K`$-group of projective right $`𝒜_W(X)`$ modules is the group of pairs $`(E,F)`$, subject to the usual equivalence relation. The Morita equivalence of $`𝒜_W`$ and $`𝒜_W^{}`$, for any two twisted bundles $`W`$ and $`W^{}`$ with the same $`H`$, is just a fancy way of describing the natural map between $`𝒜_W`$-modules and $`𝒜_W^{}`$-modules that comes from the correspondence
$$E\overline{W}E\overline{W^{}}.$$
This one-to-one map between right $`𝒜_W`$-modules and right $`𝒜_W^{}`$-modules (or the corresponding map for left modules) leads to the conclusion that $`K(𝒜_W)=K(𝒜_W^{})`$.
This work was supported in part by NSF Grant PHY-9513835 and the Caltech Discovery Fund. I am grateful to K. Hori and G. B. Segal for comments and explanations.
References
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# Single superpartner production at Tevatron Run II
## 1 Introduction
In the Minimal Supersymmetric Standard Model (MSSM), the supersymmetric (SUSY) particles must be produced in pairs. The phase space is largely suppressed in pair production of SUSY particles due to the important masses of the superpartners. The R-parity violating ($`\mathit{}_p`$ ) extension of the MSSM contains the following additional terms in the superpotential, which are trilinear in the quarks and leptons superfields,
$`W_{\mathit{}_p\text{ }}={\displaystyle \underset{i,j,k}{}}\left({\displaystyle \frac{1}{2}}\lambda _{ijk}L_iL_jE_k^c+\lambda _{ijk}^{}L_iQ_jD_k^c+{\displaystyle \frac{1}{2}}\lambda _{ijk}^{\prime \prime }U_i^cD_j^cD_k^c\right),`$ (1)
where $`i,j,k`$ are flavour indices. These $`\mathit{}_p`$ couplings offer the opportunity to produce the scalar supersymmetric particles as resonances . Although the $`\mathit{}_p`$ coupling constants are severely constrained by the low-energy experimental bounds , the resonant superpartner production reaches high cross sections both at leptonic and hadronic colliders.
The resonant production of SUSY particle has another interest: since its cross section is proportional to a power $`2`$ of the relevant $`\mathit{}_p`$ coupling, this reaction would allow an easier determination of the $`\mathit{}_p`$ couplings than the pair production provided the $`\mathit{}_p`$ coupling is large enough. As a matter of fact in the pair production study, the sensitivity on the $`\mathit{}_p`$ couplings is mainly provided by the displaced vertex analysis of the Lightest Supersymmetric Particle (LSP) decay which is difficult experimentally, especially at hadronic colliders.
Neither the Grand Unified Theories (GUT), the string theories nor the study of the discrete gauge symmetries give a strong theoretical argument in favor of the R-parity violating or R-parity conserving scenarios . Hence, the resonant production of SUSY particle through $`\mathit{}_p`$ couplings is an attractive possibility which must be considered in the phenomenology of supersymmetry.
The hadronic colliders have an advantage in detecting new particles resonance. Indeed, due to the wide energy distribution of the colliding partons, the resonance can be probed in a wide range of the new particle mass. This is in contrast with the leptonic colliders for which the center of mass energy must be fine-tuned in order to discover new narrow width resonances.
At hadronic colliders, either a slepton or a squark can be produced at the resonance respectively through a $`\lambda ^{}`$ or a $`\lambda ^{\prime \prime }`$ coupling constant. In the hypothesis of a single dominant $`\mathit{}_p`$ coupling constant, the resonant scalar particle can decay through the same $`\mathit{}_p`$ coupling as in the production, leading to a two quark final state for the hard process . In the case where both $`\lambda ^{}`$ and $`\lambda `$ couplings are non-vanishing, the slepton produced via $`\lambda ^{}`$ can decay through $`\lambda `$ giving rise to the same final state as in Drell-Yan process, namely two leptons . However, for reasonable values of the $`\mathit{}_p`$ coupling constants, the decays of the resonant scalar particle via gauge interactions are typically dominant if kinematically allowed .
The main decay of the resonant scalar particle through gauge interactions is the decay into its Standard Model partner plus a gaugino. Indeed, in the case where the resonant scalar particle is a squark, it is produced through $`\lambda ^{\prime \prime }`$ interactions so that it must be a Right squark $`\stackrel{~}{q}_R`$ and thus it cannot decay into the $`W^\pm `$-boson, which is the only other possible decay channel via gauge interactions. Besides, in the case where the resonant scalar particle is a slepton, it is a Left slepton produced via a $`\lambda ^{}`$ coupling but it cannot generally decay as $`\stackrel{~}{l}_L^\pm W^\pm \stackrel{~}{\nu }_L`$ or as $`\stackrel{~}{\nu }_LW^\pm \stackrel{~}{l}_L^{}`$. The reason is that in most of the SUSY models, as for example the supergravity or the gauge mediated models, the mass difference between the Left charged slepton and the Left sneutrino is due to the D-terms so that it is fixed by the relation $`m_{\stackrel{~}{l}_L^\pm }^2m_{\stackrel{~}{\nu }_L}^2=\mathrm{cos}2\beta M_W^2`$ and thus it does not exceed the $`W^\pm `$-boson mass. Nevertheless, we note that in the large $`\mathrm{tan}\beta `$ scenario, a resonant scalar particle of the third generation can generally decay into the $`W^\pm `$-boson due to the large mixing in the third family sfermion sector. For instance, in the SUGRA model with a large $`\mathrm{tan}\beta `$ a tau-sneutrino produced at the resonance can decay as $`\stackrel{~}{\nu }_\tau W^\pm \stackrel{~}{\tau }_1^{}`$, $`\stackrel{~}{\tau }_1^{}`$ being the lightest stau.
The resonant scalar particle production at hadronic colliders leads thus mainly to the single gaugino production, in case where the decay of the relevant scalar particle into gaugino is kinematically allowed. In this paper, we study the single gaugino productions at Tevatron Run II. The single gaugino productions at hadronic colliders were first studied in . Later, studies on the single neutralino and single chargino productions at Tevatron have been performed. The single neutralino and single chargino productions have also been considered in the context of physics at LHC. In the present article, we also study the single superpartner productions at Tevatron Run II which occur via $`22body`$ processes and do not receive contributions from resonant SUSY particle productions.
The singly produced superpartner initiates a cascade decay ended typically by the $`\mathit{}_p`$ decay of the LSP. In case of a single dominant $`\lambda ^{\prime \prime }`$ coupling constant, the LSP decays into quarks so that this cascade decay leads to multijet final states having a large QCD background . Nevertheless, if some leptonic decays, as for instance $`\stackrel{~}{\chi }^\pm l^\pm \nu \stackrel{~}{\chi }^0`$, $`\stackrel{~}{\chi }^\pm `$ being the chargino and $`\stackrel{~}{\chi }^0`$ the neutralino, enter the chain reaction, clearer leptonic signatures can be investigated . In contrast, in the hypothesis of a single dominant $`\lambda ^{}`$ coupling constant, the LSP decay into charged leptons naturally favors leptonic signatures . We will thus study the single superpartner production reaction at Tevatron Run II within the scenario of a single dominant $`\lambda _{ijk}^{}`$ coupling constant.
In section 2, we define our theoretical framework. In section 3, we present the values of the cross sections for the various single superpartner productions via $`\lambda _{ijk}^{}`$ at Tevatron Run II and we discuss the interesting multileptonic signatures that these processes can generate. In section 4, we analyse the three lepton signature induced by the single chargino production. In section 5, we study the like sign dilepton final state generated by the single neutralino and chargino productions.
## 2 Theoretical framework
Our framework throughout this paper will be the so-called minimal supergravity model (mSUGRA) which assumes the existence of a grand unified gauge theory and family universal boundary conditions on the supersymmetry breaking parameters. We choose the 5 following parameters: $`m_0`$ the universal scalars mass at the unification scale $`M_X`$, $`m_{1/2}`$ the universal gauginos mass at $`M_X`$, $`A=A_t=A_b=A_\tau `$ the trilinear Yukawa coupling at $`M_X`$, $`sign\left(\mu \right)`$ the sign of the $`\mu \left(t\right)`$ parameter ($`t=\mathrm{log}\left(M_X^2/Q^2\right)`$, $`Q`$ denoting the running scale) and $`\mathrm{tan}\beta =<H_u>/<H_d>`$ where $`<H_u>`$ and $`<H_d>`$ denote the vacuum expectation values of the Higgs fields. In this model, the higgsino mixing parameter $`\left|\mu \right|`$ is determined by the radiative electroweak symmetry breaking condition. Note also that the parameters $`m_{1/2}`$ and $`M_2\left(t\right)`$ ($`\stackrel{~}{W}`$ wino mass) are related by the solution of the one loop renormalization group equations $`m_{1/2}=\left(1\beta _at\right)M_a\left(t\right)`$ with $`\beta _a=g_X^2b_a/\left(4\pi \right)^2`$, where $`\beta _a`$ are the beta functions, $`g_X`$ is the coupling constant at $`M_X`$ and $`b_a=[3,1,11]`$, $`a=[3,2,1]`$ corresponding to the gauge group factors $`SU\left(3\right)_c,SU\left(2\right)_L,SU\left(1\right)_Y`$. We shall set the unification scale at $`M_X=\mathrm{2\; 10}^{16}GeV`$ and the running scale at the $`Z^0`$-boson mass: $`Q=m_{Z^0}`$.
We also assume the infrared fixed point hypothesis for the top quark Yukawa coupling that provides a natural explanation of a large top quark mass $`m_{top}`$. In the infrared fixed point approach, $`\mathrm{tan}\beta `$ is fixed up to the ambiguity associated with large or low $`\mathrm{tan}\beta `$ solutions. The low solution of $`\mathrm{tan}\beta `$ is fixed by the equation $`m_{top}=C\mathrm{sin}\beta `$, where $`C190210GeV`$ for $`\alpha _s\left(m_{Z^0}\right)=0.110.13`$. For instance, with a top quark mass of $`m_{top}=174.2GeV`$ , the low solution is given by $`\mathrm{tan}\beta 1.5`$. The second important effect of the infrared fixed point hypothesis is that the dependence of the electroweak symmetry breaking constraint on the $`A`$ parameter becomes weak so that $`\left|\mu \right|`$ is a known function of the $`m_0`$, $`m_{1/2}`$ and $`\mathrm{tan}\beta `$ parameters .
Finally, we consider the $`\mathit{}_p`$ extension of the mSUGRA model characterised by a single dominant $`\mathit{}_p`$ coupling constant of type $`\lambda _{ijk}^{}`$.
## 3 Single superpartner productions via $`\lambda _{ijk}^{}`$ at Tevatron Run II
### 3.1 Resonant superpartner production
At hadronic colliders, either a sneutrino ($`\stackrel{~}{\nu }`$) or a charged slepton ($`\stackrel{~}{l}`$) can be produced at the resonance via the $`\lambda _{ijk}^{}`$ coupling. As explained in Section 1, for most of the SUSY models, the slepton produced at the resonance has two possible gauge decays, namely a decay into either a chargino or a neutralino. Therefore, in the scenario of a single dominant $`\lambda _{ijk}^{}`$ coupling and for most of the SUSY models, either a chargino or a neutralino is singly produced together with either a charged lepton or a neutrino, through the resonant superpartner production at hadronic colliders. There are thus four main possible types of single superpartner production reaction involving $`\lambda _{ijk}^{}`$ at hadronic colliders which receive a contribution from resonant SUSY particle production. The diagrams associated to these four reactions are drawn in Fig.1. As can be seen in this figure, these single superpartner productions receive also some contributions from both the $`t`$ and $`u`$ channels. Note that all the single superpartner production processes drawn in Fig.1 have charge conjugated processes. We have calculated the amplitudes of the processes shown in Fig.1 and the results are given in Appendix A.
#### 3.1.1 Cross sections
In this section, we discuss the dependence of the single gaugino production cross sections on the various supersymmetric parameters. We will not assume here the radiative electroweak symmetry breaking condition in order to study the variations of the cross sections with the higgsino mixing parameter $`\mu `$.
First, we study the cross section of the single chargino production $`p\overline{p}\stackrel{~}{\chi }^+l_i^{}`$ which occurs through the $`\lambda _{ijk}^{}`$ coupling (see Fig.1(a)). The differences between the $`\stackrel{~}{\chi }^+e^{}`$, $`\stackrel{~}{\chi }^+\mu ^{}`$ and $`\stackrel{~}{\chi }^+\tau ^{}`$ production (occuring respectively through the $`\lambda _{1jk}^{}`$, $`\lambda _{2jk}^{}`$ and $`\lambda _{3jk}^{}`$ couplings with identical $`j`$ and $`k`$ indices) cross sections involve $`m_{l_i}`$ lepton mass terms (see Appendix A) and are thus negligible. The $`p\overline{p}\stackrel{~}{\chi }^+l_i^{}`$ reaction receives contributions from the $`s`$ channel sneutrino exchange and the $`t`$ and $`u`$ channels squark exchanges as shown in Fig.1. However, the $`t`$ and $`u`$ channels represent small contributions to the whole single chargino production cross section when the sneutrino exchanged in the $`s`$ channel is real, namely for $`m_{\stackrel{~}{\nu }_{iL}}>m_{\stackrel{~}{\chi }^\pm }`$. The $`t`$ and $`u`$ channels cross sections will be relevant only when the produced sneutrino is virtual since the $`s`$ channel contribution is small. In this situation the single chargino production rate is greatly reduced compared to the case where the exchanged sneutrino is produced as a resonance. Hence, The $`t`$ and $`u`$ channels do not represent important contributions to the $`\stackrel{~}{\chi }^+l_i^{}`$ production rate.
The dependence of the $`\stackrel{~}{\chi }^+l_i^{}`$ production rate on the $`A`$ coupling is weak. Indeed, the rate depends on the $`A`$ parameter only through the masses of the third generation squarks eventually exchanged in the $`t`$ and $`u`$ channels (see Fig.1). Similarly, the dependences on the $`A`$ coupling of the rates of the other single gaugino productions shown in Fig.1 are weak. Therefore, in this article we present the results for $`A=0`$. Later, we will discuss the effects of large $`A`$ couplings on the cascade decays which are similar to the effects of large $`\mathrm{tan}\beta `$ values.
$`\mathrm{tan}\beta `$ dependence: The dependence of the $`\stackrel{~}{\chi }^+l_i^{}`$ production rate on $`\mathrm{tan}\beta `$ is also weak, except for $`\mathrm{tan}\beta <10`$. This can be seen in Fig.2 where the cross section of the $`p\overline{p}\stackrel{~}{\chi }_1^+\mu ^{}`$ reaction occuring through the $`\lambda _{211}^{}`$ coupling is shown as a function of the $`\mathrm{tan}\beta `$ parameter. The choice of the $`\lambda _{211}^{}`$ coupling is motivated by the fact that the analysis in Sections 4 and 5 are explicitly made for this $`\mathit{}_p`$ coupling. In Fig.2, we have taken the $`\lambda _{211}^{}`$ value equal to its low-energy experimental bound for $`m_{\stackrel{~}{d}_R}=100GeV`$ which is $`\lambda _{211}^{}<0.09`$ .
At this stage, some remarks on the values of the cross sections presented in this section must be done. First, the single gaugino production rates must be multiplied by a factor 2 in order to take into account the charge conjugated process, which is for example in the present case $`p\overline{p}\stackrel{~}{\chi }^{}\mu ^+`$. Furthermore, the values of the cross sections for all the single gaugino productions are obtained using the CTEQ4L structure function . Choosing other parametrizations does not change significantly the results since proton structure functions in our kinematical domain in Bjorken $`x`$ are known and have been already measured. For instance, with the set of parameters $`\lambda _{211}^{}=0.09`$, $`M_2=100GeV`$, $`\mathrm{tan}\beta =1.5`$, $`m_0=300GeV`$ and $`\mu =500GeV`$, the $`\stackrel{~}{\chi }_1^+\mu ^{}`$ production cross section is $`0.503pb`$ for the CTEQ4L structure function , $`0.503pb`$ for the BEP structure function , $`0.480pb`$ for the MRS (R2) structure function and $`0.485pb`$ for the GRV LO structure function .
$`\mu `$ dependence: In Fig.3, we present the cross sections of the $`\stackrel{~}{\chi }_1^+\mu ^{}`$ and $`\stackrel{~}{\chi }_2^+\mu ^{}`$ productions as a function of the $`\mu `$ parameter. We observe in this figure the weak dependence of the cross section $`\sigma \left(p\overline{p}\stackrel{~}{\chi }_1^+\mu ^{}\right)`$ on $`\mu `$ for $`\left|\mu \right|>M_2`$. The reason is the smooth dependence of the $`\stackrel{~}{\chi }_1^\pm `$ mass on $`\mu `$ in this domain. However, the rate strongly decreases in the region $`\left|\mu \right|<M_2`$ in which the $`\stackrel{~}{\chi }_1^\pm `$ chargino is mainly composed by the higgsino. Nevertheless, the small $`\left|\mu \right|`$ domain ($`\left|\mu \right|`$ smaller than $`100GeV`$ for $`\mathrm{tan}\beta =1.41`$, $`M_2>100GeV`$, $`m_0=500GeV`$ and $`\lambda ^{}0`$) is excluded by the present experimental limits derived from the LEP data .
In contrast, the cross section $`\sigma \left(p\overline{p}\stackrel{~}{\chi }_2^+\mu ^{}\right)`$ increases in the domain $`\left|\mu \right|<M_2`$. The explanation is that the $`\stackrel{~}{\chi }_2^\pm `$ mass is enhanced as $`\left|\mu \right|`$ increases. The region in which $`\sigma \left(p\overline{p}\stackrel{~}{\chi }_2^+\mu ^{}\right)`$ becomes important is at small values of $`\left|\mu \right|`$, near the LEP limits of . We also remark in Fig.3 that the single $`\stackrel{~}{\chi }_1^+`$ production rate values remain above the single $`\stackrel{~}{\chi }_2^+`$ production rate values in all the considered range of $`\mu `$. In this figure, we also notice that the cross section is smaller when $`\mu `$ is negative. To be conservative, we will take $`\mu <0`$ in the following.
$`m_0`$ and $`M_2`$ dependences: In fact, the cross section $`\sigma \left(p\overline{p}\stackrel{~}{\chi }^+l_i^{}\right)`$ depends mainly on the $`m_0`$ and $`M_2`$ parameters. We present in Fig.4 the rate of the $`\stackrel{~}{\chi }_1^+\mu ^{}`$ production as a function of the $`m_0`$ and $`M_2`$ parameters. The rate decreases at high values of $`m_0`$ since the sneutrino becomes heavier as $`m_0`$ increases and more energetic initial partons are required in order to produce the resonant sneutrino. The decrease of the rate at large values of $`M_2`$ is due to the increase of the chargino mass and thus the reduction of the phase space factor.
In Fig.5, we show the variations of the $`\sigma \left(p\overline{p}\stackrel{~}{\chi }_1^+\mu ^{}\right)`$ cross sections with $`m_0`$ for fixed values of $`M_2`$, $`\mu `$ and $`\mathrm{tan}\beta `$. The cross sections corresponding to the $`\stackrel{~}{\chi }_1^+\mu ^{}`$ production through various $`\mathit{}_p`$ couplings of type $`\lambda _{2jk}^{}`$ are presented. In this figure, we only consider the $`\mathit{}_p`$ couplings giving the highest cross sections. The values of the considered $`\lambda _{2jk}^{}`$ couplings have been taken at their low-energy limit for a squark mass of $`100GeV`$. The rate of the $`\stackrel{~}{\chi }_2^+\mu ^{}`$ production through $`\lambda _{211}^{}`$ is also shown in this figure. We already notice that the cross section is significant for many $`\mathit{}_p`$ couplings and we will come back on this important statement in the following.
The $`\sigma \left(p\overline{p}\stackrel{~}{\chi }^+\mu ^{}\right)`$ rates decrease as $`m_0`$ increases for the same reason as in Fig.4. A decrease of the rates also occurs at small values of $`m_0`$. The reason is the following. When $`m_0`$ decreases, the $`\stackrel{~}{\nu }`$ mass is getting closer to the $`\stackrel{~}{\chi }^\pm `$ masses so that the phase space factor associated to the decay $`\stackrel{~}{\nu }_\mu \stackrel{~}{\chi }^\pm \mu ^{}`$ decreases.
We also observe that the single $`\stackrel{~}{\chi }_2^+`$ production rate is much smaller than the single $`\stackrel{~}{\chi }_1^+`$ production rate, as in Fig.3.
The differences between the $`\stackrel{~}{\chi }_1^+\mu ^{}`$ production rates occuring via the various $`\lambda _{2jk}^{}`$ couplings are explained by the different parton densities. Indeed, as shown in Fig.1 the hard process associated to the $`\stackrel{~}{\chi }_1^+\mu ^{}`$ production occuring through the $`\lambda _{2jk}^{}`$ coupling constant has a partonic initial state $`\overline{q}_jq_k`$. The $`\stackrel{~}{\chi }_1^+\mu ^{}`$ production via the $`\lambda _{211}^{}`$ coupling has first generation quarks in the initial state which provide the maximum parton density.
We now discuss the rate behaviours for the reactions $`p\overline{p}\stackrel{~}{\chi }^{}\nu _\mu `$, $`p\overline{p}\stackrel{~}{\chi }^0\mu ^{}`$ and $`p\overline{p}\stackrel{~}{\chi }^0\nu _\mu `$ which occur via $`\lambda _{211}^{}`$, in the SUSY parameter space. The dependences of these rates on the $`A`$, $`\mathrm{tan}\beta `$, $`\mu `$ and $`M_2`$ parameters are typically the same as for the $`\stackrel{~}{\chi }^+\mu ^{}`$ production rate. The variations of the $`\stackrel{~}{\chi }_1^{}\nu _\mu `$, $`\stackrel{~}{\chi }_{1,2}^0\mu ^{}`$ and $`\stackrel{~}{\chi }_1^0\nu _\mu `$ productions cross sections with the $`m_0`$ parameter are shown in Fig.6. The $`\stackrel{~}{\chi }_2^{}\nu _\mu `$, $`\stackrel{~}{\chi }_{3,4}^0\mu ^{}`$ and $`\stackrel{~}{\chi }_{3,4}^0\nu _\mu `$ production rates are comparatively negligible and thus have not been represented. We observe in this figure that the cross sections decrease at large $`m_0`$ values like the $`\stackrel{~}{\chi }^+\mu ^{}`$ production rate. However, while the single $`\stackrel{~}{\chi }_1^\pm `$ productions rates decrease at small $`m_0`$ values (see Fig.5 and Fig.6), this is not true for the single $`\stackrel{~}{\chi }_1^0`$ productions (see Fig.6). The reason is that in mSUGRA the $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{l}_{iL}`$ ($`l_i=l_i^\pm ,\nu _i`$) masses are never close enough to induce a significant decrease of the cross section associated to the reaction $`p\overline{p}\stackrel{~}{l}_{iL}\stackrel{~}{\chi }_1^0l_i`$, where $`l_i=l_i^\pm ,\nu _i`$ (see Fig.1(c)(d)), caused by a phase space factor reduction. Therefore, the resonant slepton contribution to the single $`\stackrel{~}{\chi }_1^0`$ production is not reduced at small $`m_0`$ values like the resonant slepton contribution to the single $`\stackrel{~}{\chi }_1^\pm `$ production. For the same reason, the single $`\stackrel{~}{\chi }_1^0`$ productions have much higher cross sections than the single $`\stackrel{~}{\chi }_1^\pm `$ productions in most of the mSUGRA parameter space, as illustrate Fig.5 and Fig.6. We note that in the particular case of a single dominant $`\lambda _{3jk}^{}`$ coupling constant and of large $`\mathrm{tan}\beta `$ values, the rate of the reaction $`p\overline{p}\stackrel{~}{\tau }_1^\pm \stackrel{~}{\chi }_1^0\tau ^\pm `$ (see Fig.1(d)), where $`\stackrel{~}{\tau }_1^\pm `$ is the lightest tau-slepton, can be reduced at low $`m_0`$ values since then $`m_{\stackrel{~}{\tau }_1^\pm }`$ can be closed to $`m_{\stackrel{~}{\chi }_1^0}`$ due to the large mixing occuring in the staus sector. By analysing Fig.5 and Fig.6, we also remark that the $`\stackrel{~}{\chi }^{}\nu _\mu `$ ($`\stackrel{~}{\chi }^0\mu ^{}`$) production rate is larger than the $`\stackrel{~}{\chi }^+\mu ^{}`$ ($`\stackrel{~}{\chi }^0\nu _\mu `$) one. The explanation is that in $`p\overline{p}`$ collisions the initial states of the resonant charged slepton production $`u_j\overline{d}_k,\overline{u}_jd_k`$ have higher partonic densities than the initial states of the resonant sneutrino production $`d_j\overline{d}_k,\overline{d}_jd_k`$. This phenomenon also increases the difference between the rates of the $`\stackrel{~}{\chi }_1^0\mu ^{}`$ and $`\stackrel{~}{\chi }_1^+\mu ^{}`$ productions at Tevatron.
Although the single $`\stackrel{~}{\chi }_1^\pm `$ production cross sections are smaller than the $`\stackrel{~}{\chi }_1^0`$ ones, it is interesting to study both of them since they have quite high values.
### 3.2 Non-resonant superpartner production
At hadronic colliders, the single productions of SUSY particle via $`\lambda _{ijk}^{}`$ can occur through some $`22body`$ processes which do not receive contributions from any resonant superpartner production. These non-resonant superpartner productions are (one must also add the charge conjugated processes):
* The gluino production $`\overline{u}_jd_k\stackrel{~}{g}l_i`$ via the exchange of a $`\stackrel{~}{u}_{jL}`$ ($`\stackrel{~}{d}_{kR}`$) squark in the $`t`$ ($`u`$) channel.
* The squark production $`\overline{d}_jg\stackrel{~}{d}_{kR}^{}\nu _i`$ via the exchange of a $`\stackrel{~}{d}_{kR}`$ squark ($`d_j`$ quark) in the $`t`$ ($`s`$) channel.
* The squark production $`\overline{u}_jg\stackrel{~}{d}_{kR}^{}l_i`$ via the exchange of a $`\stackrel{~}{d}_{kR}`$ squark ($`u_j`$ quark) in the $`t`$ ($`s`$) channel.
* The squark production $`d_kg\stackrel{~}{d}_{jL}\nu _i`$ via the exchange of a $`\stackrel{~}{d}_{jL}`$ squark ($`d_k`$ quark) in the $`t`$ ($`s`$) channel.
* The squark production $`d_kg\stackrel{~}{u}_{jL}l_i`$ via the exchange of a $`\stackrel{~}{u}_{jL}`$ squark ($`d_k`$ quark) in the $`t`$ ($`s`$) channel.
* The sneutrino production $`\overline{d}_jd_kZ\stackrel{~}{\nu }_{iL}`$ via the exchange of a $`d_k`$ or $`d_j`$ quark ($`\stackrel{~}{\nu }_{iL}`$ sneutrino) in the $`t`$ ($`s`$) channel.
* The charged slepton production $`\overline{u}_jd_kZ\stackrel{~}{l}_{iL}`$ via the exchange of a $`d_k`$ or $`u_j`$ quark ($`\stackrel{~}{l}_{iL}`$ slepton) in the $`t`$ ($`s`$) channel.
* The sneutrino production $`\overline{u}_jd_kW^{}\stackrel{~}{\nu }_{iL}`$ via the exchange of a $`d_j`$ quark ($`\stackrel{~}{l}_{iL}`$ sneutrino) in the $`t`$ ($`s`$) channel.
* The charged slepton production $`\overline{d}_jd_kW^+\stackrel{~}{l}_{iL}`$ via the exchange of a $`u_j`$ quark ($`\stackrel{~}{\nu }_{iL}`$ sneutrino) in the $`t`$ ($`s`$) channel.
The single gluino production cannot reach high cross sections due to the strong experimental limits on the squarks and gluinos masses which are typically about $`m_{\stackrel{~}{q}},m_{\stackrel{~}{g}}\stackrel{>}{}200GeV`$ . Indeed, the single gluino production occurs through the exchange of squarks in the $`t`$ and $`u`$ channels, as described above, so that the cross section of this production decreases as the squarks and gluinos masses increase. For the value $`m_{\stackrel{~}{q}}=m_{\stackrel{~}{g}}=250GeV`$ which is close to the experimental limits, we find the single gluino production rate $`\sigma \left(p\overline{p}\stackrel{~}{g}\mu \right)10^2pb`$ which is consistent with the results of . The cross sections given in this section are computed at a center of mass energy of $`\sqrt{s}=2TeV`$ using the version 33.18 of the COMPHEP routine with the CTEQ4m structure function and an $`\mathit{}_p`$ coupling $`\lambda _{211}^{}=0.09`$. Similarly, the single squark production cross section cannot be large: for $`m_{\stackrel{~}{q}}=250GeV`$, the rate $`\sigma \left(p\overline{p}\stackrel{~}{u}_L\mu \right)`$ is of order $`10^3pb`$. The production of a slepton together with a massive gauge boson has a small phase space factor and does not involve strong interaction couplings. The cross section of this type of reaction is thus small. For instance, with a slepton mass of $`m_{\stackrel{~}{l}}=100GeV`$ we find the cross section $`\sigma \left(p\overline{p}Z\stackrel{~}{\mu }_L\right)`$ to be of order $`10^2pb`$.
As a conclusion, the non-resonant single superpartner productions have small rates and will not be considered here. Nevertheless, some of these reactions are interesting as their cross section involves few SUSY parameters, namely only one scalar superpartner mass and one $`\mathit{}_p`$ coupling constant.
## 4 Three lepton signature analysis
### 4.1 Signal
In this section, we study the three lepton signature at Tevatron Run II generated by the single chargino production through $`\lambda _{ijk}^{}`$, $`p\overline{p}\stackrel{~}{\chi }^\pm l_i^{}`$, followed by the cascade decay, $`\stackrel{~}{\chi }^\pm \stackrel{~}{\chi }_1^0l^\pm \nu `$, $`\stackrel{~}{\chi }_1^0l_iu_j\overline{d}_k,\overline{l}_i\overline{u}_jd_k`$ (the indices $`i,j,k`$ correspond to the indices of $`\lambda _{ijk}^{}`$). In fact, the whole final state is 3 charged leptons + 2 hard jets + missing energy ($`E/`$). The two jets and the missing energy come respectively from the quarks and the neutrino produced in the cascade decay. In the mSUGRA model, which predicts the $`\stackrel{~}{\chi }_1^0`$ as the LSP in most of the parameter space, the $`p\overline{p}\stackrel{~}{\chi }^\pm l_i^{}`$ reaction is the only single gaugino production allowing the three lepton signature to be generated in a significant way. Since the $`\stackrel{~}{\chi }_1^\pm l_i^{}`$ production rate is dominant compared to the $`\stackrel{~}{\chi }_2^\pm l_i^{}`$ production rate, as discussed in Section 3.1.1, we only consider the contribution to the three lepton signature from the single lightest chargino production.
For $`m_{\stackrel{~}{\nu }},m_{\stackrel{~}{l}},m_{\stackrel{~}{q}},m_{\stackrel{~}{\chi }_2^0}>m_{\stackrel{~}{\chi }_1^\pm }`$, the branching ratio $`B\left(\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0l^\pm \nu \right)`$ is typically of order $`30\%`$ and is smaller than for the other possible decay $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0\overline{q}_pq_p^{}`$ because of the color factor.
Since in our framework the $`\stackrel{~}{\chi }_1^0`$ is the LSP, it can only decay via $`\lambda _{ijk}^{}`$, either as $`\stackrel{~}{\chi }_1^0l_iu_jd_k`$ or as $`\stackrel{~}{\chi }_1^0\nu _id_jd_k`$, with a branching ratio $`B\left(\stackrel{~}{\chi }_1^0l_iu_jd_k\right)`$ ranging between $`40\%`$ and $`70\%`$.
The three lepton signature is particularly attractive at hadronic colliders because of the possibility to reduce the associated Standard Model background. In Section 4.2 we describe this Standard Model background and in Section 4.4 we show how it can be reduced.
### 4.2 Standard Model background of the 3 lepton signature at Tevatron
The first source of Standard Model background for the three leptons final state is the top quark pair production $`q\overline{q}t\overline{t}`$ or $`ggt\overline{t}`$. Since the top quark life time is smaller than its hadronisation time, the top decays and its main channel is the decay into a $`W`$ gauge boson and a bottom quark as $`tbW`$. The $`t\overline{t}`$ production can thus give rise to a $`3l`$ final state if the $`W`$ bosons and one of the b-quarks undergo leptonic decays simultaneously. The cross section, calculated at leading order with PYTHIA using the CTEQ2L structure function, times the branching fraction is $`\sigma \left(p\overline{p}t\overline{t}\right)\times B^2\left(Wl_p\nu _p\right)863fb`$ ($`704fb`$) with $`p=1,2,3`$ at $`\sqrt{s}=2TeV`$ for a top quark mass of $`m_{top}=170GeV`$ ($`175GeV`$).
The other major source of Standard Model background is the $`W^\pm Z^0`$ production followed by the leptonic decays of the gauge bosons, namely $`Wl\nu `$ and $`Zl\overline{l}`$. The value for the cross section times the branching ratios is $`\sigma \left(p\overline{p}WZ\right)\times B\left(Wl_p\nu _p\right)\times B\left(Zl_p\overline{l}_p\right)82fb`$ ($`p=1,2,3`$) at leading order with a center of mass energy of $`\sqrt{s}=2TeV`$.
The $`W^\pm Z^0`$ production gives also a small contribution to the 3 leptons background through the decays: $`Wbu_p`$ and $`Zb\overline{b}`$, $`Wl\nu `$ and $`Zb\overline{b}`$ or $`Wbu_p`$ and $`Zl\overline{l}`$, if a lepton is produced in each of the b jets.
Similarly, the $`Z^0Z^0`$ production followed by the decays $`Zl\overline{l}`$ ($`l=e,\mu `$), $`Z\tau \overline{\tau }`$, where one of the $`\tau `$ decays into lepton while the other decays into jet, leads to three leptons in the final state. Within the same framework as above, the cross section is of order $`\sigma \left(p\overline{p}ZZ3l\right)2fb`$.
The $`Z^0Z^0`$ production can also contribute weakly to the 3 leptons background via the decays: $`Zl\overline{l}`$ and $`Zb\overline{b}`$ or $`Zb\overline{b}`$ and $`Zb\overline{b}`$, since a lepton can be produced in a b jet.
It has been pointed out recently that the $`WZ^{}`$ (throughout this paper a star indicates a virtual particle) and the $`W\gamma ^{}`$ productions could represent important contributions to the trilepton background . The complete list of contributions to the 3 leptons final state from the $`WZ`$,$`W\gamma ^{}`$ and $`ZZ`$ productions, including cases where either one or both of the gauge bosons can be virtual, has been calculated in . The authors of have found that the $`WZ`$, $`W\gamma ^{}`$ and $`ZZ`$ backgrounds (including virtual boson(s)) at the upgraded Tevatron have together a cross section of order $`0.5fb`$ after the following cuts have been implemented: $`P_t\left(l_1\right)>20GeV`$, $`P_t\left(l_2\right)>15GeV`$, $`P_t\left(l_3\right)>10GeV`$; $`\left|\eta (l_1,l_{2,3})\right|<1.0,2.0`$; $`ISO_{\delta R=0.4}<2GeV`$; $`E/_T>25GeV`$; $`81GeV<M_{inv}\left(l\overline{l}\right)<101GeV`$; $`12GeV<M_{inv}\left(l\overline{l}\right)`$; $`60GeV<m_T(l,E/_T)<85GeV`$.
We note that there is at most one hard jet in the 3 leptons backgrounds generated by the $`WZ`$, $`W\gamma `$ and $`ZZ`$ productions (including virtual boson(s)). Since the number of hard jets is equal to 2 in our signal (see Section 4.1), a jet veto can thus reduce this Standard Model background with respect to the signal.
Other small sources of Standard Model background have been estimated in : The productions like $`Zb`$, $`Wt`$ or $`Wt\overline{t}`$. After applying cuts on the geometrical acceptance, the transverse momentum and the isolation, these backgrounds are expected to be at most of order $`10^4pb`$ in $`p\overline{p}`$ collisions with a center of mass energy of $`\sqrt{s}=2TeV`$. We have checked that the $`Zb`$ production gives a negligible contribution to the 3 lepton signature.
There are finally some non-physics sources of background. First, the 4 leptons signal, which can be generated by the $`Z^0Z^0`$ and $`t\overline{t}`$ productions, appears as a 3 leptons signature if one of the leptons is missed. Besides, the processes $`p\overline{p}Z+X,DrellYan+X`$ would mimic a trilepton signal if $`X`$ fakes a lepton. Monte Carlo simulations using simplified detector simulation, like for example SHW as in the present study (see Section 4.4), cannot give a reliable estimate of this background. A knowledge of the details of the detector response as well as the jet fragmentation is necessary in order to determinate the probability to fake a lepton. In , using standard cuts the background coming from $`p\overline{p}Z+X,DrellYan+X`$ has been estimated to be of order $`2fb`$ at Tevatron with $`\sqrt{s}=2TeV`$. The authors of have also estimated the background from the three-jet events faking trilepton signals to be around $`10^3fb`$.
Hence for the study of the Standard Model background associated to the 3 lepton signature at Tevatron Run II, we consider the $`W^\pm Z^0`$ production and both the physics and non-physics contributions generated by the $`Z^0Z^0`$ and $`t\overline{t}`$ productions.
### 4.3 Supersymmetric background of the 3 lepton signature at Tevatron
If an excess of events is observed in the three lepton channel at Tevatron, one would wonder what is the origin of those anomalous events. One would thus have to consider all of the supersymmetric productions leading to the three lepton signature. In the present context of R-parity violation, multileptonic final states can be generated by the single chargino production involving $`\mathit{}_p`$ couplings, but also by the supersymmetric particle pair production which involves only gauge couplings . In $`\mathit{}_p`$ models, the superpartner pair production can even lead to the trilepton signature . As a matter of fact, both of the produced supersymmetric particles decay, either directly or through cascade decays, into the LSP which is the neutralino in our framework. In the hypothesis of a dominant $`\lambda ^{}`$ coupling constant, each of the 2 produced neutralinos can decay into a charged lepton and two quarks: at least two charged leptons and four jets in the final state are produced. The third charged lepton can be generated in the cascade decays as for example at the level of the chargino decay $`\stackrel{~}{\chi }^\pm \stackrel{~}{\chi }^0l^\pm \nu `$.
In Table 1, we show for different mSUGRA points the cross section of the sum of all superpartner pair productions, namely the $`R_p`$ conserving SUSY background of the 3 lepton signature generated by the single chargino production. As can be seen in this table, the summed superpartner pair production rate decreases as $`m_0`$ and $`m_{1/2}`$ increase. This is due to the increase of the superpartner masses as the $`m_0`$ or $`m_{1/2}`$ parameter increases. The SUSY background will be important only for low values of $`m_0`$ and $`m_{1/2}`$ as we will see in the following.
### 4.4 Cuts
In order to simulate the single chargino production $`p\overline{p}\stackrel{~}{\chi }_1^\pm l^{}`$ at Tevatron, the matrix elements (see Appendix A) of this process have been implemented in a version of the SUSYGEN event generator allowing the generation of $`p\overline{p}`$ reactions . The Standard Model background ($`W^\pm Z^0`$, $`Z^0Z^0`$ and $`t\overline{t}`$ productions) has been simulated using the PYTHIA event generator and the SUSY background (all SUSY particles pair productions) using the HERWIG event generator . SUSYGEN, PYTHIA and HERWIG have been interfaced with the SHW detector simulation package , which mimics an average of the CDF and D0 Run II detector performance.
We have developped a series of cuts in order to enhance the signal-to-background ratio.
First, we have selected the events with at least three leptons where the leptons are either an electron, a muon or a tau reconstructed from a jet, namely $`N_l3\left[l=e,\mu ,\tau \right]`$. We have also considered the case where the selected leptons are only electrons and muons, namely $`N_l3\left[l=e,\mu \right]`$.
The selection criteria on the jets was to have a number of jets greater or equal to two, where the jets have a transverse momentum higher than $`10GeV`$, namely $`N_j2`$ with $`P_t\left(j\right)>10GeV`$. This jet veto reduces the 3 lepton backgrounds coming from the $`W^\pm Z^0`$ and $`Z^0Z^0`$ productions. Indeed, the $`W^\pm Z^0`$ production generates no hard jets and the $`Z^0Z^0`$ production generates at most one hard jet. Moreover, the hard jet produced in the $`Z^0Z^0`$ background is generated by a tau decay (see Section 4.2) and can thus be identified as a tau.
Besides, some effective cuts concerning the energies of the produced leptons have been applied. In Fig.7, we show the distributions of the third leading lepton energy in the 3 lepton events produced by the Standard Model background ($`W^\pm Z^0`$, $`Z^0Z^0`$ and $`t\overline{t}`$) and the SUSY signal. Based on those kinds of distributions, we have chosen the following cut on the third leading lepton energy: $`E\left(l_3\right)>10GeV`$. Similarly, we have required that the energies of the 2 leading leptons verify $`E\left(l_2\right)>20GeV`$ and $`E\left(l_1\right)>20GeV`$.
We will refer to all the selection criteria described above, namely $`N_l3\left[l=e,\mu ,\tau \right]`$ with $`E\left(l_1\right)>20GeV`$, $`E\left(l_2\right)>20GeV`$, $`E\left(l_3\right)>10GeV`$, and $`N_j2`$ with $`P_t\left(j\right)>10GeV`$, as cut $`1`$.
Finally, since the leptons originating from the hadron decays (as in the $`t\overline{t}`$ production) are not well isolated, we have applied some cuts on the lepton isolation. We have imposed the isolation cut $`\mathrm{\Delta }R=\sqrt{\delta \varphi ^2+\delta \theta ^2}>0.4`$ where $`\varphi `$ is the azimuthal angle and $`\theta `$ the polar angle between the 3 most energetic charged leptons and the 2 hardest jets. Such a cut is for instance motivated by the distributions shown in Fig.8 of the $`\mathrm{\Delta }R`$ angular difference between the third leading lepton and the second leading jet, in the 3 lepton events generated by the SUSY signal and Standard Model background. We call cut $`\mathrm{\Delta }R>0.4`$ together with cut $`1`$, cut $`2`$.
In order to eliminate poorly isolated leptons, we have also required that $`E<2GeV`$, where $`E`$ represents the summed energies of the jets being close to a muon or an electron, namely the jets contained in the cone centered on a muon or an electron and defined by $`\mathrm{\Delta }R<0.25`$. This cut is not applied for taus candidates as they have hadronic decays. It is quite efficient (see Fig.21 for the 2 lepton case) since we sum over all jet energies in the cone. The Standard Model background shows more jets and less separation between jets and leptons in $`(\theta ,\varphi )`$ in final state than the single productions <sup>1</sup><sup>1</sup>1This cut will have to be fine tuned with real events since it will depend on the energy flow inside the detector, the overlap and minimum biased events.. We denote cut $`E<2GeV`$ plus cut $`2`$ as cut $`3`$ <sup>2</sup><sup>2</sup>2Although it has not been applied, we mention another kind of isolation cut which allows to further reduce the Standard Model background: $`\delta \varphi >70^{}`$ between the leading charged lepton and the 2 hardest jets..
The selected events require high energy charged leptons and jets and can thus easily be triggered at Tevatron. This is illustrated in Fig.9 where we show the energy distributions of the 3 leptons and the second leading jet in the 3 leptons events selected by applying cut $`3`$ and generated by the SUSY signal and Standard Model background.
In Table 2, we give the numbers of three lepton events expected from the Standard Model background at Tevatron Run II with the various cuts described above. We see in Table 2 that the main source of Standard Model background to the three lepton signature at Tevatron is the $`t\overline{t}`$ production. This is due to the important cross section of the $`t\overline{t}`$ production compared to the other Standard Model backgrounds (see Section 4.2). Table 2 also shows that the $`t\overline{t}`$ background is relatively more suppressed than the other sources of Standard Model background by the lepton isolation cuts. The reason is that in the $`t\overline{t}`$ background, one of the 3 charged leptons of the final state is generated in a $`b`$-jet and is thus not well isolated.
In Table 3, we give the number of three lepton events generated by the SUSY background (all superpartners pair productions) at Tevatron Run II as a function of the $`m_0`$ and $`m_{1/2}`$ parameters for the cut 3. This number of events decreases as $`m_0`$ and $`m_{1/2}`$ increase due to the behaviour of the summed superpartners pair productions cross section in the SUSY parameter space (see Section 4.3).
### 4.5 Results
#### 4.5.1 Discovery potential for the $`\lambda _{2jk}^{}`$ coupling constant
We first present the reach in the mSUGRA parameter space obtained from the analysis of the trilepton signature at Tevatron Run II generated by the single chargino production through the $`\lambda _{211}^{}`$ coupling, namely $`p\overline{p}\stackrel{~}{\chi }_1^\pm \mu ^{}`$. The sensitivity that can be obtained on the $`\lambda _{2jk}^{}`$ ($`j`$ and $`k`$ being not equal to $`1`$ simultaneously) couplings based on the $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ production analysis will be discussed at the end of this section for a given mSUGRA point. We give more detailed results for the case of a single dominant $`\lambda _{211}^{}`$ coupling since this $`\mathit{}_p`$ coupling gives the highest partonic luminosity to the $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ production (see Section 3.1.1) and leads thus to the highest sensitivities.
In Fig.10, we present the $`3\sigma `$ and $`5\sigma `$ discovery contours and the limits at $`95\%`$ confidence level in the plane $`m_0`$ versus $`m_{1/2}`$, for $`sign\left(\mu \right)<0`$, $`\mathrm{tan}\beta =1.5`$ and using a set of values for $`\lambda _{211}^{}`$ and the luminosity. This discovery potential was obtained by considering the $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ production and the background originating from the Standard Model. The signal and background were selected by using cut $`3`$ described in Section 4.4. The results presented for a luminosity of $`=0.5fb^1`$ in Fig.10 and Fig.11 were obtained with cut 2 only in order to optimize the sensitivity on the SUSY parameters. The reduction of the sensitivity on $`\lambda _{211}^{}`$ observed in Fig.10 when either $`m_0`$ or $`m_{1/2}`$ increases is due to the decrease of the $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ production cross section with $`m_0`$ or $`m_{1/2}`$ (or equivalently $`M_2`$), which can be observed in Fig.4. In Fig.10, we also see that for all the considered values of $`\lambda _{211}^{}`$ and the luminosity, the sensitivity on $`m_{1/2}`$ is reduced to low masses in the domain $`m_0\stackrel{<}{}200GeV`$. This important reduction of the sensitivity as $`m_0`$ decreases is due to the decrease of the phase space factor associated to the decay $`\stackrel{~}{\nu }_\mu \stackrel{~}{\chi }^\pm \mu ^{}`$ (see Section 3.1.1). Finally, we note from Fig.3 that for $`sign\left(\mu \right)>0`$ the $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ production cross section, and thus the sensitivities presented in Fig.10, would incur a little increase compared to the case $`sign\left(\mu \right)<0`$.
In Fig.11, the discovery potential is shown in the $`\lambda _{211}^{}`$-$`m_0`$ plane for different values of $`M_2`$ and the luminosity. For a given value of $`M_2`$, we note that the sensitivity on the $`\lambda _{211}^{}`$ coupling decreases at high and low values of $`m_0`$. The main explanation is the decrease of the $`p\overline{p}\stackrel{~}{\chi }_1^\pm \mu ^{}`$ rate at high and low values of $`m_0`$ which appears clearly in Fig.5. We also observe, as in Fig.10, a decrease of the sensitivity on the $`\lambda _{211}^{}`$ coupling when $`M_2`$ (or equivalently $`m_{1/2}`$) increases for a fixed value of $`m_0`$.
The strongest bounds on the supersymmetric masses obtained at LEP in an $`\mathit{}_p`$ model with a non-vanishing $`\lambda ^{}`$ Yukawa coupling are $`m_{\stackrel{~}{\chi }_1^0}>26GeV`$ (for $`m_0=200GeV`$ and $`\mathrm{tan}\beta =\sqrt{2}`$ ), $`m_{\stackrel{~}{\chi }_1^\pm }>100GeV`$, $`m_{\stackrel{~}{l}}>93GeV`$, $`m_{\stackrel{~}{\nu }}>86GeV`$ . For the minimum values of $`m_0`$ and $`m_{1/2}`$ spanned by the parameter space described in Figures 10 and 11, namely $`m_0=100GeV`$ and $`M_2=100GeV`$, the mass spectrum is $`m_{\stackrel{~}{\chi }_1^\pm }=113GeV`$, $`m_{\stackrel{~}{\chi }_1^0}=54GeV`$, $`m_{\stackrel{~}{\nu }_L}=127GeV`$, $`m_{\stackrel{~}{l}_L}=137GeV`$, $`m_{\stackrel{~}{l}_R}=114GeV`$, so that we are well above these limits. Since both the scalar and gaugino masses increase with $`m_0`$ and $`m_{1/2}`$, the parameter space described in Figures 10 and 11 lies outside the SUSY parameters ranges excluded by LEP data . Therefore, the discovery potential of Figures 10 and 11 represents an important improvement with respect to the supersymmetric masses limits derived from LEP data . Figures 10 and 11 show also that the low-energy bound on the considered $`\mathit{}_p`$ coupling, $`\lambda _{211}^{}<0.09\left(m_{\stackrel{~}{d}_R}/100GeV\right)`$ at $`1\sigma `$ (from $`\pi `$ decay) , can be greatly improved.
Interesting sensitivities on the SUSY parameters can already be obtained within the first year of Run II at Tevatron with a low luminosity ($`=0.5fb^1`$) and no reconstruction of the tau-jets. To illustrate this point, we present in Fig.12 and Fig.13 the same discovery potentials as in Fig.10 and Fig.11, respectively, obtained without reconstruction of the tau leptons decaying into jets. By comparing Fig.10, Fig.11 and Fig.12, Fig.13, we observe that the sensitivity on the SUSY parameters is weakly affected by the reconstruction of the tau-jets <sup>3</sup><sup>3</sup>3This is actually an artefact of the method: cut 3 is our most efficient cut to reduce the Standard Model background with electrons and muons but is not applied with taus. Thus, the relative ratio signal over background is not so good with taus. Finding another efficient cut could improve our discovery potential and limits using taus..
Using the ratios of the cross sections for the $`\stackrel{~}{\chi }_1^+\mu ^{}`$ productions via different $`\lambda _{2jk}^{}`$ couplings, one can deduce from the sensitivity obtained on $`\lambda _{211}^{}`$ via the 3 lepton final state analysis an estimation of the sensitivity on any $`\lambda _{2jk}^{}`$ coupling. For instance, let us consider the SUSY point $`m_0=180GeV`$, $`M_2=200GeV`$, $`\mathrm{tan}\beta =1.5`$ and $`\mu =200GeV`$ ($`m_{\stackrel{~}{u}_L}=601GeV`$, $`m_{\stackrel{~}{d}_L}=603GeV`$, $`m_{\stackrel{~}{u}_R}=582GeV`$, $`m_{\stackrel{~}{d}_R}=580GeV`$, $`m_{\stackrel{~}{l}_L}=253GeV`$, $`m_{\stackrel{~}{l}_R}=205GeV`$ $`m_{\stackrel{~}{\nu }_L}=248GeV`$, $`m_{\stackrel{~}{\chi }_1^\pm }=199GeV`$, $`m_{\stackrel{~}{\chi }_1^0}=105GeV`$) which corresponds, as can be seen in Fig.11, to the point where the sensitivity on $`\lambda _{211}^{}`$ is maximized for $`M_2=200GeV`$. We can see on Fig.5 that for this SUSY point, the ratio between the rates of the $`\stackrel{~}{\chi }_1^+\mu ^{}`$ productions via $`\lambda _{211}^{}`$ and $`\lambda _{221}^{}`$ is $`\sigma \left(\lambda _{211}^{}\right)/\sigma \left(\lambda _{221}^{}\right)7.9`$ for same values of the $`\mathit{}_p`$ couplings. Therefore, since the single chargino production rate scales as $`\lambda ^2`$ (see Appendix A), the sensitivity on $`\lambda _{221}^{}`$ at this SUSY point is equal to the sensitivity obtained on $`\lambda _{211}^{}`$ ($`0.02`$ at $`95\%CL`$ with $`=2fb^1`$ as shows Fig.11) multiplied by the factor $`\sqrt{7.9}`$, namely $`0.05`$. This result represents a significant improvement with respect to the low-energy indirect limit $`\lambda _{221}^{}<0.18\left(m_{\stackrel{~}{d}_R}/100GeV\right)`$ . Using the same method, we find at the same SUSY point the sensitivities on the $`\lambda _{2jk}^{}`$ coupling constants given in Table 4.
All the sensitivities on the $`\lambda _{2jk}^{}`$ coupling constants given in Table 4 are stronger than the low-energy bounds of which we rewrite here: $`\lambda _{21k}^{}<0.09\left(m_{\stackrel{~}{d}_{kR}}/100GeV\right)`$ at $`1\sigma `$ ($`\pi `$ decay), $`\lambda _{22k}^{}<0.18\left(m_{\stackrel{~}{d}_{kR}}/100GeV\right)`$ at $`1\sigma `$ ($`D`$ decay), $`\lambda _{231}^{}<0.22\left(m_{\stackrel{~}{b}_L}/100GeV\right)`$ at $`2\sigma `$ ($`\nu _\mu `$ deep inelastic scattering), $`\lambda _{232}^{}<0.36\left(m_{\stackrel{~}{q}}/100GeV\right)`$ at $`1\sigma `$ ($`R_\mu `$), $`\lambda _{233}^{}<0.36\left(m_{\stackrel{~}{q}}/100GeV\right)`$ at $`1\sigma `$ ($`R_\mu `$).
In the case of a single dominant $`\lambda _{2j3}^{}`$ coupling, the neutralino decays as $`\stackrel{~}{\chi }_1^0\mu u_jb`$ and the semileptonic decay of the b-quark could affect the analysis efficiency. Therefore in this case, the precise sensitivity cannot be simply calculated by scaling the value obtained for $`\lambda _{211}^{}`$. Nevertheless, the order of magnitude of the sensitivity which can be inferred from our analysis should be correct.
The range of SUSY parameters in which the constraint on a given $`\lambda _{2jk}^{}`$ coupling constant obtained via the three leptons analysis is stronger than the relevant low-energy bound depends on the low-energy bound itself as well as on the values of the cross section for the single chargino production via the considered $`\lambda _{2jk}^{}`$ coupling.
Finally, we remark that while the low-energy constraints on the $`\lambda _{2jk}^{}`$ couplings become weaker as the squark masses increase, the sensitivities on those couplings obtained from the three leptons analysis are essentially independent of the squark masses as long as $`m_{\stackrel{~}{q}}>m_{\stackrel{~}{\chi }_1^\pm }`$ (recall that the branching ratio of the decay $`\stackrel{~}{\chi }_1^\pm q\overline{q}\stackrel{~}{\chi }_1^0`$ is greatly enhanced when $`m_{\stackrel{~}{q}}<m_{\stackrel{~}{\chi }_1^\pm }`$).
We end this section by some comments on the effect of the supersymmetric $`R_p`$ conserving background to the 3 lepton signature. In order to illustrate this discussion, we consider the results on the $`\lambda _{211}^{}`$ coupling constant.
We see from Table 3 that the SUSY background to the 3 lepton final state can affect the sensitivity on the $`\lambda _{211}^{}`$ coupling constant obtained by considering only the Standard Model background, which is shown in Fig.10, only in the region of small superpartner masses, namely in the domain $`m_{1/2}\stackrel{<}{}300GeV`$ for $`\mathrm{tan}\beta =1.5`$, $`sign\left(\mu \right)<0`$ and assuming a luminosity of $`=1fb^1`$.
In contrast with the SUSY signal amplitude which is increased if $`\lambda _{211}^{}`$ is enhanced, the SUSY background amplitude is typically independent on the value of the $`\lambda _{211}^{}`$ coupling constant since the superpartner pair production does not involve $`\mathit{}_p`$ couplings. Therefore, even if we consider the SUSY background in addition to the Standard Model one, it is still true that large values of the $`\lambda _{211}^{}`$ coupling can be probed over a wider domain of the SUSY parameter space than low values, as can be observed in Fig.10 for $`m_{1/2}\stackrel{>}{}300GeV`$. Note that in Fig.10 larger values of $`\lambda _{211}^{}`$ could have been considered as the low-energy bound on this $`\mathit{}_p`$ coupling, namely $`\lambda _{211}^{}<0.09\left(m_{\stackrel{~}{d}_R}/100GeV\right)`$ , is proportional to the squark mass.
Finally, we mention that further cuts, as for instance some cuts based on the superpartner mass reconstructions (see Section 4.5.4), could allow to reduce the SUSY background to the 3 lepton signature.
#### 4.5.2 High $`\mathrm{tan}\beta `$ scenario
In mSUGRA, for large values of $`\mathrm{tan}\beta `$ and small values of $`m_0`$ compared to $`m_{1/2}`$, due to the large mixing in the third generation sfermions, the mass of the lighter $`\stackrel{~}{\tau }_1`$ slepton can become smaller than $`m_{\stackrel{~}{\chi }_1^\pm }`$, with the sneutrino remaining heavier than the $`\stackrel{~}{\chi }_1^\pm `$ so that the $`\stackrel{~}{\chi }_1^\pm l^{}`$ production rate can still be significant. In this situation, the efficiency for the 3 lepton signature arising mainly through, $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\tau }_1^\pm \nu _\tau ,\stackrel{~}{\tau }_1^\pm \stackrel{~}{\chi }_1^0\tau ^\pm ,\stackrel{~}{\chi }_1^0l_i^\pm u_jd_k`$, can be enhanced compared to the case where the 3 lepton signal comes from, $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0l^\pm \nu ,\stackrel{~}{\chi }_1^0l_i^\pm u_jd_k`$. Indeed, the branching ratio $`B\left(\stackrel{~}{\chi }_1^\pm \stackrel{~}{\tau }_1^\pm \nu _\tau \right)`$ can reach $`100\%`$, $`B\left(\stackrel{~}{\tau }_1^\pm \stackrel{~}{\chi }_1^0\tau ^\pm \right)100\%`$, since the $`\stackrel{~}{\chi }_1^0`$ is the LSP, $`B\left(\tau l\nu _l\nu _\tau \right)=35\%`$ ($`l=e,\mu `$) and the $`\tau `$-jets can be reconstructed at Tevatron Run II. However, in such a scenario the increased number of tau leptons in the final state leads to a softer charged lepton spectrum which tends to reduce the efficiency after cuts. Therefore, for relatively small values of $`m_0`$ compared to $`m_{1/2}`$, the sensitivity obtained in the high $`\mathrm{tan}\beta `$ scenario is essentially unaffected with respect to the low $`\mathrm{tan}\beta `$ situation, unless $`m_0`$ is small enough to render $`m_{\stackrel{~}{\tau }_1}`$ and $`m_{\stackrel{~}{\chi }_1^0}`$ almost degenerate. As a matter of fact, in such a situation, the energy of the tau produced in the decay $`\stackrel{~}{\tau }_1^\pm \stackrel{~}{\chi }_1^0\tau ^\pm `$ often falls below the analysis cuts. Therefore, this degeneracy results in a loss of signal efficiency after cuts, at small values of $`m_0`$ compared to $`m_{1/2}`$, and thus in a loss of sensitivity, with respect to the low $`\mathrm{tan}\beta `$ situation. This can be seen by comparing Fig.10, Fig.11 and Fig.14, Fig.15. Indeed, the decrease of the sensitivity on $`m_{1/2}`$ at low $`m_0`$ is stronger for high $`\mathrm{tan}\beta `$ (see Fig.14) than for low $`\mathrm{tan}\beta `$ (see Fig.10). Similarly, the decrease of the sensitivity on $`\lambda _{211}^{}`$ at low $`m_0`$ is stronger for high $`\mathrm{tan}\beta `$ (see Fig.15) than for low $`\mathrm{tan}\beta `$ (see Fig.11).
The effect on the discovery potential of the single chargino production rate increase at large $`\mathrm{tan}\beta `$ values shown in Fig.2 is hidden by the large $`\mathrm{tan}\beta `$ scenario influences on the cascade decays described above.
In contrast with the low $`\mathrm{tan}\beta `$ scenario (see Section 4.5.1), the sensitivity on the SUSY parameters depends in a significant way on the reconstruction of the tau-jets in case where $`\mathrm{tan}\beta `$ is large, as can be seen in Fig.14 and Fig.15. The reason is the increased number of tau leptons among the final state particles in a large $`\mathrm{tan}\beta `$ model. This is due to the decrease of the lighter stau mass which tends to increase the $`B\left(\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0\tau ^\pm \nu _\tau \right)`$ branching ratio.
#### 4.5.3 Discovery potential for the $`\lambda _{1jk}^{}`$ and $`\lambda _{3jk}^{}`$ coupling constants
In Fig.16, we present the $`3\sigma `$ and $`5\sigma `$ discovery contours and the limits at $`95\%`$ confidence level in the plane $`m_0`$ versus $`m_{1/2}`$, for $`sign\left(\mu \right)<0`$, $`\mathrm{tan}\beta =1.5`$, $`\lambda _{311}^{}=0.10`$ and various values of the luminosity. In Fig.17, the discovery potential is shown in the $`\lambda _{311}^{}`$-$`m_0`$ plane for $`M_2=200GeV`$. Comparing Fig.16, Fig.17 and Fig.10, Fig.11, we see that the sensitivity on the SUSY parameters is weaker in the case of a single dominant $`\lambda _{311}^{}`$ coupling than in the case of a single dominant $`\lambda _{211}^{}`$ coupling. The reason is that in the case of a single dominant $`\lambda _{3jk}^{}`$ coupling constant, tau leptons are systematically produced at the chargino production level $`p\overline{p}\stackrel{~}{\chi }_1^\pm \tau ^{}`$ (see Fig.1(a)) as well as in the LSP decay $`\stackrel{~}{\chi }_1^0\tau u_jd_k`$ (see Section 4.1), so that the number of tau leptons among the 3 charged leptons of the final state is increased compared to the dominant $`\lambda _{2jk}^{}`$ case. The decrease in sensitivity is due to the fact that a lepton (electron or muon) generated in a tau decay has an higher probability not to pass the analysis requirements concerning the particle energy and that the reconstruction efficiency for a tau decaying into a jet is limited.
Nevertheless, the discovery potentials of Fig.16 and Fig.17 represent also an important improvement with respect to the experimental mass limits from LEP measurements and to the low-energy indirect constraint $`\lambda _{311}^{}<0.10\left(m_{\stackrel{~}{d}_R}/100GeV\right)`$ at $`1\sigma `$ (from $`\tau ^{}\pi ^{}\nu _\tau `$) .
We also observe in Fig.16 and Fig.17 that the results obtained from the $`\stackrel{~}{\chi }_1^\pm \tau ^{}`$ production analysis in the case of a single dominant $`\lambda _{3jk}^{}`$ coupling depend strongly on the reconstruction of the tau-jets. This is due to the large number of tau leptons among the 3 charged leptons of the considered final state.
Using the same method and same SUSY point as in Section 4.5.1, we have estimated the sensitivity on all the $`\lambda _{3jk}^{}`$ coupling constants from the sensitivity obtained on $`\lambda _{311}^{}`$ at $`95\%CL`$ for a luminosity of $`=2fb^1`$. The results are given in Table 5.
All the sensitivities on the $`\mathit{}_p`$ couplings presented in Table 5, except those on $`\lambda _{32k}^{}`$, are stronger than the present indirect limits on the same $`\mathit{}_p`$ couplings: $`\lambda _{31k}^{}<0.10\left(m_{\stackrel{~}{d}_{kR}}/100GeV\right)`$ at $`1\sigma `$ ($`\tau ^{}\pi ^{}\nu _\tau `$), $`\lambda _{32k}^{}<0.20`$ (for $`m_{\stackrel{~}{l}}=m_{\stackrel{~}{q}}=100GeV`$) at $`1\sigma `$ ($`D^0\overline{D}^0`$ mix), $`\lambda _{33k}^{}<0.48\left(m_{\stackrel{~}{q}}/100GeV\right)`$ at $`1\sigma `$ ($`R_\tau `$) .
We mention that in the case of a single dominant $`\lambda _{3j3}^{}`$ coupling, the neutralino decays as $`\stackrel{~}{\chi }_1^0\tau u_jb`$ so that the b semileptonic decay could affect a little the analysis efficiency.
We discuss now the sensitivities that could be obtained on a single dominant $`\lambda _{1jk}^{}`$ coupling constant via the analysis of the reaction $`p\overline{p}\stackrel{~}{\chi }_1^\pm e^{}`$ (see Fig.1(a)). Since the cross section of the $`\stackrel{~}{\chi }_1^\pm e^{}`$ production through $`\lambda _{1jk}^{}`$ is equal to the rate of the $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ production via $`\lambda _{2jk}^{}`$, for same $`j`$ and $`k`$ indices (see Section 3.1.1), the sensitivity obtained on a $`\lambda _{1jk}^{}`$ coupling constant is expected to be identical to the sensitivity on $`\lambda _{2jk}^{}`$. If we assume that the sensitivities obtained on the $`\lambda _{1jk}^{}`$ couplings are equal to those presented in Table 4, we remark that for the SUSY point chosen in this table only the sensitivities expected for the $`\lambda _{112}^{}`$, $`\lambda _{113}^{}`$, $`\lambda _{121}^{}`$, $`\lambda _{131}^{}`$ and $`\lambda _{132}^{}`$ couplings are stronger than the corresponding low-energy bounds: $`\lambda _{11k}^{}<0.02\left(m_{\stackrel{~}{d}_{kR}}/100GeV\right)`$ at $`2\sigma `$ (Charged current universality), $`\lambda _{1j1}^{}<0.035\left(m_{\stackrel{~}{q}_{jL}}/100GeV\right)`$ at $`2\sigma `$ (Atomic parity violation), $`\lambda _{132}^{}<0.34`$ at $`1\sigma `$ for $`m_{\stackrel{~}{q}}=100GeV`$ ($`R_e`$) . The reason is that the low-energy constraints on the $`\lambda _{1jk}^{}`$ couplings are typically more stringent than the limits on the $`\lambda _{2jk}^{}`$ couplings .
#### 4.5.4 Mass reconstructions
The $`\stackrel{~}{\chi }_1^0`$ neutralino decays in our framework as $`\stackrel{~}{\chi }_1^0l_iu_jd_k`$ through the $`\lambda _{ijk}^{}`$ coupling constant. The invariant mass distribution of the lepton and the 2 jets coming from this decay channel is peaked at the $`\stackrel{~}{\chi }_1^0`$ mass. The experimental analysis of this invariant mass distribution would thus be particularly interesting since it would allow a model independent determination of the lightest neutralino mass.
We have performed the $`\stackrel{~}{\chi }_1^0`$ mass reconstruction based on the 3 lepton signature analysis. The difficulty of this mass reconstruction lies in the selection of the lepton and the 2 jets coming from the $`\stackrel{~}{\chi }_1^0`$ decay. In the signal we are considering, the only jets come from the $`\stackrel{~}{\chi }_1^0`$ decay, and of course from the initial and final QCD radiations. Therefore, if there are more than 2 jets in the final state we have selected the 2 hardest ones. It is more subtle for the selection of the lepton since our signal contains 3 leptons. We have considered the case of a single dominant coupling of type $`\lambda _{2jk}^{}`$ and focused on the $`e\mu \mu `$ final state. In these events, one of the $`\mu ^\pm `$ is generated in the decay of the produced sneutrino as $`\stackrel{~}{\nu }_\mu \stackrel{~}{\chi }_1^\pm \mu ^{}`$ and the other one in the decay of the $`\stackrel{~}{\chi }_1^0`$ as $`\stackrel{~}{\chi }_1^0\mu ^\pm u_jd_k`$, while the electron comes from the chargino decay $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0e^\pm \nu _e`$. Indeed, the dominant contribution to the single chargino production is the resonant sneutrino production (see Fig.1). In order to select the muon from the $`\stackrel{~}{\chi }_1^0`$ decay we have chosen the softer muon, since for relatively important values of the $`m_{\stackrel{~}{\nu }_\mu }m_{\stackrel{~}{\chi }_1^\pm }`$ mass difference the muon generated in the sneutrino decay is the most energetic. Notice that for too degenerate $`\stackrel{~}{\nu }_\mu `$ and $`\stackrel{~}{\chi }_1^\pm `$ masses, the sensitivity on the SUSY parameters suffers a strong decrease as shown in Section 4.5.1.
We present in Fig.18 the invariant mass distribution of the muon and the 2 jets produced in the $`\stackrel{~}{\chi }_1^0`$ decay. This distribution has been obtained by using the selection criteria described above and by considering the mSUGRA point: $`m_0=200GeV`$, $`M_2=150GeV`$, $`\mathrm{tan}\beta =1.5`$, $`sign\left(\mu \right)<0`$ and $`\lambda _{211}^{}=0.09`$ ($`m_{\stackrel{~}{\chi }_1^0}=77.7GeV`$, $`m_{\stackrel{~}{\chi }_1^\pm }=158.3GeV`$, $`m_{\stackrel{~}{\nu }_L}=236GeV`$). We also show on the plot of Fig.18 the fit of the invariant mass distribution. As can be seen from this fit, the distribution is well peaked around the $`\stackrel{~}{\chi }_1^0`$ generated mass. The average reconstructed $`\stackrel{~}{\chi }_1^0`$ mass is of $`71\pm 9GeV`$.
We have also performed the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_\mu `$ mass reconstructions based on the 3 lepton signature analysis in the scenario of a single dominant coupling of type $`\lambda _{2jk}^{}`$. The $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_\mu `$ masses reconstructions are based on the 4-momentum of the neutrino present in the $`3l+2j+\nu `$ final state (see Section 4.1). The transverse component of this momentum can be deduced from the momentum of the charged leptons and jets present in the final state. However, the longitudinal component of the neutrino momentum remains unknown due to the poor detection at small polar angle values. Therefore, in this study we have assumed a vanishing longitudinal component of the neutrino momentum. Besides, we have focused on the $`e\mu \mu `$ events as in the $`\stackrel{~}{\chi }_1^0`$ mass reconstruction study. In this context, the cascade decay initiated by the produced lightest chargino is $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0e^\pm \nu _e`$, $`\stackrel{~}{\chi }_1^0\mu ^\pm u_jd_k`$. Therefore, the $`\stackrel{~}{\chi }_1^\pm `$ has been reconstructed from the softer muon, the 2 jets, the electron and the neutrino present in the final state, since the softer muon is mainly generated in the $`\stackrel{~}{\chi }_1^0`$ decay as explained above. The $`\stackrel{~}{\nu }_\mu `$ has then been reconstructed from the $`\stackrel{~}{\chi }_1^\pm `$ and the leading muon of the final state. This was motivated by the fact that the dominant contribution to the single chargino production is the reaction $`p\overline{p}\stackrel{~}{\nu }_\mu \stackrel{~}{\chi }_1^\pm \mu ^{}`$ (see Fig.1).
In Fig.19, we present the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_\mu `$ mass reconstructions performed through the method presented above. We also show on the plots of Fig.19 the fits of the invariant mass distributions. As can be seen from those fits, the distributions are well peaked around the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_{\mu L}`$ generated masses. The average reconstructed masses are $`m_{\stackrel{~}{\chi }_1^\pm }=171\pm 35GeV`$ and $`m_{\stackrel{~}{\nu }_{\mu L}}=246\pm 32GeV`$. This study on the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_{\mu L}`$ masses shows that based on a simplified mass reconstruction analysis promising results are obtained from the 3 lepton signature generated by the single $`\stackrel{~}{\chi }_1^\pm `$ production. The $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_{\mu L}`$ mass reconstructions can be improved using constrained fits.
In the hypothesis of a single dominant coupling constant of type $`\lambda _{1jk}^{}`$, exactly the same kind of $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_\mu `$ mass reconstructions can be performed by selecting the $`e+e+\mu +2j+\nu `$ events. In contrast, the case of a single dominant $`\lambda _{3jk}^{}`$ coupling requires more sophisticated methods.
As a conclusion, in the case of a single dominant coupling constant of type $`\lambda _{1jk}^{}`$ or $`\lambda _{2jk}^{}`$, the $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_\mu `$ mass reconstructions based on the 3 lepton signature generated by the single $`\stackrel{~}{\chi }_1^\pm `$ production at Tevatron can easily give precise results, in contrast with the mass reconstructions performed in the superpartner pair production analysis at hadronic colliders which suffer a high combinatorial background .
#### 4.5.5 Model dependence of the results
In this Section, we discuss qualitatively the impact on our results of the choice of our theoretical model, namely mSUGRA with the infrared fixed point hypothesis for the top quark Yukawa coupling. We focus on the discovery potentials obtained in Sections 4.5.1, 4.5.2 and 4.5.3, since the choice of the theoretical framework does not influence the study of the neutralino mass reconstruction made in Section 4.5.4 which is model independent.
The main effect of the infrared fixed point approach is to fix the value of the $`\mathrm{tan}\beta `$ parameter, up to the ambiguity on the low or high solution. Therefore, the infrared fixed point hypothesis has no important effects on the results since the dependences of the single gaugino productions rates on $`\mathrm{tan}\beta `$ are smooth, in the high $`\mathrm{tan}\beta `$ scenario (see Section 3.1.1).
As we have mentioned in Section 2, in the mSUGRA scenario, the $`\left|\mu \right|`$ parameter is fixed. This point does not influence much our results since the single gaugino production cross sections vary weakly with $`\left|\mu \right|`$ as shown in Section 3.1.1.
Another particularity of the mSUGRA model is that the LSP is the $`\stackrel{~}{\chi }_1^0`$ in most of the parameter space. For instance, in a model where the LSP would be the lightest chargino or a squark, the contribution to the three lepton signature from the $`\stackrel{~}{\chi }_1^\pm l^{}`$ production would vanish.
Finally in mSUGRA, the squark masses are typically larger than the lightest chargino mass so that the decay $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0l^\pm \nu `$ has a branching ratio of at least $`30\%`$ (see Section 4.1). In a scenario where $`m_{\stackrel{~}{\chi }_1^\pm }>m_{\stackrel{~}{q}}`$, the two-body decay $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{q}q`$ would be dominant so that the contribution to the three lepton signature from the $`\stackrel{~}{\chi }_1^\pm l^{}`$ production would be small. Besides, in mSUGRA, the $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0`$ mass difference is typically large enough to avoid significant branching ratio for the $`\mathit{}_p`$ decay of the lightest chargino which would result in a decrease of the sensitivities on the SUSY parameters presented in Sections 4.5.1, 4.5.2 and 4.5.3.
In a model where the contribution to the three lepton signature from the $`\stackrel{~}{\chi }^\pm l^{}`$ production would be suppressed, the three lepton final state could be generated in a significant way by other single gaugino productions, namely the $`\stackrel{~}{\chi }^\pm \nu `$, $`\stackrel{~}{\chi }^0l^{}`$ or $`\stackrel{~}{\chi }^0\nu `$ productions.
## 5 Like sign dilepton signature analysis
### 5.1 Signal
Within the context of the mSUGRA model, three of the single gaugino productions via $`\lambda _{ijk}^{}`$ presented in Section 3.1 can generate a final state containing a pair of same sign leptons. As a matter of fact, the like sign dilepton signature can be produced through the reactions $`p\overline{p}\stackrel{~}{\chi }_1^0l_i^\pm `$; $`p\overline{p}\stackrel{~}{\chi }_2^0l_i^\pm `$, $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0+X`$ ($`Xl^\pm `$); $`p\overline{p}\stackrel{~}{\chi }_1^\pm l_i^{}`$, $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0q\overline{q}`$ and $`p\overline{p}\stackrel{~}{\chi }_1^\pm \nu _i`$, $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0l^\pm \nu `$, $`i`$ corresponding to the flavour index of the $`\lambda _{ijk}^{}`$ coupling. Indeed, since the $`\stackrel{~}{\chi }_1^0`$ is a Majorana particle, it decays via $`\lambda _{ijk}^{}`$ into a lepton, as $`\stackrel{~}{\chi }_1^0l_iu_j\overline{d}_k`$, and into an anti-lepton, as $`\stackrel{~}{\chi }_1^0\overline{l}_i\overline{u}_jd_k`$, with the same probability. The $`\stackrel{~}{\chi }_{3,4}^0l_i^\pm `$, $`\stackrel{~}{\chi }_2^\pm l_i^{}`$ and $`\stackrel{~}{\chi }_2^\pm \nu _i`$ productions do not bring significant contributions to the like sign dilepton signature due to their relatively small cross sections (see Section 3.1.1).
In mSUGRA, the most important contribution to the like sign dilepton signature originates from the $`\stackrel{~}{\chi }_1^0l_i^\pm `$ production since this reaction has a dominant cross section in most of the mSUGRA parameter space, as shown in Section 3.1.1. The other reason is that if $`\stackrel{~}{\chi }_1^0`$ is the LSP, the $`\stackrel{~}{\chi }_1^0l_i^\pm `$ production rate is not affected by branching ratios of any cascade decay since the $`\stackrel{~}{\chi }_1^0`$ only decays through $`\mathit{}_p`$ coupling.
### 5.2 Standard Model background of the like sign dilepton signature at Tevatron
The $`b\overline{b}`$ production can lead to the like sign dilepton signature if both of the b quarks decay semi-leptonically. The leading order cross section of the $`\overline{b}b`$ production at Tevatron for an energy of $`\sqrt{s}=2TeV`$ is $`\sigma \left(p\overline{p}b\overline{b}\right)\mathrm{4.654\; 10}^{10}fb`$. This rate has been calculated with PYTHIA using the CTEQ2L structure function.
The $`t\overline{t}`$ production, followed by the decays $`tW^+bl^+\nu b`$, $`\overline{t}W^{}\overline{b}\overline{q}q\overline{b}\overline{q}ql^+\nu \overline{c}`$, or $`tW^+b\overline{q}qb\overline{q}ql^{}\overline{\nu }c`$, $`\overline{t}W^{}\overline{b}l^{}\overline{\nu }\overline{b}`$, also generates a final state with two same sign leptons. The leading order cross section of the $`t\overline{t}`$ production at $`\sqrt{s}=2TeV`$, including the relevant branching ratios, is $`\sigma \left(p\overline{p}t\overline{t}\right)\times 2\times B\left(Wl_p\nu _p\right)\times B\left(Wq_p\overline{q}_p^{}\right)3181fb`$ ($`2800fb`$) for $`m_{top}=170GeV`$ ($`175GeV`$) with $`p,p^{}=1,2,3`$.
The third important source of Standard Model background is the $`t\overline{b}/\overline{t}b`$ production since the (anti-)$`b`$ quark can undergo a semi-leptonic decay as $`bl^{}\overline{\nu }c`$ ($`\overline{b}l^+\nu \overline{c}`$) and the (anti-) top quark can decay simultaneously as $`tbW^+bl^+\nu `$ ($`\overline{t}\overline{b}W^{}\overline{b}l^{}\overline{\nu }`$). The leading order cross section at $`\sqrt{s}=2TeV`$ including the branching fraction is $`\sigma \left(p\overline{p}tq,\overline{t}q\right)\times B\left(Wl_p\nu _p\right)802fb`$ ($`687fb`$) for $`m_{top}=170GeV`$ ($`175GeV`$) with $`p=1,2,3`$.
Other small sources of Standard Model background are the $`W^\pm W^{}`$ production, followed by the decays: $`Wl\nu `$ and $`Wbu_p`$ ($`p=1,2`$) or $`Wbu_p`$ and $`Wbu_p`$ ($`p=1,2`$), the $`W^\pm Z^0`$ production, followed by the decays: $`Wl\nu `$ and $`Zb\overline{b}`$ or $`Wq_p\overline{q}_p^{}`$ and $`Zb\overline{b}`$, and the $`Z^0Z^0`$ production, followed by the decays: $`Zl\overline{l}`$ and $`Zb\overline{b}`$ or $`Zq_p\overline{q}_p`$ and $`Zb\overline{b}`$.
Finally, the 3 lepton final states generated by the $`Z^0Z^0`$ and $`W^\pm Z^0`$ productions (see Section 4.2) can be mistaken for like sign dilepton events in case where one of the leptons is lost in the detection. Non-physics sources of background can also be caused by some fake leptons or by the misidentification of the charge of a lepton.
Therefore for the study of the Standard Model background associated to the like sign dilepton signal at Tevatron Run II, we consider the $`b\overline{b}`$, the $`t\overline{t}`$, the $`W^\pm W^{}`$ and the single top production and both the physics and non-physics contributions generated by the $`W^\pm Z^0`$ and $`Z^0Z^0`$ productions.
### 5.3 Supersymmetric background of the like sign dilepton signature at Tevatron
All the pair productions of superpartners are a source of SUSY background for the like sign dilepton signature originating from the single gaugino productions. Indeed, both of the produced superpartners initiate a cascade decay ended by the $`\mathit{}_p`$ decay of the LSP through $`\lambda _{ijk}^{}`$, and if the two LSP’s undergo the same decay $`\stackrel{~}{\chi }_1^0l_iu_j\overline{d}_k`$ or $`\stackrel{~}{\chi }_1^0\overline{l}_i\overline{u}_jd_k`$, two same sign charged leptons are generated. Another possible way for the SUSY pair production to generate the like sign dilepton signature is that only one of the LSP’s decays into a charged lepton of a given sign, the other decaying as $`\stackrel{~}{\chi }_1^0\nu _id_jd_k`$, and a second charged lepton of the same sign is produced in the cascade decays.
The cross sections of the superpartners pair productions have been studied in Section 4.3.
### 5.4 Cuts
In order to simulate the single chargino productions $`p\overline{p}\stackrel{~}{\chi }_1^\pm l^{}`$, $`p\overline{p}\stackrel{~}{\chi }_1^\pm \nu `$ and the single neutralino production $`p\overline{p}\stackrel{~}{\chi }_1^0l^{}`$ at Tevatron, the matrix elements (see Appendix A) of these processes have been implemented in a version of the SUSYGEN event generator allowing the generation of $`p\overline{p}`$ reactions . The Standard Model background ($`W^\pm W^{}`$, $`W^\pm Z^0`$, $`Z^0Z^0`$, $`t\overline{b}/\overline{t}b`$, $`t\overline{t}`$ and $`b\overline{b}`$ productions) has been simulated using the PYTHIA event generator and the SUSY background (all SUSY particles pair productions) using the HERWIG event generator . SUSYGEN, PYTHIA and HERWIG have been interfaced with the SHW detector simulation package (see Section 4.4).
Several selection criteria have been applied in order to reduce the background.
First, we have selected the events containing two same sign muons. The reason is that in the like sign dilepton signature analysis we have focused on the case of a single dominant $`\mathit{}_p`$ coupling constant of the type $`\lambda _{2jk}^{}`$. In such a scenario, the two same charge leptons generated in the $`\stackrel{~}{\chi }_1^0l^{}`$ production, which represents the main contribution to the like sign dilepton final state (see Section 5.1), are muons (see Fig.1 and Section 5.1). This requirement that the 2 like sign leptons have the same flavour allows to reduce the Standard Model background with respect to the signal.
We require a number of jets greater or equal to two with a transverse momentum higher than $`10GeV`$, namely $`N_j2`$ with $`P_t\left(j\right)>10GeV`$. This jet veto reduces the non-physics backgrounds generated by the $`W^\pm Z^0`$ and $`Z^0Z^0`$ productions (see Section 5.2) which produce at most one hard jet (see Section 4.4).
Besides, some effective cuts concerning the energies of the 2 selected muons have been applied. In Fig.20, we present the distributions of the 2 muon energies in the like sign dimuon events generated by the Standard Model background ($`W^\pm W^{}`$, $`W^\pm Z^0`$, $`Z^0Z^0`$, $`t\overline{t}`$, $`t\overline{b}/\overline{t}b`$ and $`b\overline{b}`$) and the SUSY signal. Based on these distributions, we have chosen the following cuts on the muon energies: $`E\left(\mu _2\right)>20GeV`$ and $`E\left(\mu _1\right)>20GeV`$.
We will refer to all the selection criteria described above, namely 2 same sign muons with $`E\left(\mu _2\right)>20GeV`$ and $`E\left(\mu _1\right)>20GeV`$, and $`N_j2`$ with $`P_t\left(j\right)>10GeV`$, as cut $`1`$.
Let us explain the origin of the two peaks in the upper left plot of Fig.20. This will be helpful for the mass reconstruction study of Section 5.5.2.
The main contribution to the like sign dimuon signature from the SUSY signal is the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production (see Section 5.1) in the case of a single dominant $`\lambda _{2jk}^{}`$ coupling. Furthermore, the dominant contribution to this production is the reaction $`p\overline{p}\stackrel{~}{\mu }_L^\pm \stackrel{~}{\chi }_1^0\mu ^\pm `$. In this reaction, the $`\mu ^\pm `$ produced together with the $`\stackrel{~}{\chi }_1^0`$ has an energy around $`E\left(\mu ^\pm \right)\left(m_{\stackrel{~}{\mu }_L^\pm }^2+m_{\mu ^\pm }^2m_{\stackrel{~}{\chi }_1^0}^2\right)/2m_{\stackrel{~}{\mu }_L^\pm }=121.9GeV`$ for the SUSY point considered in Fig.20, namely $`M_2=250GeV`$, $`m_0=200GeV`$, $`\mathrm{tan}\beta =1.5`$ and $`sign\left(\mu \right)<0`$, which gives rise to the mass spectrum: $`m_{\stackrel{~}{\chi }_1^0}=127.1GeV`$, $`m_{\stackrel{~}{\chi }_2^0}=255.3GeV`$, $`m_{\stackrel{~}{\chi }_1^\pm }=255.3GeV`$, $`m_{\stackrel{~}{l}_L^\pm }=298GeV`$ and $`m_{\stackrel{~}{\nu }_L^\pm }=294GeV`$. This energy value corresponds approximatively to the mean value of the right peak of the leading muon energy distribution presented in the upper left plot of Fig.20. This is due to the fact that the leading muon in the dimuon events generated by the reaction $`p\overline{p}\stackrel{~}{\chi }_1^0\mu ^\pm `$ is the $`\mu ^\pm `$ produced together with the $`\stackrel{~}{\chi }_1^0`$ for relatively important values of the $`m_{\stackrel{~}{\mu }_L^\pm }m_{\stackrel{~}{\chi }_1^0}`$ mass difference. The right peak in the upper left plot of Fig.20 is thus associated to the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production.
Similarly, the left peak in the upper left plot of Fig.20 corresponds to the reactions $`p\overline{p}\stackrel{~}{\mu }_L^\pm \stackrel{~}{\chi }_2^0\mu ^\pm `$ and $`p\overline{p}\stackrel{~}{\nu }_{\mu L}\stackrel{~}{\chi }_1^\pm \mu ^{}`$ which produce $`\mu ^\pm `$ of energies around $`E\left(\mu ^\pm \right)\left(m_{\stackrel{~}{\mu }_L^\pm }^2+m_{\mu ^\pm }^2m_{\stackrel{~}{\chi }_2^0}^2\right)/2m_{\stackrel{~}{\mu }_L^\pm }=39.6GeV`$ and $`E\left(\mu ^\pm \right)\left(m_{\stackrel{~}{\nu }_{\mu L}}^2+m_{\mu ^\pm }^2m_{\stackrel{~}{\chi }_1^\pm }^2\right)/2m_{\stackrel{~}{\nu }_{\mu L}}=36.2GeV`$, respectively. The $`\stackrel{~}{\chi }_1^\pm \nu _\mu `$ production represents a less important contribution to the like sign dimuon events compared to the 3 above single gaugino productions since the 2 same sign leptons generated in this production are not systematically muons and the involved branching ratios have smaller values (see Section 5.1).
Finally, since the leptons produced in the quark $`b`$ decays are not well isolated (as in the $`W^\pm W^{}`$, $`W^\pm Z^0`$, $`Z^0Z^0`$, $`t\overline{t}`$, $`t\overline{b}/\overline{t}b`$ and $`b\overline{b}`$ productions), we have applied some cuts on the lepton isolation. We have imposed the isolation cut $`\mathrm{\Delta }R=\sqrt{\delta \varphi ^2+\delta \theta ^2}>0.4`$ where $`\varphi `$ is the azimuthal angle and $`\theta `$ the polar angle between the 2 same sign muons and the 2 hardest jets. This cut is for example motivated by the distributions shown in Fig.21 of the $`\mathrm{\Delta }R`$ angular difference between the second leading muon and the second leading jet, in the like sign dimuons events generated by the SUSY signal and Standard Model background. We call cut $`\mathrm{\Delta }R>0.4`$ together with cut $`1`$, cut $`2`$.
In order to eliminate poorly isolated muons, we have also imposed that $`E<2GeV`$, where $`E`$ represents the summed energies of the jets being close to a muon, namely the jets contained in the cone centered on a muon and defined by $`\mathrm{\Delta }R<0.25`$. This cut is for instance motivated by the distributions shown in Fig.22 which represent the summed energies $`E`$ of the jets being close to the second leading muon in the like sign dimuons events generated by the SUSY signal and Standard Model background. We denote cut $`E<2GeV`$ plus cut $`2`$ as cut $`3`$.
The selected events require high energy charged leptons and jets and can thus be easily triggered at Tevatron. Moreover, the considered charged leptons and jets are typically emitted at intermediate polar angles and would thus be often detected at Tevatron. These points are illustrated in Fig.23 where are shown the energy and polar angle distributions of the leading muon and the leading jet in the like sign dimuons events selected by applying cut $`3`$ and generated by the SUSY signal and Standard Model background.
In Table 6, we give the numbers of like sign dilepton events expected from the Standard Model background at Tevatron Run II with the various cuts described above. We see in Table 6 that the main source of Standard Model background to the like sign dilepton signature at Tevatron is the $`t\overline{t}`$ production. This is due to its important cross section compared to the other Standard Model backgrounds (see Section 5.2) and to the fact that in the $`t\overline{t}`$ background, in contrast with the $`b\overline{b}`$ background, only one charged lepton of the final state is produced in a $`b`$-jet and is thus not isolated.
In Table 7, we give the number of like sign dilepton events generated by the SUSY background (all superpartners pair productions) at Tevatron Run II as a function of the $`m_0`$ and $`m_{1/2}`$ parameters for cut 3. This number of events decreases as $`m_0`$ and $`m_{1/2}`$ increase due to the behaviour of the summed superpartners pair production cross section in the SUSY parameter space (see Section 4.3).
### 5.5 Results
#### 5.5.1 Discovery potential
We first present the reach in the mSUGRA parameter space obtained from the analysis of the like sign dilepton final state at Tevatron Run II produced by the single neutralino and chargino productions via $`\lambda _{211}^{}`$: $`p\overline{p}\stackrel{~}{\chi }_{1,2}^0\mu ^\pm `$, $`p\overline{p}\stackrel{~}{\chi }_1^\pm \mu ^{}`$ and $`p\overline{p}\stackrel{~}{\chi }_1^\pm \nu _\mu `$. The sensitivities that can be obtained on the $`\lambda _{2jk}^{}`$ ($`j`$ and $`k`$ being not equal to $`1`$ simultaneously), $`\lambda _{1jk}^{}`$ and $`\lambda _{3jk}^{}`$ coupling constants will be discussed at the end of this section.
In Fig.24, we present the $`3\sigma `$ and $`5\sigma `$ discovery contours and the limits at $`95\%`$ confidence level in the plane $`m_0`$ versus $`m_{1/2}`$, for $`sign\left(\mu \right)<0`$, $`\mathrm{tan}\beta =1.5`$, $`\lambda _{211}^{}=0.05`$ and using a set of values for the luminosity. Those discovery potentials were obtained by considering the $`\stackrel{~}{\chi }_{1,2}^0\mu ^\pm `$, $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ and $`\stackrel{~}{\chi }_1^\pm \nu _\mu `$ productions and the background originating from the Standard Model. The signal and background were selected by using cut $`3`$ described in Section 5.4. The reduction of the sensitivity on $`m_{1/2}`$ observed in Fig.24 as $`m_0`$ increases is due to the decrease of the $`\stackrel{~}{\chi }_{1,2}^0\mu ^\pm `$, $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ and $`\stackrel{~}{\chi }_1^\pm \nu _\mu `$ productions cross sections with the $`m_0`$ increase observed in Fig.5 and Fig.6. In Fig.24, we also see that the sensitivity on $`m_{1/2}`$ is reduced in the domain $`m_0\stackrel{<}{}200GeV`$. This reduction of the sensitivity is due to the fact that in mSUGRA at low $`\mathrm{tan}\beta `$ and for large values of $`m_{1/2}`$ and small values of $`m_0`$, the LSP is the Right slepton $`\stackrel{~}{l}_{iR}^\pm `$ ($`i=1,2,3`$). Therefore, in this mSUGRA region the dominant decay channel of the lightest neutralino is $`\stackrel{~}{\chi }_1^0\stackrel{~}{l}_{iR}^\pm l_i^{}`$ ($`i=1,2,3`$) so that the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production, which is the main contribution to the like sign dilepton signature, leads to the $`2\mu ^\pm +2jets`$ final state only in a few cases. There are two reasons. First, in this mSUGRA scenario the charged lepton produced in the main $`\stackrel{~}{\chi }_1^0`$ decay is not systematically a muon. Secondly, if the LSP is the Right slepton $`\stackrel{~}{l}_{iR}^\pm `$ it cannot decay in the case of a single dominant $`\lambda _{ijk}^{}`$ coupling constant and it is thus a stable particle.
The sensitivities presented in the discovery reach of Fig.24 which are obtained from the like sign dilepton signature analysis are higher than the sensitivities shown in Fig.10 which correspond to the trilepton final state analysis. This is due to the 3 following points. First, the rate of the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production (recall that it represents the main contribution to the like sign dilepton final state) is larger than the $`\sigma \left(p\overline{p}\stackrel{~}{\chi }_1^\pm \mu ^{}\right)`$ cross section in most of the mSUGRA parameter space (see Section 3.1.1). Secondly, the $`\stackrel{~}{\chi }_1^0`$ decay leading to the like sign dilepton final state in the case of the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production has a larger branching ratio than the cascade decay initiated by the $`\stackrel{~}{\chi }_1^\pm `$ which generates the trilepton final state (see Sections 4.1 and 5.1). Finally, at Tevatron Run II the Standard Model background of the like sign dilepton signature is weaker than the trilepton Standard Model background (see Tables 3 and 7).
It is clear from Fig.24 that at low values of the $`m_0`$ and $`m_{1/2}`$ parameters, high sensitivities can be obtained on the $`\lambda _{211}^{}`$ coupling constant. We have found that for instance at the mSUGRA point defined as $`m_0=200GeV`$, $`m_{1/2}=200GeV`$, $`sign\left(\mu \right)<0`$ and $`\mathrm{tan}\beta =1.5`$, $`\lambda _{211}^{}`$ values of $`0.03`$ can be probed through the like sign dilepton analysis at Tevatron Run II assuming a luminosity of $`=1fb^1`$. This result was obtained by applying cut $`3`$ described in Section 5.4 on the SUSY signal ($`\stackrel{~}{\chi }_{1,2}^0\mu ^\pm `$, $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ and $`\stackrel{~}{\chi }_1^\pm \nu _\mu `$ productions) and the Standard Model background.
We expect that, as in the three lepton signature analysis, interesting sensitivities could be obtained on other $`\lambda _{2jk}^{}`$ coupling constants.
The sensitivities obtained on the $`\lambda _{3jk}^{}`$ couplings from the like sign dilepton signature analysis should be weaker than the sensitivities on the $`\lambda _{2jk}^{}`$ couplings deduced from the same study. Indeed, in the case of a single dominant $`\lambda _{3jk}^{}`$ coupling the same sign leptons generated by the $`\stackrel{~}{\chi }_1^0\tau ^\pm `$ production would be 2 tau leptons (see Fig.1(d) and Section 5.1). Therefore, the like sign dileptons ($`e^\pm e^\pm `$ or $`\mu ^\pm \mu ^\pm `$) produced by the $`\mathit{}_p`$ signal would be mainly generated in tau decays and would thus have higher probabilities to not pass the analysis cuts on the particle energy. Moreover, the requirement of $`e^\pm e^\pm `$ or $`\mu ^\pm \mu ^\pm `$ events would decrease the efficiency after cuts of the $`\mathit{}_p`$ signal due to the hadronic decay of the tau. Finally, the selection of two same flavour like sign dileptons ($`e^\pm e^\pm `$ or $`\mu ^\pm \mu ^\pm `$) would reduce the $`\mathit{}_p`$ signal, since each of the 2 produced taus could decay either into an electron or a muon, and hence would not be an effective cut anymore.
The sensitivities obtained on the $`\lambda _{1jk}^{}`$ couplings from the like sign dilepton signature study are expected to be identical to the sensitivities on the $`\lambda _{2jk}^{}`$ couplings obtained from the same study. Indeed, in the case of a single dominant $`\lambda _{1jk}^{}`$ coupling constant, the only difference in the like sign dilepton signature analysis would be that $`e^\pm e^\pm `$ events should be selected instead of $`\mu ^\pm \mu ^\pm `$ events (see Fig.1(d) and Section 5.1). Nevertheless, a smaller number of $`\lambda _{1jk}^{}`$ couplings is expected to be probed since the low-energy constraints on the $`\lambda _{1jk}^{}`$ couplings are generally stronger than the limits on the $`\lambda _{2jk}^{}`$ couplings .
In the high $`\mathrm{tan}\beta `$ case, the lightest stau $`\stackrel{~}{\tau }_1`$ can become the LSP instead of the lightest neutralino, due to a large mixing in the third generation of charged sleptons. In such a situation, the dominant decay channel of the lightest neutralino is $`\stackrel{~}{\chi }_1^0\stackrel{~}{\tau }_1^\pm \tau ^{}`$. Two scenarios must then be discussed: if the single dominant $`\mathit{}_p`$ coupling is not of the type $`\lambda _{3jk}^{}`$, the $`\stackrel{~}{\tau }_1^\pm `$-LSP is a stable particle so that the reaction $`p\overline{p}\stackrel{~}{\chi }_1^0l_i^\pm `$, representing the main contribution to the like sign dilepton final state, does not often lead to the $`2\mu ^\pm +2jets`$ signature. If the single dominant $`\mathit{}_p`$ coupling is of the type $`\lambda _{3jk}^{}`$, the $`\stackrel{~}{\chi }_1^0\tau ^\pm `$ production can receive a contribution from the resonant $`\stackrel{~}{\tau }_2^\pm `$ production (see Fig.1(d)) and the $`\stackrel{~}{\tau }_1^\pm `$-LSP decays via $`\lambda _{3jk}^{}`$ as $`\stackrel{~}{\tau }_1^\pm u_jd_k`$ so that the $`2\mu ^\pm +2jets`$ signature can still be generated in a significant way by the $`p\overline{p}\stackrel{~}{\chi }_1^0\tau ^\pm `$ reaction.
We end this Section by some comments on the effect of the supersymmetric $`R_p`$ conserving background to the like sign dilepton signature. In order to illustrate this discussion, we consider the results on the $`\lambda _{211}^{}`$ coupling constant.
We see from Table 7 that the SUSY background to the like sign dilepton final state can affect the sensitivity on the $`\lambda _{211}^{}`$ coupling constant obtained by considering only the Standard Model background, which is shown in Fig.24, only in the region of small superpartners masses, namely in the domain $`m_{1/2}\stackrel{<}{}300GeV`$ for $`\mathrm{tan}\beta =1.5`$, $`sign\left(\mu \right)<0`$ and assuming a luminosity of $`=1fb^1`$.
In contrast with the SUSY signal amplitude which is increased if $`\lambda _{211}^{}`$ is enhanced, the SUSY background amplitude is typically independent on the value of the $`\lambda _{211}^{}`$ coupling constant since the superpartner pair production does not involve $`\mathit{}_p`$ couplings. Therefore, even if we consider the SUSY background in addition to the Standard Model one, it is still true that large values of the $`\lambda _{211}^{}`$ coupling can be probed over a wider domain of the SUSY parameter space than low values, as can be observed in Fig.24 for $`m_{1/2}\stackrel{>}{}300GeV`$. Note that in Fig.24 larger values of $`\lambda _{211}^{}`$ still respecting the indirect limit could have been considered.
Finally, we mention that further cuts, as for instance some cuts based on the superpartners mass reconstructions (see Section 5.5.2), could allow to reduce the SUSY background to the like sign dilepton signature.
#### 5.5.2 Mass reconstructions
The $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{l}_L^\pm `$ mass reconstructions can be performed in a model independent way via the like sign dilepton analysis. We have simulated these mass reconstructions based on the like sign dimuon events generated in the scenario of a single dominant $`\lambda _{2jk}^{}`$ coupling constant. In this scenario, the main SUSY contribution to the like sign dilepton signature, namely the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production, has the final state $`\mu ^\pm +\mu ^\pm +2jets`$ (see Section 5.1). Indeed, the produced $`\stackrel{~}{\chi }_1^0`$ decays into $`\mu ^\pm u_jd_k`$ through $`\lambda _{2jk}^{}`$. The muon generated together with the $`\stackrel{~}{\chi }_1^0`$ can be identified as the leading muon for relatively large $`m_{\stackrel{~}{\mu }_L^\pm }m_{\stackrel{~}{\chi }_1^0}`$ mass differences (see Section 5.4). Note that for nearly degenerate values of $`m_{\stackrel{~}{\mu }_L^\pm }`$ and $`m_{\stackrel{~}{\chi }_1^0}`$ the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production rate and thus the sensitivity on the SUSY parameters would be reduced (see Section 3.1.1). The muon created in the $`\stackrel{~}{\chi }_1^0`$ decay can thus be identified as the softer muon so that the $`\stackrel{~}{\chi }_1^0`$ can be reconstructed from the the softer muon and the 2 jets present in the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production final state. The other contributions to the like sign dimuons events can lead to some missing energy and at most 4 jets in the final state (see Section 5.1). Hence, we have chosen to reconstruct the $`\stackrel{~}{\chi }_1^0`$ from the 2 leading jets when the final state contains more than 2 jets. Once the $`\stackrel{~}{\chi }_1^0`$ has been reconstructed, the $`\stackrel{~}{\mu }_L^\pm `$ has been reconstructed from the $`\stackrel{~}{\chi }_1^0`$ and the leading muon since the dominant contribution to the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production is the reaction $`p\overline{p}\stackrel{~}{\mu }_L^\pm \stackrel{~}{\chi }_1^0\mu ^\pm `$. These mass reconstructions are represented in Fig.25. In this figure, we also represent the same mass reconstructions obtained by applying a cut in the upper left plot of Fig.20 excluding the peak associated to the $`\stackrel{~}{\chi }_2^0\mu ^\pm `$ and $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ productions (see Section 5.4). The interest of this cut, as can be seen in Fig.25, is to select the $`\stackrel{~}{\chi }_1^0\mu ^\pm `$ production and thus to improve the accuracy on the $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\mu }_L^\pm `$ reconstructions which are based on this production. We observe in Fig.25 that the $`\stackrel{~}{\chi }_1^0`$ reconstruction has less combinatorial background than the $`\stackrel{~}{\mu }_L^\pm `$ reconstruction. This comes from the fact that the selection of the softer muon and the 2 leading jets allows to reconstruct the $`\stackrel{~}{\chi }_1^0`$ even in the dimuon events generated by the $`\stackrel{~}{\chi }_2^0\mu ^\pm `$ and $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ productions, while the selection of the 2 muons and the 2 leading jets does not allow to reconstruct the $`\stackrel{~}{\mu }_L^\pm `$ in the dimuon events generated by the $`\stackrel{~}{\chi }_2^0\mu ^\pm `$ and $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ productions (see Section 5.1). We have represented on the plots of Fig.25 the fits of the invariant mass distributions. We see from these fits that the distributions are well peaked around the $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\mu }_L^\pm `$ generated masses. The average reconstructed masses are $`m_{\stackrel{~}{\chi }_1^0}=116\pm 11GeV`$ and $`m_{\stackrel{~}{\mu }_L^\pm }=285\pm 20GeV`$.
We note that the accuracy on the $`\stackrel{~}{\chi }_1^0`$ (and thus on the $`\stackrel{~}{\mu }_L^\pm `$) mass reconstruction could be improved if the distributions in the upper plots of Fig.25 were recalculated by selecting the muon giving the $`\stackrel{~}{\chi }_1^0`$ mass the closer to the mean value of the peak obtained in the relevant upper plot of Fig.25.
In the hypothesis of a single dominant coupling constant of type $`\lambda _{1jk}^{}`$ or $`\lambda _{3jk}^{}`$, exactly the same kind of $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\mu }_L^\pm `$ mass reconstructions can be performed by selecting the $`e^\pm +e^\pm +jets+E/`$ or $`l_i^\pm +l_j^\pm +jets+E/`$ events, respectively.
As a conclusion, the $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\mu }_L^\pm `$ mass reconstructions based on the like sign dilepton signature generated by the $`\stackrel{~}{\chi }_{1,2}^0\mu ^\pm `$, $`\stackrel{~}{\chi }_1^\pm \mu ^{}`$ and $`\stackrel{~}{\chi }_1^\pm \nu _\mu `$ productions at Tevatron can easily give precise results, in contrast with the mass reconstructions performed in the superpartner pair production analysis at hadronic colliders which suffer an high combinatorial background .
#### 5.5.3 Model dependence of the results
In our theoretical framework (see Section 2), the values of the $`\left|\mu \right|`$ and $`\mathrm{tan}\beta `$ (up to the ambiguity of low/high solution) parameters are predicted. This has no important effects on the results presented in Sections 5.5.1 as the single gaugino production cross sections vary weakly with these parameters (see Section 3.1.1).
However, since we have worked within the mSUGRA model, the $`\stackrel{~}{l}_L^\pm `$ mass was typically larger than the $`\stackrel{~}{\chi }_1^0`$ mass. In a situation where $`m_{\stackrel{~}{l}_L^\pm }`$ would approach $`m_{\stackrel{~}{\chi }_1^0}`$, the rate of the $`\stackrel{~}{\chi }_1^0l_i^\pm `$ production, representing in mSUGRA the main contribution to the like sign dilepton signature (see Section 5.1), would decrease. Therefore, within a model allowing degenerate $`\stackrel{~}{l}_L^\pm `$ and $`\stackrel{~}{\chi }_1^0`$ masses or even a $`\stackrel{~}{l}_L^\pm `$ lighter than the $`\stackrel{~}{\chi }_1^0`$, other single gaugino productions than the $`p\overline{p}\stackrel{~}{\chi }_1^0l_i^\pm `$ reaction could represent the major contribution to the like sign dilepton signature in some parts of the SUSY parameter space.
Besides, in a situation where the LSP would not be the $`\stackrel{~}{\chi }_1^0`$, the branching ratios of the $`\stackrel{~}{\chi }_1^0`$ decays violating $`R_p`$ would be reduced with respect to the case where the LSP is the $`\stackrel{~}{\chi }_1^0`$, as often occurs in mSUGRA. However, in such a situation, the like sign dilepton signature could receive a significant contribution from a decay of the $`\stackrel{~}{\chi }_1^0`$ different from the $`\mathit{}_p`$ channel. In those kinds of scenarios where the LSP is not the $`\stackrel{~}{\chi }_1^0`$, the $`\stackrel{~}{\chi }_1^0l_i^\pm `$ production would not represent systematically the main contribution to the like sign dilepton signature.
In the several scenarios described above where the $`\stackrel{~}{\chi }_1^0l_i^\pm `$ production is not the major contribution to the like sign dilepton signature, this signature could receive quite important contribution from the other single gaugino productions described in Section 3.1.
## 6 Conclusion
The single gaugino productions at Tevatron reach important cross sections thanks to the contributions of the resonant slepton productions. Hence, the analysis of the 3 charged leptons and like sign dilepton signatures generated by the single gaugino productions at Tevatron Run II would allow to obtain high sensitivities on many $`\mathit{}_p`$ coupling constants, compared to the low-energy limits, in wide domains of the SUSY parameter space. This is also due to the fact that the Standard Model backgrounds associated to the 3 charged leptons and like sign dilepton final states at Tevatron can be greatly suppressed.
From the supersymmetry discovery point of view, superpartner masses well beyond the present experimental limits could be tested through the analysis of the the 3 charged leptons and like sign dilepton signatures generated by the single gaugino productions at Tevatron Run II. If some of the $`\mathit{}_p`$ coupling constants values were close to their low-energy bounds, the single gaugino productions study based on the 3 charged leptons and like sign dilepton signatures would even allow to extend the region in the $`m_0`$-$`m_{1/2}`$ plane probed by the superpartner pair production analyses in the 3 charged leptons and like sign dilepton channels at Tevatron Run II. The reason is that the single superpartner production has a larger phase space factor than the superpartner pair production.
Besides, the 3 charged leptons and like sign dilepton signatures generated by the single gaugino productions at Tevatron Run II would allow to reconstruct in a model independent way the $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_1^\pm `$, $`\stackrel{~}{\nu }_L`$ and $`\stackrel{~}{l}_L^\pm `$ masses with a smaller combinatorial background than in the superpartner pair production analysis.
We end this summary by a comparison between the results obtained from the studies of the 3 charged lepton and like sign dilepton signatures generated by the single gaugino productions at Tevatron Run II. In the mSUGRA model, the like sign dilepton signature analysis would give rise to higher sensitivities on the SUSY parameters than the study of the 3 charged lepton final state. This comes notably from the fact that in mSUGRA, the $`\stackrel{~}{\chi }_1^0`$ is lighter than the $`\stackrel{~}{\chi }_1^\pm `$ so that the cross section of the $`\stackrel{~}{\chi }_1^0l^\pm `$ production, which is the main contribution to the like sign dilepton signature, reaches larger values than the cross section of the $`\stackrel{~}{\chi }_1^\pm l^{}`$ production, representing the main contribution to the 3 charged lepton final state.
Other interesting prospective studies concerning hadronic colliders are the analyses of the single gaugino productions occuring through resonant squark productions via $`\lambda ^{\prime \prime }`$ coupling constants which we will perform in the next future.
## 7 Acknowledgments
We would like to thank Emmanuelle Perez, Robi Peschanski and Auguste Besson for fruitful discussions and reading the manuscript.
## Appendix A Formulas for spin summed amplitudes
In this Appendix, we give the amplitudes for all the single productions of supersymmetric particle at hadronic colliders, which can receive a contribution from a slepton resonant production. These single productions occur via the $`\mathit{}_p`$ coupling $`\lambda _{ijk}^{}`$ and correspond to the four reactions, $`q\overline{q}\stackrel{~}{\chi }_a^+\overline{\nu }_i`$, $`q\overline{q}\stackrel{~}{\chi }_a^0\overline{\nu }_i`$, $`q\overline{q}\stackrel{~}{\chi }_a^0\overline{l}_i`$, $`q\overline{q}\stackrel{~}{\chi }_a^{}\overline{l}_i`$. Each of those four processes receives contributions from both the t and u channel (see Fig.1) and have charge conjugated diagrams. Note also that the contributions coming from the exchange of a right squark in the u channel involve the higgsino components of the gauginos. These contributions, in the case of the single chargino production, do not interfere with the s channel slepton exchange since the initial or final states are different (see Fig.1). In the following, we give the formulas for the probability amplitudes, squared and summed over the polarizations. Our notations closely follow the notations of . In particular, the matrix elements $`N_{ij}^{}`$ are defined in the basis of the photino and the zino, as in .
$`|M_s(u^j\overline{d}^k\stackrel{~}{\chi }_a^+\overline{\nu }_i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}g^2|U_{a1}|^2}{12(sm_{\stackrel{~}{l}_L^i}^2)^2}}(m_{u^j}^2+m_{d^k}^2s)(m_{\stackrel{~}{\chi }_a^+}^2s)`$ (A.1)
$`|M_t(u^j\overline{d}^k\stackrel{~}{\chi }_a^+\overline{\nu }_i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}g^2}{12(tm_{\stackrel{~}{d}_L^j}^2)^2}}(m_{d^k}^2t)[(|U_{a1}|^2+{\displaystyle \frac{m_{u^j}^2|V_{a2}|^2}{2m_W^2\mathrm{sin}^2\beta }})(m_{u^j}^2+m_{\stackrel{~}{\chi }_a^+}^2t)`$ (A.2)
$``$ $`{\displaystyle \frac{4m_{u^j}^2m_{\stackrel{~}{\chi }_a^+}Re(U_{a1}V_{a2})}{\sqrt{2}m_W\mathrm{sin}\beta }}]`$ (A.3)
$`|M_u(u^k\overline{d}^j\stackrel{~}{\chi }_a^+\nu _i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}g^2m_{d^k}^2|U_{a2}|^2}{24m_W^2\mathrm{cos}\beta ^2(um_{\stackrel{~}{d}_R^k}^2)^2}}(m_{\stackrel{~}{\chi }_a^+}^2+m_{u^k}^2u)(m_{d^j}^2u)`$ (A.4)
$`2Re[M_sM_t^{}(\stackrel{~}{\chi }_a^+\overline{\nu }_i)]`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}g^2}{6(sm_{\stackrel{~}{l}_L^i}^2)(tm_{\stackrel{~}{d}_L^j}^2)}}[{\displaystyle \frac{|U_{a1}|^2}{2}}[(m_{u^j}^2+m_{\stackrel{~}{\chi }_a^+}^2t)(m_{d^k}^2t)`$ (A.5)
$`+`$ $`(m_{u^j}^2+m_{d^k}^2s)(m_{\stackrel{~}{\chi }_a^+}^2s)(m_{u^j}^2u)(m_{\stackrel{~}{\chi }_a^+}^2+m_{d^k}^2u)]`$ (A.6)
$``$ $`(m_{d^k}^2t){\displaystyle \frac{Re(U_{a1}V_{a2})m_{\stackrel{~}{\chi }_a^+}m_{u^j}^2}{\sqrt{2}m_W\mathrm{sin}\beta }}],`$ (A.7)
where, $`s=(p(u^j)p(\overline{d}_k))^2`$, $`t=(p(u^j)p(\stackrel{~}{\chi }_a^+))^2`$ and $`u=(p(\overline{d}^j)p(\nu _i))^2`$.
$`|M_s(d_j\overline{d}_k\stackrel{~}{\chi }_a^0\overline{\nu }_i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}g^2|N_{a2}^{}|^2}{24\mathrm{cos}^2\theta _W(sm_{\stackrel{~}{\nu }_L^i}^2)^2}}(sm_{d^k}^2m_{d^j}^2)(sm_{\stackrel{~}{\chi }_a^0}^2)`$ (A.8)
$`|M_t(d_j\overline{d}_k\stackrel{~}{\chi }_a^0\overline{\nu }_i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}g^2}{6(tm_{\stackrel{~}{d}_L^j}^2)^2}}(m_{d^k}^2t)[(m_{d^j}^2+m_{\stackrel{~}{\chi }_a^0}^2t)({\displaystyle \frac{g^2m_{d^j}^2|N_{a3}^{}|^2}{4m_W^2\mathrm{cos}^2\beta }}+{\displaystyle \frac{e^2}{9}}|N_{a1}^{}|^2`$ (A.9)
$`+`$ $`{\displaystyle \frac{g^2|N_{a2}^{}|^2(\mathrm{sin}^2\theta _W/31/2)^2}{\mathrm{cos}^2\theta _W}}{\displaystyle \frac{2egRe(N_{a1}^{}N_{a2}^{})(\mathrm{sin}^2\theta _W/31/2)}{3\mathrm{cos}\theta _W}})`$ (A.10)
$`+`$ $`{\displaystyle \frac{2m_{\stackrel{~}{\chi }_a^0}m_{d^j}^2g}{m_W\mathrm{cos}\beta }}({\displaystyle \frac{eRe(N_{a1}^{}N_{a3}^{})}{3}}+{\displaystyle \frac{gRe(N_{a2}^{}N_{a3}^{})}{\mathrm{cos}\theta _W}}({\displaystyle \frac{\mathrm{sin}^2\theta _W}{3}}{\displaystyle \frac{1}{2}}))]`$ (A.11)
$`|M_u(d_j\overline{d}_k\stackrel{~}{\chi }_a^0\overline{\nu }_i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}}{6(um_{\stackrel{~}{d}_R^k}^2)^2}}(m_{d^j}^2u)[(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u)({\displaystyle \frac{g^2m_{d^k}^2|N_{a3}^{}|^2}{4m_W^2\mathrm{cos}^2\beta }}+{\displaystyle \frac{e^2|N_{a1}^{}|^2}{9}}`$ (A.12)
$`+`$ $`{\displaystyle \frac{g^2\mathrm{sin}^4\theta _W|N_{a2}^{}|^2}{9\mathrm{cos}^2\theta _W}}{\displaystyle \frac{2egRe(N_{a1}^{}N_{a2}^{})\mathrm{sin}^2\theta _W}{9\mathrm{cos}\theta _W}})`$ (A.13)
$``$ $`{\displaystyle \frac{2m_{\stackrel{~}{\chi }_a^0}m_{d^k}^2g}{m_W\mathrm{cos}\beta }}({\displaystyle \frac{eRe(N_{a1}^{}N_{a3}^{})}{3}}+{\displaystyle \frac{g\mathrm{sin}^2\theta _WRe(N_{a2}^{}N_{a3}^{})}{3\mathrm{cos}\theta _W}})]`$ (A.14)
$`2Re[M_sM_t^{}(\stackrel{~}{\chi }_a^0\overline{\nu }_i)]`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}g}{12\mathrm{cos}\theta _W(sm_{\stackrel{~}{\nu }_L^i}^2)(tm_{\stackrel{~}{d}_L^j}^2)}}[(m_{d^k}^2t){\displaystyle \frac{m_{\stackrel{~}{\chi }_a^0}m_{d^j}^2gRe(N_{a2}^{}N_{a3}^{})}{m_W\mathrm{cos}\beta }}`$ (A.15)
$`+`$ $`({\displaystyle \frac{eRe(N_{a1}^{}N_{a2}^{})}{3}}+{\displaystyle \frac{g|N_{a2}^{}|^2}{\mathrm{cos}\theta _W}}({\displaystyle \frac{\mathrm{sin}^2\theta _W}{3}}{\displaystyle \frac{1}{2}}))[(m_{d^j}^2+m_{\stackrel{~}{\chi }_a^0}^2t)(m_{d^k}^2t)`$ (A.16)
$`+`$ $`(m_{d^j}^2+m_{d^k}^2s)(m_{\stackrel{~}{\chi }_a^0}^2s)(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u)(m_{d^j}^2u)]]`$ (A.17)
$`2Re[M_tM_u^{}(\stackrel{~}{\chi }_a^0\overline{\nu }_i)]`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}}{6(um_{\stackrel{~}{d}_R^k}^2)(tm_{\stackrel{~}{d}_L^j}^2)}}[(m_{d^k}^2t){\displaystyle \frac{gm_{\stackrel{~}{\chi }_a^0}m_{d^j}^2}{m_W\mathrm{cos}\beta }}({\displaystyle \frac{g\mathrm{sin}^2\theta _WRe(N_{a2}^{}N_{a3}^{})}{3\mathrm{cos}\theta _W}}{\displaystyle \frac{eRe(N_{a1}^{}N_{a3}^{})}{3}})`$ (A.18)
$`+`$ $`[(m_{d^j}^2u)(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u)+(m_{d^k}^2t)(m_{d^j}^2+m_{\stackrel{~}{\chi }_a^0}^2t)(m_{\stackrel{~}{\chi }_a^0}^2s)(m_{d^j}^2+m_{d^k}^2s)]`$ (A.20)
$`\left({\displaystyle \frac{egRe(N_{a1}^{}N_{a2}^{})}{3\mathrm{cos}\theta _W}}({\displaystyle \frac{2\mathrm{sin}^2\theta _W}{3}}{\displaystyle \frac{1}{2}})+{\displaystyle \frac{e^2|N_{a1}^{}|^2}{9}}+{\displaystyle \frac{g^2\mathrm{sin}^2\theta _W|N_{a2}^{}|^2}{3\mathrm{cos}^2\theta _W}}({\displaystyle \frac{\mathrm{sin}^2\theta _W}{3}}{\displaystyle \frac{1}{2}})\right)`$
$``$ $`{\displaystyle \frac{m_{\stackrel{~}{\chi }_a^0}m_{d^k}^2g}{m_W\mathrm{cos}\beta }}\left({\displaystyle \frac{eRe(N_{a1}^{}N_{a3}^{})}{3}}+{\displaystyle \frac{gRe(N_{a2}^{}N_{a3}^{})}{\mathrm{cos}\theta _W}}({\displaystyle \frac{\mathrm{sin}^2\theta _W}{3}}{\displaystyle \frac{1}{2}})\right)(m_{d^j}^2u)`$ (A.21)
$`+`$ $`{\displaystyle \frac{m_{d^j}^2m_{d^k}^2g^2|N_{a3}^{}|^2}{2m_W^2\mathrm{cos}^2\beta }}(m_{\stackrel{~}{\chi }_a^0}^2s)]`$ (A.22)
$`2Re[M_sM_u^{}(\stackrel{~}{\chi }_a^0\overline{\nu }_i)]`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}g}{12\mathrm{cos}\theta _W(sm_{\stackrel{~}{\nu }_L^i}^2)(um_{\stackrel{~}{d}_R^k}^2)}}[{\displaystyle \frac{m_{\stackrel{~}{\chi }_a^0}m_{d^k}^2gRe(N_{a2}^{}N_{a3}^{})}{m_W\mathrm{cos}\beta }}(m_{d^j}^2u)`$ (A.23)
$`+`$ $`({\displaystyle \frac{eRe(N_{a1}^{}N_{a2}^{})}{3}}+{\displaystyle \frac{|N_{a2}^{}|^2g\mathrm{sin}^2\theta _W}{3\mathrm{cos}\theta _W}})[(m_{d^j}^2+m_{d^k}^2s)(m_{\stackrel{~}{\chi }_a^0}^2s)`$ (A.24)
$`+`$ $`(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u)(m_{d^j}^2u)(m_{d^j}^2+m_{\stackrel{~}{\chi }_a^0}^2t)(m_{d^k}^2t)]],`$ (A.25)
where, $`s=(p(d^j)p(\overline{d}_k))^2`$, $`t=(p(d^j)p(\stackrel{~}{\chi }_a^0))^2`$ and $`u=(p(d^j)p(\overline{\nu }_i))^2`$.
$`|M_s(u_j\overline{d}_k\stackrel{~}{\chi }_a^0\overline{l}_i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}}{6(sm_{\stackrel{~}{l}_L^i}^2)^2}}(sm_{u^j}^2m_{d^k}^2)[({\displaystyle \frac{g^2m_{l^i}^2|N_{a3}^{}|^2}{4m_W^2\mathrm{cos}^2\beta }}+e^2|N_{a1}^{}|^2+{\displaystyle \frac{g^2|N_{a2}^{}|^2}{\mathrm{cos}^2\theta _W}}(\mathrm{sin}^2\theta _W{\displaystyle \frac{1}{2}})^2`$ (A.26)
$``$ $`{\displaystyle \frac{2egRe(N_{a1}^{}N_{a2}^{})}{\mathrm{cos}\theta _W}}(\mathrm{sin}^2\theta _W{\displaystyle \frac{1}{2}}))(sm_{l^i}^2m_{\stackrel{~}{\chi }_a^0}^2){\displaystyle \frac{2gm_{\stackrel{~}{\chi }_a^0}m_{l^i}^2}{m_W\mathrm{cos}\beta }}(eRe(N_{a1}^{}N_{a3}^{})`$ (A.27)
$`+`$ $`{\displaystyle \frac{gRe(N_{a2}^{}N_{a3}^{})}{\mathrm{cos}\theta _W}}(\mathrm{sin}^2\theta _W{\displaystyle \frac{1}{2}}))]`$ (A.28)
$`|M_t(u_j\overline{d}_k\stackrel{~}{\chi }_a^0\overline{l}_i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}}{6(tm_{\stackrel{~}{u}_L^j}^2)^2}}(t+m_{l^i}^2+m_{d^k}^2)[({\displaystyle \frac{g^2m_{u^j}^2|N_{a4}^{}|^2}{4m_W^2\mathrm{sin}^2\beta }}+{\displaystyle \frac{4e^2|N_{a1}^{}|^2}{9}}`$ (A.29)
$`+`$ $`{\displaystyle \frac{g^2|N_{a2}^{}|^2}{\mathrm{cos}^2\theta _W}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{2\mathrm{sin}^2\theta _W}{3}})^2+{\displaystyle \frac{4egRe(N_{a1}^{}N_{a2}^{})}{3\mathrm{cos}\theta _W}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{2\mathrm{sin}^2\theta _W}{3}}))(t+m_{u^j}^2+m_{\stackrel{~}{\chi }_a^0}^2)`$ (A.30)
$`+`$ $`{\displaystyle \frac{2gm_{u^j}^2m_{\stackrel{~}{\chi }_a^0}}{m_W\mathrm{sin}\beta }}({\displaystyle \frac{2eRe(N_{a1}^{}N_{a4}^{})}{3}}+{\displaystyle \frac{gRe(N_{a2}^{}N_{a4}^{})}{\mathrm{cos}\theta _W}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{2\mathrm{sin}^2\theta _W}{3}}))]`$ (A.31)
$`|M_u(u_j\overline{d}_k\stackrel{~}{\chi }_a^0\overline{l}_i)|^2`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}}{6(um_{\stackrel{~}{d}_R^k}^2)^2}}(m_{u^j}^2+m_{l^i}^2u)[({\displaystyle \frac{e^2|N_{a1}^{}|^2}{9}}+{\displaystyle \frac{g^2\mathrm{sin}^4\theta _W|N_{a2}^{}|^2}{9\mathrm{cos}^2\theta _W}}{\displaystyle \frac{2egRe(N_{a1}^{}N_{a2}^{})\mathrm{sin}^2\theta _W}{9\mathrm{cos}\theta _W}}`$ (A.32)
$`+`$ $`{\displaystyle \frac{g^2m_{d^k}^2|N_{a3}^{}|^2}{4m_W^2\mathrm{cos}^2\beta }})(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u){\displaystyle \frac{2gm_{d^k}^2m_{\stackrel{~}{\chi }_a^0}}{m_W\mathrm{cos}\beta }}({\displaystyle \frac{eRe(N_{a1}^{}N_{a3}^{})}{3}}`$ (A.33)
$`+`$ $`{\displaystyle \frac{g\mathrm{sin}^2\theta _WRe(N_{a2}^{}N_{a3}^{})}{3\mathrm{cos}\theta _W}})]`$ (A.34)
$`2Re[M_sM_t^{}(\stackrel{~}{\chi }_a^0\overline{l}_i)]`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}}{6(sm_{\stackrel{~}{l}_L^i}^2)(tm_{\stackrel{~}{u}_L^j}^2)}}[{\displaystyle \frac{m_{l^i}^2m_{u^j}^2g^2Re(N_{a3}^{}N_{a4}^{})}{2m_W^2\mathrm{sin}\beta \mathrm{cos}\beta }}(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u)`$ (A.35)
$`+`$ $`({\displaystyle \frac{2e^2|N_{a1}^{}|^2}{3}}+{\displaystyle \frac{egRe(N_{a1}^{}N_{a2}^{})}{3\mathrm{cos}\theta _W}}(4\mathrm{sin}^2\theta _W{\displaystyle \frac{5}{2}})`$ (A.36)
$`+`$ $`{\displaystyle \frac{g^2|N_{a2}^{}|^2}{\mathrm{cos}^2\theta _W}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{2\mathrm{sin}^2\theta _W}{3}})(\mathrm{sin}^2\theta _W{\displaystyle \frac{1}{2}}))`$ (A.38)
$`[(m_{u^j}^2+m_{d^k}^2s)(m_{\stackrel{~}{\chi }_a^0}^2+m_{l^i}^2s)+(m_{u^j}^2+m_{\stackrel{~}{\chi }_a^0}^2t)(m_{l^i}^2+m_{d^k}^2t)`$
$``$ $`(m_{u^j}^2+m_{l^i}^2u)(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u)]+{\displaystyle \frac{gm_{u^j}^2m_{\stackrel{~}{\chi }_a^0}}{m_W\mathrm{sin}\beta }}(eRe(N_{a1}^{}N_{a4}^{})+{\displaystyle \frac{gRe(N_{a2}^{}N_{a4}^{})}{\mathrm{cos}\theta _W}}`$ (A.40)
$`(\mathrm{sin}^2\theta _W{\displaystyle \frac{1}{2}}))(m_{l^i}^2+m_{d^k}^2t)(sm_{u^j}^2m_{d^k}^2){\displaystyle \frac{gm_{l^i}^2m_{\stackrel{~}{\chi }_a^0}}{m_W\mathrm{cos}\beta }}({\displaystyle \frac{2eRe(N_{a1}^{}N_{a3}^{})}{3}}`$
$`+`$ $`{\displaystyle \frac{gRe(N_{a2}^{}N_{a3}^{})}{\mathrm{cos}\theta _W}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{2\mathrm{sin}^2\theta _W}{3}}))]`$ (A.41)
$`2Re[M_tM_u^{}(\stackrel{~}{\chi }_a^0\overline{l}_i)]`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}}{6(um_{\stackrel{~}{d}_R^k}^2)(tm_{\stackrel{~}{u}_L^j}^2)}}[{\displaystyle \frac{gm_{u^j}^2m_{\stackrel{~}{\chi }_a^0}}{m_W\mathrm{sin}\beta }}(m_{l^i}^2+m_{d^k}^2t)({\displaystyle \frac{eRe(N_{a1}^{}N_{a4}^{})}{3}}`$ (A.42)
$`+`$ $`{\displaystyle \frac{g\mathrm{sin}^2\theta _WRe(N_{a2}^{}N_{a4}^{})}{3\mathrm{cos}\theta _W}}){\displaystyle \frac{m_{\stackrel{~}{\chi }_a^0}gm_{d^k}^2}{m_W\mathrm{cos}\beta }}({\displaystyle \frac{2eRe(N_{a1}^{}N_{a3}^{})}{3}}+{\displaystyle \frac{gRe(N_{a2}^{}N_{a3}^{})}{\mathrm{cos}\theta _W}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{2\mathrm{sin}^2\theta _W}{3}}))`$ (A.44)
$`(m_{l^i}^2+m_{u^j}^2u){\displaystyle \frac{g^2Re(N_{a3}^{}N_{a4}^{})m_{u^j}^2m_{d^k}^2}{2m_W^2\mathrm{cos}\beta \mathrm{sin}\beta }}(sm_{l^i}^2m_{\stackrel{~}{\chi }_a^0}^2)+({\displaystyle \frac{2e^2|N_{a1}^{}|^2}{9}}`$
$`+`$ $`{\displaystyle \frac{egRe(N_{a1}^{}N_{a2}^{})}{3\mathrm{cos}\theta _W}}({\displaystyle \frac{1}{2}}+{\displaystyle \frac{4\mathrm{sin}^2\theta _W}{3}})+{\displaystyle \frac{g^2\mathrm{sin}^2\theta _W|N_{a2}^{}|^2}{3\mathrm{cos}^2\theta _W}}`$ (A.46)
$`({\displaystyle \frac{1}{2}}{\displaystyle \frac{2\mathrm{sin}^2\theta _W}{3}}))[(m_{l^i}^2+m_{u^j}^2u)(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u)+(m_{l^i}^2+m_{d^k}^2t)(m_{\stackrel{~}{\chi }_a^0}^2+m_{u^j}^2t)`$
$``$ $`(m_{l^i}^2+m_{\stackrel{~}{\chi }_a^0}^2s)(m_{d^k}^2+m_{u^j}^2s)]]`$ (A.47)
$`2Re[M_sM_u^{}(\stackrel{~}{\chi }_a^0\overline{l}_i)]`$ $`=`$ $`{\displaystyle \frac{\lambda _{ijk}^{}{}_{}{}^{2}}{6(sm_{\stackrel{~}{l}_L^i}^2)(um_{\stackrel{~}{d}_R^k}^2)}}[{\displaystyle \frac{gm_{l^i}^2m_{\stackrel{~}{\chi }_a^0}}{m_W\mathrm{cos}\beta }}({\displaystyle \frac{eRe(N_{a1}^{}N_{a3}^{})}{3}}+{\displaystyle \frac{g\mathrm{sin}^2\theta _WRe(N_{a2}^{}N_{a3}^{})}{3\mathrm{cos}\theta _W}})`$ (A.50)
$`(sm_{d^k}^2m_{u^j}^2){\displaystyle \frac{gm_{d^k}^2m_{\stackrel{~}{\chi }_a^0}}{m_W\mathrm{cos}\beta }}\left(eRe(N_{a1}^{}N_{a3}^{})+{\displaystyle \frac{gRe(N_{a2}^{}N_{a3}^{})}{\mathrm{cos}\theta _W}}(\mathrm{sin}\theta _W^2{\displaystyle \frac{1}{2}})\right)`$
$`(m_{l^i}^2+m_{u^j}^2u)+{\displaystyle \frac{g^2m_{l^i}^2m_{d^k}^2|N_{a3}^{}|^2}{2m_W^2\mathrm{cos}^2\beta }}(m_{\stackrel{~}{\chi }_a^0}^2+m_{u^j}^2t)+({\displaystyle \frac{e^2|N_{a1}^{}|^2}{3}}`$
$``$ $`{\displaystyle \frac{egRe(N_{a1}^{}N_{a2}^{})}{3\mathrm{cos}\theta _W}}(2\mathrm{sin}\theta _W^2{\displaystyle \frac{1}{2}})+{\displaystyle \frac{g^2|N_{a2}^{}|^2\mathrm{sin}^2\theta _W}{3\mathrm{cos}^2\theta _W}}(\mathrm{sin}^2\theta _W{\displaystyle \frac{1}{2}}))`$ (A.52)
$`[(m_{l^i}^2+m_{u^j}^2u)(m_{\stackrel{~}{\chi }_a^0}^2+m_{d^k}^2u)(m_{l^i}^2+m_{d^k}^2t)(m_{\stackrel{~}{\chi }_a^0}^2+m_{u^j}^2t)`$
$`+`$ $`(m_{l^i}^2+m_{\stackrel{~}{\chi }_a^0}^2s)(m_{d^k}^2+m_{u^j}^2s)],]`$ (A.53)
where, $`s=(p(u^j)p(\overline{d}_k))^2`$, $`t=(p(u^j)p(\stackrel{~}{\chi }_a^0))^2`$ and $`u=(p(u^j)p(\overline{l}_i))^2`$.
$`|M_s(d_j\overline{d}_k\stackrel{~}{\chi }_a^{}\overline{l}_i)|^2`$ $`=`$ $`{\displaystyle \frac{g^2\lambda _{ijk}^{}{}_{}{}^{2}}{6(sm_{\stackrel{~}{\nu }_L^i}^2)^2}}(sm_{d^j}^2m_{d^k}^2)[({\displaystyle \frac{m_{l^i}^2|U_{a2}|^2}{4m_W^2\mathrm{cos}^2\beta }}+{\displaystyle \frac{|V_{a1}|^2}{2}})(sm_{\stackrel{~}{\chi }_a^+}^2m_{l^i}^2)`$ (A.54)
$`+`$ $`{\displaystyle \frac{\sqrt{2}Re(V_{a1}U_{a2})m_{l^i}^2m_{\stackrel{~}{\chi }_a^+}}{m_W\mathrm{cos}\beta }}]`$ (A.55)
$`|M_t(d_j\overline{d}_k\stackrel{~}{\chi }_a^{}\overline{l}_i)|^2`$ $`=`$ $`{\displaystyle \frac{g^2\lambda _{ijk}^{}{}_{}{}^{2}}{3(tm_{\stackrel{~}{u}_L^j}^2)^2}}(tm_{d^k}^2m_{l^i}^2)[(tm_{\stackrel{~}{\chi }_a^+}^2m_{d^j}^2)({\displaystyle \frac{|V_{a1}|^2}{4}}+{\displaystyle \frac{m_{d^j}^2|U_{a2}|^2}{8M_W^2\mathrm{cos}^2\beta }})`$ (A.56)
$`+`$ $`{\displaystyle \frac{Re(V_{a1}U_{a2})m_{\stackrel{~}{\chi }_a^+}m_{d^j}^2}{\sqrt{2}m_W\mathrm{cos}\beta }}]`$ (A.57)
$`|M_u(\overline{u}_ku_j\stackrel{~}{\chi }_a^{}\overline{l}_i)|^2`$ $`=`$ $`{\displaystyle \frac{g^2\lambda _{ijk}^{}{}_{}{}^{2}}{24(um_{\stackrel{~}{d}_R^k}^2)^2}}(m_{\stackrel{~}{\chi }_a^+}^2+m_{u^k}^2u)(m_{l^i}^2+m_{u^j}^2u){\displaystyle \frac{|U_{a2}|^2m_{d^k}^2}{m_W^2\mathrm{cos}^2\beta }}`$ (A.58)
$`2Re[M_sM_t^{}(\stackrel{~}{\chi }_a^{}\overline{l}_i)]`$ $`=`$ $`{\displaystyle \frac{g^2\lambda _{ijk}^{}{}_{}{}^{2}}{12(sm_{\stackrel{~}{\nu }_L^i}^2)(tm_{\stackrel{~}{u}_L^j}^2)}}[|V_{a1}|^2[(m_{l^i}^2+m_{d^j}^2u)(m_{\stackrel{~}{\chi }_a^+}^2+m_{d^k}^2u)`$ (A.59)
$`+`$ $`(m_{l^i}^2+m_{d^k}^2t)(m_{\stackrel{~}{\chi }_a^+}^2+m_{d^j}^2t)+(m_{l^i}^2+m_{\stackrel{~}{\chi }_a^+}^2s)(m_{d^k}^2+m_{d^j}^2s)]`$ (A.60)
$`+`$ $`{\displaystyle \frac{Re(V_{a1}U_{a2})m_{\stackrel{~}{\chi }_a^+}\sqrt{2}}{m_W\mathrm{cos}\beta }}[m_{l^i}^2(sm_{d^j}^2m_{d^k}^2)m_{d^j}^2(m_{l^i}^2+m_{d^k}^2t)]`$ (A.61)
$``$ $`{\displaystyle \frac{|U_{a2}|^2m_{l^i}^2m_{d^j}^2}{m_W^2\mathrm{cos}^2\beta }}(m_{\stackrel{~}{\chi }_a^+}^2+m_{d^k}^2u)],`$ (A.62)
where, $`s=(p(d^j)p(\overline{d}_k))^2`$, $`t=(p(d^j)p(\stackrel{~}{\chi }_a^{}))^2`$ and $`u=(p(u^j)p(\overline{l}_i))^2`$.
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# Some conjectures about the Hilbert series of generic ideals in the exterior algebra
## 1. Introduction
In the symmetric algebra $`K[x_1,\mathrm{},x_n]`$, the set of Hilbert series coming from homogeneous quotients are classified by Macaulays theorem \[Mac27, Eis95, Big95\]. There is an infinite number of possible series, but if we fix positive integers $`d_1,\mathrm{},d_r`$, and restrict our study to quotients by homogeneous ideals $`I`$ of “type” or “numerical character” $`(d_1,\mathrm{},d_r)`$, ie generated by forms of those prescribed degrees, then there are only finitely many Hilbert series. Furthermore, in the affine space parametrising these homogeneous ideals, there is a Zariski-open subset of ideals with the same Hilbert series, and the Hilbert series obtained on this open set is minimal \[FL91, Frö85\].
Unfortunately, even though we know the set of *all* Hilbert series, we do not know what Hilbert series arise from ideals of numerical character $`(d_1,\mathrm{},d_r)`$. In fact, we do not even know the “generic” series, but it is conjectured \[MS91, Frö85\] that it is $`(1t)^n_{i=1}^r(1t^{d_i})`$; the brackets mean “truncate before the first non-positive coefficient”.
In the exterior algebra $`V_n`$, we also know the set of all Hilbert series of homogeneous quotients, by the so-called Kruskal-Katona theorem \[Kru63, Kat68, CL69, AHH97\]. Here, this set is finite, so one would think that it should be easy to find the subset of Hilbert series coming from quotients by ideals having a prescribed numerical character. In particular, it should be easy to find the generic value. However, very little is known.
In this article, we give one new result (the series for a quotient by *one* form of *even* degree) and several conjectures, supported by extensive computer calculations.
It is worthwhile to point out that the problem of determining the Hilbert series of quotients by generic *quadratic* forms is especially interesting, since it determines the Koszulness of the quadratic algebras in question. We refer to the recent article by Fröberg and Löfwall \[FL00\].
## 2. Notation
Let $`K`$ be a field of characteristic 0. Then $``$ is the prime subfield of $`K`$. For any positive integer $`n`$, let $`V=V_n`$ be an $`n`$-dimensional vector space over $`K`$, with a distinguished basis $`X_n=\{x_1,\mathrm{},x_n\}`$. Let $`K[x_1,\mathrm{},x_n]`$ denote the symmetric algebra on $`V_n`$, and let $`V_n`$ denote the exterior algebra on $`V_n`$. We define $`𝔖(V_n)`$, the *square-free* algebra on $`V_n`$, to be the commutative $`K`$-algebra generated by $`X_n`$, with the relations $`x_i^2=0`$; in other words, $`𝔖(V_n)=\frac{K[x_1,\mathrm{},x_n]}{(x_1^2,\mathrm{},x_n^2)}`$. There is an isomorphism of graded vector spaces between $`V_n`$ and $`𝔖(V_n)`$, but they are not isomorphic as $`K`$-algebras, since the exterior algebra is skew-commutative and $`𝔖(V_n)`$ is commutative.
We shall need the following operations for formal power series.
###### Definition 2.1.
Let $`f(t)=_{i=0}^{\mathrm{}}a_it^i[[t]]`$, $`g(t)=_{i=0}^{\mathrm{}}b_it^i[[t]]`$. We say that $`fg`$ if $`a_ib_i`$ for all $`i`$. We define
$`\mathrm{max}(f(t),g(t))`$ $`={\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}\mathrm{max}(a_i,b_i)`$
$`f(t)`$ $`={\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}a_it^i,\mathrm{}=\mathrm{max}(\left\{ia_j>0\text{ for }ji\right\})`$
$`f(t)`$ $`={\displaystyle \underset{i=\mathrm{}}{\overset{\mathrm{}}{}}}a_it^i,\mathrm{}=\mathrm{min}(\left\{ia_j>0\text{ for }ji\right\})`$
We use the conventions $`\mathrm{max}(\mathrm{})=1=\mathrm{min}()`$, $`\mathrm{min}(\mathrm{})=+\mathrm{}=\mathrm{max}()`$.
Let $`[X_n]`$ denote the free abelian monoid on $`X_n`$, and denote by $`Y_n`$ the subset of square-free monomials. Then $`Y_n`$ is a $`K`$-basis for both $`V_n`$ and $`𝔖(V_n)`$. We define the *degree* of a monomial in $`[X_n]`$ (and in $`Y_n`$) in the usual way, and denote by $`[X_n]^d`$ and $`Y_n^d`$ the subset of monomials (square-free monomials) of degree $`d`$.
A form $`K[x_1,\mathrm{},x_n]f=_{m[X_n]^d}c_mm`$ is said to be *generic* if the coefficients $`c_mK`$ fulfil the following conditions:
1. $`c_m`$,
2. $`mm^{}c_mc_m^{}`$,
3. The set of all $`c_m`$’s is algebraically independent over $``$.
A homogeneous ideal $`IK[x_1,\mathrm{},x_n]`$ is called generic if it can be minimally generated by a finite set of generic forms, so that all of the occuring coefficients of the forms are different, and so that the set of all occuring coefficients is algebraically independent over $``$. If the forms have degrees $`d_1,\mathrm{},d_r`$, then we say that $`I`$ has “numerical character” $`(d_1,\mathrm{},d_r)`$. It is an important fact that any two generic ideals of the same numerical character have the same initial ideal and the same Hilbert series.
Now consider the affine space $`V=𝐀^{\left(\genfrac{}{}{0pt}{}{n+d_11}{d_1}\right)}\times \mathrm{}\times 𝐀^{\left(\genfrac{}{}{0pt}{}{n+d_r1}{d_r}\right)}`$ parametrising the set of homogeneous ideals of numerical character $`(d_1,\mathrm{},d_r)`$. Since there are countably many conditions to be fulfilled for an ideal to be generic, the subset of the parameter space corresponding to generic ideals is not open, but a countable intersection of open sets, hence dense. However, in $`V`$ there is a Zariski-open subset corresponding to ideals with the same Hilbert function, and the generic ideals are contained in this subset \[FL91\].
We make similar definitions for the square-free algebra, and for the exterior algebra. Here, a generic form is a generic linear combination of *square-free* monomials of a certain degree. It is still true that the generic Hilbert series is attained on an open component of the parameter space, and that the generic ideals are contained in this component.
## 3. Hilbert series for generic principal ideals in the symmetric and square-free algebra
### 3.1. Principal ideals in the symmetric algebra
If $`fK[x_1,\mathrm{},x_n]`$ is a non-zero form of degree $`d`$, not necessarily homogeneous, then clearly the Hilbert series of the quotient $`\frac{K[x_1,\mathrm{},x_n]}{(f)}`$ is $`(1t)^n(1t^d)`$.
### 3.2. Principal ideals in the square-free algebra
If $`f𝔖(V_n)`$ is a generic form of degree $`d`$, then there is a similar simple formula for $`\frac{𝔖(V_n)}{(f)}(t)`$ (the Hilbert series of the quotient). To state the formula, we need some additional notation.
###### Definition 3.1.
We denote the zero series by $`0`$, and define
$`\mathrm{\Delta }_{n,d}(t)`$ $`=(t^d1)(1+t)^n`$
$`={\displaystyle \underset{v=(nd)/2}{\overset{n}{}}}\left(\left({\displaystyle \genfrac{}{}{0pt}{}{n}{v}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n}{v+d}}\right)\right)t^v`$
$`\delta _{n,d}(t)`$ $`=(1+t)^n(1t^d)`$
$`={\displaystyle \underset{v=0}{\overset{(nd)/2}{}}}(\left({\displaystyle \genfrac{}{}{0pt}{}{n}{v+d}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n}{v}}\right))t^v`$
The following result is due to Fröberg \[FH94\].
###### Theorem 3.2.
Let $`f𝔖(V_n)`$ be a generic form of degree $`d`$. Then
$$\frac{𝔖(V_n)}{(f)}(t)=\delta _{n,d}(t)$$
(1)
###### Proof.
By considering the graded exact sequence
$$0\mathrm{ann}(f)(d)𝔖(V_n)(d)\stackrel{f}{}𝔖(V_n)\frac{𝔖(V_n)}{(f)}0$$
(2)
in each degree $`r`$, we see that (1) holds if and only if multiplication by $`f`$, regarded as a linear map $`\varphi _r`$ from $`𝔖(V_n)_r`$ to $`𝔖(V_n)_{r+d}`$, is injective when $`\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{n}{r+d}\right)`$, and surjective when $`\left(\genfrac{}{}{0pt}{}{n}{r}\right)\left(\genfrac{}{}{0pt}{}{n}{r+d}\right)`$.
Write $`f=_{mY_n^d}c_mm`$. For $`0rnd`$, $`Y_n^r`$ is a basis of $`𝔖(V_n)_r`$, and $`Y_n^{r+d}`$ is a basis of $`𝔖(V_n)_{r+d}`$. Thus, we must show that for each $`r`$, the matrix of $`\varphi _r`$ in this basis has maximal rank. This matrix has rows indexed by $`Y_n^{r+d}`$ and columns indexed by $`Y_n^r`$. The entry at row $`R`$, column $`C`$ is
$$\{\begin{array}{cc}0\hfill & CR\hfill \\ c_m\hfill & R=mC\hfill \end{array}$$
If we specialise this matrix, the rank can only decrease, so if we can prove that some specialised matrix has full rank, then we are done. Putting all $`c_m=1`$, we obtain the incidence matrix of $`r`$-subsets of $`[n]`$ into $`r+d`$-subsets of $`[n]`$, that is, the rows are indexed by $`r`$-subsets and the columns by $`r+d`$-subsets, with a $`1`$ at the $`a,b`$’th position iff $`ab`$, and 0 otherwise. It has been shown by combinatorialists that this matrix has full rank \[Wil73, Kan72, GJ73\]. ∎
## 4. Principal ideals in the exterior algebra — the difference between even and odd degree
Let $`fV_n`$ be a generic form of degree $`d`$. Denote the Hilbert series of $`\frac{{\scriptscriptstyle V_n}}{f}`$ by $`q_{n,d}(t)`$, that of the annihilator of $`f`$ by $`a_{n,d}(t)`$, and that of the principal ideal $`(f)`$ by $`p_{n,d}(t)`$. From the the graded exact sequence
$$0\mathrm{ann}(f)(d)V_n(d)\stackrel{f}{}V_n\frac{V_n}{(f)}0$$
(3)
we get that
$$\begin{array}{cc}\hfill q_{n,d}(t)& =t^da_{n,d}(t)t^d(1+t)^n+(1+t)^n\hfill \\ & =t^da_{n,d}(t)+(1+t)^n(1t^d)\hfill \\ \hfill a_{n,d}(t)& =t^d\left(q_{n,d}(t)(1+t)^n(1t^d)\right)\hfill \end{array}$$
(4)
If $`d`$ is even, we shall prove that the vector space map
$$\stackrel{v}{}V_n\stackrel{f}{}\stackrel{v+d}{}V_n$$
(5)
is injective “when it can be”, ie when $`\left(\genfrac{}{}{0pt}{}{n}{v}\right)\left(\genfrac{}{}{0pt}{}{n}{v+d}\right)`$, and surjective “when it can be”, ie when $`\left(\genfrac{}{}{0pt}{}{n}{v}\right)\left(\genfrac{}{}{0pt}{}{n}{v+d}\right)`$. This leads immediately to the formulæ
$$\begin{array}{cc}\hfill q_{n,d}(t)& =(1+t)^n(1t^d)=\delta _{n,d}(t)\hfill \\ \hfill a_{n,d}(t)& =t^d\left(q_{n,d}(t)(1t^d)(1+t)^n\right)\hfill \\ & =t^d\left(\delta _{n,d}(t)(1t^d)(1+t)^n\right)\hfill \\ & =t^d\underset{r=0}{\overset{n}{}}\left[\mathrm{max}(0,\left(\genfrac{}{}{0pt}{}{n}{r+d}\right)\left(\genfrac{}{}{0pt}{}{n}{r}\right))\left(\left(\genfrac{}{}{0pt}{}{n}{r+d}\right)\left(\genfrac{}{}{0pt}{}{n}{r}\right)\right)\right]t^r\hfill \\ & =t^d\underset{r=0}{\overset{n}{}}\mathrm{max}(0,\left(\genfrac{}{}{0pt}{}{n}{r+d}\right)+\left(\genfrac{}{}{0pt}{}{n}{r}\right))t^r\hfill \\ & =t^d\mathrm{\Delta }_{n,d}(t)\hfill \end{array}$$
(6)
In particular, as $`n\mathrm{}`$, $`(1+t)^nq_{n,d}(t)(1t^d)`$, and $`a_{n,d}(t)0`$, with respect to the $`t`$-adic norm on $`[[t]]`$.
If $`d`$ is odd, then we have that $`f^2=0`$, hence $`fg=0`$ whenever $`g(f)`$, hence $`\mathrm{ann}(f)(f)`$, hence $`a_{n,d}(t)p_{n,d}(t)`$. In other words, there is a graded complex
$$\left(V\right)(d)\stackrel{f}{}V\stackrel{f}{}\left(V\right)(d)$$
(7)
the graded homology of which determines $`a_{n,d}(t)p_{n,d}(t)`$. In the (not very interesting) case $`d=1`$, then we know from \[AAH99\] that this homology vanishes. For odd $`d>1`$, we guess that for a fixed degree $`r`$, and $`n`$ very large, this homology vanishes. Hence, in degree $`r`$, the “obstruction to injectivity” in (5) is as small as possible. An equivalent formulation: consider the start of a minimal free graded $`V_n`$-resolution of $`\frac{{\scriptscriptstyle V_n}}{(f)}`$,
$$\frac{V_n}{(f)}V_n\stackrel{f}{}V_n\underset{j=1}{\overset{r}{}}(V_n)(\beta _{2,i}),$$
where $`\beta _{2,i}`$ are the graded Betti numbers. Then we guess that as $`n`$ increases, and for a fixed $`i2d`$, $`\beta _{2,i}=0`$. On the other hand, for sufficiently large $`n`$, we guess that $`\beta _{2,2d}=1`$. Since $`\beta _{2,i}`$ is the dimension of the degree $`id`$ part of a certain Tor group, this conjecture can also be stated in terms of Cartan homology (see \[AHH97\]).
We show the order (ie the smallest $`\mathrm{}`$ for which $`t^{\mathrm{}}`$ occurs with non-zero coefficient) of $`a_{n,d}(t)p_{n,d}(t)`$ for small $`n,d`$ in Table 1.
It would seem that the order of the difference grows linearly in $`n`$, so that $`a_{n,d}(t)p_{n,d}(t)0`$ rather rapidly.
Let us turn to the consequences of this conjecture. We get that $`a_{n,d}(t)p_{n,d}(t)`$ with respect to the $`(t)`$-adic filtration. It then follows from (4) that
$$q_{n,d}(t)t^dp_{n,d}(t)+(1+t)^n(1t^d)$$
(8)
Substituting $`p_{n,d}(t)=(1+t)^nq_{n,d}(t)`$ and solving for $`q_{n,d}(t)`$ we get that
$$q_{n,d}(t)\frac{(1+t)^nt^d+(1+t)^n(1t^d)}{(1+t^d)}=\frac{(1+t)^n}{(1+t^d)},$$
(9)
hence
$$\frac{\frac{{\scriptscriptstyle V_n}}{(f)}(t)}{(V_n)(t)}=\frac{q_{n,d}(t)}{(1+t)^n}\frac{1}{1+t^d}\text{ as }n\mathrm{}.$$
(10)
## 5. Principal ideals on generic forms of even degree in the exterior algebra
If $`d=2`$ then we can change coordinates on $`V`$ and replace $`f`$ with the form $`x_1x_2+x_3x_4+\mathrm{}`$, as is demonstrated in \[Bou59\]. The Hilbert series of the quotient can now be easily calculated. We get that $`\frac{{\scriptscriptstyle (V_n)}}{(f)}(t)=(1+t)^n(1t^2)`$, which is the same as the Hilbert series for the corresponding quotient in the square-free algebra.
###### Remark 5.1.
It is *not true* that if $`f_e=_{1i<jn}\alpha _{ij}x_ix_j`$ is a non-generic quadratic form in $`V_n`$, and $`f_s=_{1i<jn}\alpha _{ij}x_ix_j`$ is the corresponding form in $`𝔖(V_n)`$, then $`\frac{{\scriptscriptstyle V_n}}{(f_e)}`$ and $`\frac{𝔖(V_n)}{(f_s)}`$ have the same Hilbert series. For an example, consider the form $`x_1x_2+x_1x_3+x_1x_4+x_3x_4`$. The quotient of $`V_4`$ by this form has Hilbert series $`5t^2+4t+1,`$ but the corresponding quotient of $`𝔖(V_4)`$ has series $`t^3+5t^2+4t+1.`$
We next show that if the degree $`d`$ of $`f`$ is even, then the Hilbert series of the quotient $`\frac{{\scriptscriptstyle V_n}}{(f)}`$ is the same as for the square-free algebra. To this end, we need some combinatorial results, which we have collected in the appendix. With the aid of these, we can prove:
###### Theorem 5.2.
Let $`f^dV`$, with $`d`$ even, be a generic form. Then the linear transformation
$$^rV\stackrel{f}{}^{r+d}V$$
(11)
is injective for $`2r+dn`$, and surjective for $`2r+dn`$.
###### Proof.
We put $`k=r+d`$. Suppose that
$$f=\underset{K\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)}{}c_Kx_K$$
(12)
The matrix of the map (11) is an $`\left(\genfrac{}{}{0pt}{}{n}{r+d}\right)\times \left(\genfrac{}{}{0pt}{}{n}{d}\right)`$ matrix, $`\stackrel{~}{M_{r,r+d,n}}`$, where the rows are indexed by $`(r+d)`$-subsets $`K`$, and the columns by $`d`$-subsets $`T`$. The entry at position $`(K,T)`$ is
$$\{\begin{array}{cc}0\hfill & \text{ if }TK\hfill \\ \sigma (T,K)c_T\hfill & \text{ if }TK\hfill \end{array}$$
(13)
We must prove that this matrix has maximal rank. Clearly, the rank can not increase under specialisation, so if we prove that the matrix obtained by replacing each $`c_T`$ with 1 has maximal rank, then so does $`\stackrel{~}{M_{r,r+d,n}}`$. However, the specialised matrix is nothing but the matrix $`M_{r,r+d,n}`$ of Theorem A.6, so it has full rank. ∎
###### Theorem 5.3.
Let $`fV_n`$ be a generic form of degree $`d`$, with $`d`$ even. Then
$$\frac{V_n}{(f)}(t)=(1+t)^n(1t^d)=\delta _{n,d}(t)$$
(14)
###### Proof.
This follows from Theorem 5.2, together with (3). ∎
## 6. Principal ideals on generic forms of odd degree in the exterior algebra
Let $`d`$ be an odd integer. Recall that we’ve conjectured that $`a_{n,d}(t)p_{n,d}(t)0`$ as $`n\mathrm{}`$, and that this conjecture leads to the conclusions that $`p_{n,d}(t)(1+t)^n(1+t^d)^1`$. In this section, we shall try to guess the exact value of $`q_{n,d}(t)`$.
Since $`a_{n,d}(t)p_{n,d}(t)`$, $`a_{n,d}(t)\mathrm{\Delta }_{n,d}(t)`$, it follows that $`a_{n,d}(t)\mathrm{max}(p_{n,d}(t),\mathrm{\Delta }_{n,d}(t))`$. We tabulate the difference $`a_{n,d}(t)\mathrm{max}(p_{n,d}(t),\mathrm{\Delta }_{n,d}(t))`$ in Table 2 and Table 3.
Using the data of Table 3, we make the following conjecture:
###### Conjecture 6.1.
Let $`d`$ be an odd integer $`>3`$. Then, putting $`\tau _{n,d}(t)=a_{n,d}(t)\mathrm{max}(p_{n,d}(t),\mathrm{\Delta }_{n,d}(t))`$,
$$\tau _{n,d}(t)=\{\begin{array}{cc}t^{v(v1)/2}\hfill & v,s:v>0,nd=1+\frac{5}{2}v+\frac{1}{2}v^2,d=5+2vs\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}$$
(15)
This conjecture yields a formula for the Hilbert series, but since said formula is very complicated, we do not write it down; instead we show how to derive $`q_{n,d}(t)`$. From
$$\begin{array}{cc}\hfill a_{n,d}(t)& =\tau _{n,d}(t)+\mathrm{max}(p_{n,d}(t),\mathrm{\Delta }_{n,d}(t))\hfill \\ \hfill q_{n,d}(t)& =a_{n,d}(t)t^d+(1+t)^n(1t^d)\hfill \\ \hfill p_{n,d}(t)& =(1+t)^nq_{n,d}(t)\hfill \end{array}$$
(16)
we get
$$\begin{array}{cc}\hfill p_{n,d}(t)& =(1+t)^nq_{n,d}(t)\hfill \\ & =(1+t)^na_{n,d}(t)t^d(1+t)^n(1t^d)\hfill \\ & =(1+t)^nt^d\tau _{n,d}(t)t^d\mathrm{max}(p_{n,d}(t),\mathrm{\Delta }_{n,d}(t))(1+t)^n(1t^d)\hfill \\ & =t^d\left((1+t)^n\tau _{n,d}(t)\mathrm{max}(p_{n,d}(t),\mathrm{\Delta }_{n,d}(t))\right)\hfill \end{array}$$
(17)
Hence, writing $`p_{n,d}(t)=_{i=0}^na_it^i`$, with the $`a_i`$’s as undetermined coefficients, and denoting the $`t^i`$-coefficient of $`\tau _{n,d}(t)`$ by $`b_i`$, we get the equation
$$a_{\mathrm{}}=\left(\genfrac{}{}{0pt}{}{n}{\mathrm{}d}\right)b_i\mathrm{}\mathrm{max}(a_\mathrm{}d,\left(\genfrac{}{}{0pt}{}{n}{\mathrm{}d}\right)\left(\genfrac{}{}{0pt}{}{n}{\mathrm{}}\right))$$
(18)
which we can solve recursively, using the initial values
$$a_0=\mathrm{}=a_{d1}=0,a_d=a_n=1.$$
For the case $`d=3`$, we proceed differently: we tabulate $`q_{n,3}(t)w_{n,3}(t)`$ in Table 4, and from that, make the following conjecture:
###### Conjecture 6.2.
The Hilbert series of $`\frac{{\scriptscriptstyle V_n}}{(f)}`$, where $`f`$ is a generic cubic form, is given by
$$\begin{array}{cc}\hfill p_{n,3}(t)& =\frac{t^dL_n(t)+(1+t)^n}{1+t^d}\hfill \\ \hfill L_n(t)& =\{\begin{array}{cc}(3t)^{2\mathrm{}1}(1+t)^2\hfill & n=4\mathrm{}\hfill \\ c_1(n)t^{2\mathrm{}1}(1+t)(1+(3^{c_2(n)}1)t+t^2)\hfill & n=4\mathrm{}+1\hfill \\ (3t)^2\mathrm{}(1+t)^2\hfill & n=4\mathrm{}+2\hfill \\ (3t)^{2\mathrm{}+1}(1+t)\hfill & n=4\mathrm{}+3\hfill \end{array}\hfill \end{array}$$
(19)
where $`c_1(n),c_2(n)`$ are some positive integers.
## 7. Hilbert series for generic non-principal ideals in the symmetric and square-free algebra
Let $`I=(f_1,\mathrm{},f_r)`$ be a generic ideals in $`K[x_1,\mathrm{},x_n]`$, generated by forms of degree $`d_1,\mathrm{},d_r`$. There is a famous conjecture \[MS91, Frö85\] for the Hilbert series of the quotient $`\frac{K[x_1,\mathrm{},x_n]}{I_n}`$.
###### Conjecture 7.1.
Let $`I=(f_1,\mathrm{},d_r)K[x_1,\mathrm{},x_n]`$ be a generic ideal, with $`|f_i|=d_1`$ for $`1ir`$. Then the Hilbert series of the graded algebra $`\frac{K[x_1,\mathrm{},x_n]}{I_n}`$ is given by
$$(1t)^n\underset{i=1}{\overset{r}{}}(1t^{d_i})$$
(20)
It is easy to see that if $`rn`$, the generators form a regular sequence, and hence that
$$\frac{K[x_1,\mathrm{},x_n]}{I_n}(t)=(1t)^n\underset{i=1}{\overset{r}{}}(1t^{d_i}),\text{ for }nr$$
(21)
In particular, the conjecture holds for $`rn`$. The conjecture is also know to be true for $`r=n+1`$.
We note that (21) implies that
$$\underset{n\mathrm{}}{lim}\frac{\frac{K[x_1,\mathrm{},x_n]}{I_n}(t)}{K[x_1,\mathrm{},x_n](t)}=\underset{i=1}{\overset{r}{}}(1t^{d_i})$$
(22)
Now suppose that $`I=(f_1,\mathrm{},f_r)`$ is a generic ideal in the square-free algebra, and that $`f_i`$ is a generic form of degree $`d_i`$. Then
$$\frac{𝔖(V_n)}{(f_1,\mathrm{},f_r)}\frac{K[x_1,\mathrm{},x_n]}{(f_1^{},\mathrm{},f_r^{},x_1^2,\mathrm{},x_n^2)}$$
where $`f_i^{}`$ can be taken to be a generic form in $`K[x_1,\mathrm{},x_n]`$ which maps to $`f_i`$ under the canonical epimorphism $`K[x_1,\mathrm{},x_n]𝔖(V_n)`$. It seems reasonable to assume that the Hilbert series of the quotient will not change if we replace the squares of variables with generic quadratic forms. Conjecture 7.1 then leads to the following:
###### Conjecture 7.2.
Let $`r,n,d_1,\mathrm{},d_r`$, and let $`I_n`$ be a generic ideal i $`𝔖(V_n)`$ with generators of degrees $`d_1,\mathrm{},d_r`$. Then
$$\frac{𝔖(V_n)}{I_n}(t)=(1+t)^n\underset{i=1}{\overset{r}{}}(1t^{d_i})$$
(23)
If this conjecture holds (our computations support this), then it follows that
$$\underset{n\mathrm{}}{lim}\frac{\frac{𝔖(V_n)}{I_n}(t)}{𝔖(V_n)(t)}=\underset{i=1}{\overset{r}{}}(1t^{d_i})$$
(24)
This is analogous to (22).
## 8. Hilbert series for generic non-principal ideals in the exterior algebra
We now throw all caution to the wind to make some bold conjectures about the Hilbert series of non-principal generic ideals. Let $`I_n=(f_1,\mathrm{},f_r)`$ be a generic ideal in $`V_n`$, with $`|f_i|=d_i`$, and consider the exact sequence
$$0\mathrm{ann}(f_r)(d_r)\frac{V_n}{(f_1,\mathrm{},f_{r1})}(d_r)\stackrel{f_r}{}\frac{V_n}{(f_1,\mathrm{},f_{r1})}\frac{V_n}{(I)}0$$
(25)
We denote the Hilbert series of $`\frac{{\scriptscriptstyle V_n}}{(I)}`$ by $`q_n(t)`$, that of $`\frac{{\scriptscriptstyle V_n}}{(f_1,\mathrm{},f_{r1})}`$ by $`u_n(t)`$, and that of $`\mathrm{ann}(f_r)`$ by $`a_n(t)`$. Then
$$q_n(t)=u_n(t)t^{d_r}u_n(t)+t^da_n(t).$$
(26)
If $`d_r`$ is even, we conjecture that $`a_n(t)0`$, hence
$$q_n(t)(1t^{d_r})u_n(t)$$
(27)
If $`d_r`$ is odd, we conjecture that the annihilator of $`f_r`$ is “close” to the principal ideal on $`f_r`$, hence that $`a_n(t)(u_n(t)q_n(t))`$, which yields
$$q_n(t)(1+t^d)u_n(t)$$
(28)
By induction, we arrive at the following conjecture:
$$\underset{n\mathrm{}}{lim}\frac{q_n(t)}{(1+t)^n}=\underset{i=1}{\overset{r}{}}\left(1(1)^{d_i}t^{d_i}\right)^{(1)^{d_i}}[[t]],$$
(29)
with respect to the $`(t)`$-adic topology.
One would be tempted to guess that if all $`d_i`$’s are even, the Hilbert series of $`\frac{{\scriptscriptstyle V_n}}{(f_1,\mathrm{},f_r)}`$ should be *exactly*
$$(1+t)^n\underset{i=1}{\overset{r}{}}(1t^{d_i})$$
(30)
However, this is not true, even for the simplest case $`r=2`$ and $`d_1=d_2=2`$. In Table 5 we tabulate the difference between the true Hilbert series and (30).
## Appendix A The signed incidence matrix has full rank when the difference in cardinality is even
We prove a “signed version” of the well-known theorem that the incidence matrix of $`r`$-subsets of $`[n]=\{1,\mathrm{},n\}`$ into $`d+r`$-subsets have full rank. Our proof is a modification of the one by Wilson \[Wil73\].
To begin, let us define the “signs” involved.
###### Definition A.1.
Let $`[n]=\{1,\mathrm{},n\}`$, and let $`𝒞`$ and $``$ be two subsets of $`[n]`$, with
$$\begin{array}{cc}\hfill 𝒞& =\{t_1,\mathrm{},t_a\},t_1<\mathrm{}<t_a\hfill \\ \hfill & =\{k_1,\mathrm{},k_b\},k_1<\mathrm{}<k_b\hfill \end{array}$$
Then define $`\sigma (𝒞,)`$ to be zero if $`𝒞`$, and otherwise the sign of the permutation which sorts $`[𝒞,𝒞]`$ in ascending order. In other words, if $`𝒞`$ then $`\sigma (𝒞,)`$ is the sign of the uniquely determined permutation $`\gamma `$ such that
$$\begin{array}{cc}\hfill t_{\gamma (i)}& =k_i,1ia\hfill \\ \hfill k_{\gamma (j)}& =k_{a+j},1jb\hfill \end{array}$$
###### Definition A.2.
Let $`[n]=\{1,\mathrm{},n\}`$, and let $`A`$, $`B`$ be two subsets of $`[n]`$, of cardinality $`a`$ and $`b`$, with $`0a<b`$. For $`ar<b`$, we define
$$s_r(A,B,n)=\underset{\begin{array}{c}𝒞\left(\genfrac{}{}{0pt}{}{[n]}{r}\right)\\ A𝒞B\end{array}}{}\sigma (𝒞,B)$$
(31)
For $`0dn`$, we define
$$s_{d,n}=\underset{R\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)}{}\sigma (R,[n])=s_d(\mathrm{},[n],n)$$
(32)
###### Lemma A.3.
With the notations of Definition A.2, put $`d=br`$. We have that
$$s_r(A,B,n)=\{\begin{array}{cc}0\hfill & AB\hfill \\ (1)^ds_{d,ba}\hfill & AB\hfill \end{array}$$
(33)
###### Proof.
Put $`d=br`$. If $`AB`$ then clearly $`s_r(A,B,n)=0`$. Suppose that $`AB`$. Then
$$s_r(A,B,n)=\underset{\begin{array}{c}𝒞\left(\genfrac{}{}{0pt}{}{[n]}{r}\right)\\ A𝒞B\end{array}}{}\sigma (𝒞,B)=\underset{\begin{array}{c}𝒞\left(\genfrac{}{}{0pt}{}{B}{r}\right)\\ A𝒞\end{array}}{}\sigma (𝒞,B),$$
so the sum is independent of $`n`$. Furthermore, we can write $`A𝒞\left(\genfrac{}{}{0pt}{}{B}{r}\right)`$ as a disjoint union $`𝒞=A(𝒞A)`$, hence the sum can be written
$$\underset{S\left(\genfrac{}{}{0pt}{}{BA}{ra}\right)}{}\sigma (SA,B)=\underset{S\left(\genfrac{}{}{0pt}{}{BA}{ra}\right)}{}\sigma (S,BA).$$
Now, since $`S`$ has cardinality $`ra`$, the set $`(BA)S`$ has cardinality $`ba(ra)=br=d`$, so the permutation which transforms $`[S,BA]`$ to $`[BA,S]`$ has cardinality $`(1)^d`$. Hence, by substituting $`R=(BA)S`$, we get that the sum is equal to
$$\begin{array}{c}(1)^d\underset{S\left(\genfrac{}{}{0pt}{}{BA}{va}\right)}{}\sigma ((BA)S,BA)=(1)^d\underset{R\left(\genfrac{}{}{0pt}{}{BA}{d}\right)}{}\sigma (R,BA)\hfill \\ \hfill =(1)^d\underset{R\left(\genfrac{}{}{0pt}{}{[ba]}{d}\right)}{}\sigma (R,[ba]),\end{array}$$
which is the desired result. ∎
###### Lemma A.4.
Suppose that $`0<dn`$, and that $`d`$ is even. Then $`s_{d,n}>0`$.
###### Proof.
The lemma is trivially true for $`d=n`$. If $`d=2`$, we note that $`\sigma (\{v,v+1\},[n])=1`$ for $`1v<n`$, since the permutation transforming $`[v,v+1,1,2,\mathrm{},v1,v+2,v+3,\mathrm{},n]`$ to $`[1,2,\mathrm{},n]`$ is even. Furthermore, the signs of $`\sigma (\{v,v+\mathrm{}\},[n])`$ alternate in sign as $`\mathrm{}`$ goes from $`1`$ to $`nv`$. Thus, for a fixed $`v`$, there are either as many positive as negative $`\sigma (\{v,v+\mathrm{}\},[n])`$, or 1 more positive than negative, depending on the parity of $`nv`$. By summing over all $`v`$, we conclude that there are always strictly more positive than negative signs.
Now suppose that we have shown that $`s_{2k^{},n^{}}>0`$ for all $`k^{},n^{}`$ such that $`k^{}<k`$. We want to show that that $`s_{2k,n}>0`$. We have that
$$s_{2k,n}=\underset{R\left(\genfrac{}{}{0pt}{}{[n]}{2k}\right)}{}\sigma (R,[n]),$$
and writing $`R`$ as a disjoint union of its first two element, and the remaining elements, this becomes
$$\begin{array}{c}\underset{1k<\mathrm{}n2}{}\underset{R_2\left(\genfrac{}{}{0pt}{}{\{\mathrm{}+1,\mathrm{}+2,\mathrm{},n\}}{2k2}\right)}{}\sigma (\{k,\mathrm{}\}R_2,[n])\hfill \\ \hfill =\underset{1k<\mathrm{}n2}{}\underset{R_2\left(\genfrac{}{}{0pt}{}{\{\mathrm{}+1,\mathrm{}+2,\mathrm{},n\}}{2k2}\right)}{}\sigma (R_2,\{\mathrm{}+1,\mathrm{}+2,\mathrm{},n\})=\underset{1k<\mathrm{}n2}{}s_{2k2,n\mathrm{}}>0.\end{array}$$
Next, we define the signed incidence matrix.
###### Definition A.5.
Let $`0<a<bn`$ be integers. Then $`M_{a,b,n}`$ is the $`\left(\genfrac{}{}{0pt}{}{n}{b}\right)\times \left(\genfrac{}{}{0pt}{}{n}{a}\right)`$ matrix where the rows are indexed by $`b`$-subsets of $`[n]`$, the columns by $`a`$-subsets of $`[n]`$, and where the entry in row $`B`$, column $`A`$ is $`\sigma (A,B)`$.
###### Theorem A.6.
Let $`0<a<bn`$ be integers. If $`d=ba`$ is even, then $`M_{a,b,n}`$ has full rank.
###### Proof.
Denote the row indexed by $`\left(\genfrac{}{}{0pt}{}{[n]}{b}\right)`$ by $`\tau _{}`$, then $`\tau _{}`$ can be regarded as an element in $`V_a([n])`$, the free $``$-vector space on the $`a`$-subsets of $`[n]`$. If we denote the basis element corresponding to a $`a`$-subset $`𝒞`$ by $`ϵ_𝒞`$, then
$$\tau _{}=\underset{𝒞\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)}{}\sigma (𝒞,)ϵ_𝒞.$$
The number of rows in $`M_{a,b,n}`$ is $`\left(\genfrac{}{}{0pt}{}{n}{b}\right)`$, and the number of columns is $`\left(\genfrac{}{}{0pt}{}{n}{a}\right)`$. There are less rows than columns if $`a+b>n`$, as many rows as columns if $`a+b=n`$, and more rows than columns if $`a+b<n`$.
1. If $`\mathrm{a}+\mathrm{b}\mathrm{n}`$, we must prove that the rows are linearly independent. Suppose that there is a linear relation among the $`\tau _{}`$’s, so that
$$\underset{\left(\genfrac{}{}{0pt}{}{[n]}{b}\right)}{}a_{}\tau _{}=0$$
(34)
for some numbers $`a_{}`$. We shall prove that all $`a_{}=0`$.
Choose an $`I\left(\genfrac{}{}{0pt}{}{[n]}{i}\right)`$, $`0ia`$, and define a linear functional $`H_I:V_a([n])`$ by
$$f_I(ϵ_𝒞)=\{\begin{array}{cc}1\hfill & I𝒞\hfill \\ 0\hfill & I𝒞\hfill \end{array}$$
(35)
Then if $`\left(\genfrac{}{}{0pt}{}{[n]}{b}\right)`$ we have that
$$\begin{array}{c}f_I(\tau _{})=f_I\left(\underset{𝒞\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)}{}\sigma (𝒞,)ϵ_𝒞\right)=\underset{𝒞\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)}{}\sigma (𝒞,)f_I(ϵ_𝒞)\hfill \\ \hfill =\underset{I𝒞}{}\sigma (𝒞,)=s_a(I,,n)=\{\begin{array}{cc}s_{d,bi}\hfill & I\hfill \\ 0\hfill & I\hfill \end{array}\end{array}$$
(36)
The last step follows from Lemma A.3. Applying $`f_I`$ to (34) we get that
$$\begin{array}{c}0=f_I\left(\underset{\left(\genfrac{}{}{0pt}{}{[n]}{b}\right)}{}a_{}\tau _{}\right)=\underset{\left(\genfrac{}{}{0pt}{}{[n]}{b}\right)}{}a_{}f_I(\tau _{})\hfill \\ \hfill =\underset{\left(\genfrac{}{}{0pt}{}{[n]}{b}\right)}{}a_{}s_a(I,)=s_{d,bi}\underset{\begin{array}{c}\left(\genfrac{}{}{0pt}{}{[n]}{b}\right)I\end{array}}{}a_{}\end{array}$$
(37)
Since Lemma A.4 tells us that $`s_{d,bi}0`$, we conclude that
$$\underset{I}{}a_{}=0$$
(38)
Now, for any $`J[n]`$ we have, by exclusion-inclusion, that
$$\underset{J=\mathrm{}}{}a_{}=\underset{IJ}{}(1)^{|I|}\underset{I}{}a_{}$$
(39)
Fix $`_0\left(\genfrac{}{}{0pt}{}{[n]}{b}\right)`$ and put $`J_0=[n]_0`$. Since $`|J_0|=nba`$ we have, using (38) that
$$a__0=\underset{J_0=\mathrm{}}{}a_{}=\underset{IJ_0}{}(1)^{|I|}\underset{I}{}a_{}=0$$
(40)
Since $`a__0`$ was arbitrary, all $`a_{}`$ are zero. This shows that the $`\tau _{}`$ are linearly independent.
2. If $`\mathrm{n}=\mathrm{a}+\mathrm{b}`$, then $`M`$ is a square matrix. By the previous case, the vectors $`\tau _{}`$ are linearly independent, but since there are $`\left(\genfrac{}{}{0pt}{}{n}{a}\right)=\left(\genfrac{}{}{0pt}{}{n}{b}\right)`$ such vectors, they form a basis of $`V_a([n])`$; in particular, they span this vector space.
3. Finally, let us consider the remaining case $`\mathrm{n}>\mathrm{a}+\mathrm{b}`$, so that there are more rows than columns. We must prove that the rows span $`V_a([n])`$. We prove this by induction over $`nab`$. The case $`nab=0`$ is already proved, and forms the basis of the induction. We assume $`a,b`$ fixed, and that the assertion has been proved for all $`a+bn^{}<n`$.
Let $`\mathrm{\Gamma }\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)`$ be arbitrary. If we can express $`\alpha =ϵ_\mathrm{\Gamma }`$ as a linear combination of the $`\tau _{}`$’s, we are done. To this end, put
$$\alpha ^{}=\underset{\begin{array}{c}S\left(\genfrac{}{}{0pt}{}{[n1]}{a1}\right)\\ S\{n\}=\mathrm{\Gamma }\end{array}}{}ϵ_SV_{a1}([n1])$$
(41)
Since $`a1+b<n1`$, it follows by induction that there are scalars $`\left\{d_J\text{ }J\left(\genfrac{}{}{0pt}{}{[n1]}{a1}\right)\right\}`$ such that
$$\alpha ^{}=\underset{J\left(\genfrac{}{}{0pt}{}{[n1]}{a1}\right)}{}d_J\tau _J^{},\tau _J^{}=\underset{\begin{array}{c}S\left(\genfrac{}{}{0pt}{}{[n1]}{a1}\right)\\ SJ\end{array}}{}ϵ_S$$
(42)
For $`\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)`$, $`n`$, put $`c_{}^{}=d_{}\{n\}`$. Define
$$\alpha _0=\underset{\begin{array}{c}\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)\\ n\end{array}}{}c_{}^{}\tau _{}V_a([n])$$
(43)
If we write
$$\alpha _0=\underset{𝒞\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)}{}a_𝒞^{}ϵ_𝒞$$
we have that for $`𝒞\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)`$, $`n𝒞`$, that
$$a_𝒞^{}=\{\begin{array}{cc}1\hfill & 𝒞=\mathrm{\Gamma }\hfill \\ 0\hfill & 𝒞\mathrm{\Gamma }\hfill \end{array}$$
which implies that
$$\alpha _0=\{\begin{array}{cc}\alpha \hfill & n\mathrm{\Gamma }\hfill \\ 0\hfill & n\mathrm{\Gamma }\hfill \end{array}$$
In either case, $`\alpha \alpha _0`$ has coordinate 0 in component $`\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)`$, unless $`n`$. Hence, $`\alpha \alpha _0`$ may be regarded as a vector in $`V_a([n1])`$. By the induction hypothesis, there exist $`c_{}^{\prime \prime }`$ such that
$$\alpha \alpha _0=\underset{\left(\genfrac{}{}{0pt}{}{[n1]}{a}\right)}{}c_{}^{\prime \prime }\tau _{}$$
(44)
Defining
$$c_{}=\{\begin{array}{cc}c_{}^{}\hfill & n\hfill \\ c_{}^{\prime \prime }\hfill & n\hfill \end{array}$$
we get that
$$\alpha =\alpha _0+(\alpha \alpha _0)=\underset{\begin{array}{c}\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)\\ n\end{array}}{}c_{}^{}\tau _{}+\underset{\begin{array}{c}\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)\\ n\end{array}}{}c_{}^{\prime \prime }\tau _{}=\underset{\left(\genfrac{}{}{0pt}{}{[n]}{a}\right)}{}c_{}\tau _{}$$
## Appendix B Calculations
The computer calculations were done on the computers of the UMS Medicis, École Polytechnique, and on the computers at the Department of Mathematics, Stockholm University. We have used the programme Macaulay 2 \[GS\] to calculate Hilbert series and minimal free resolutions. To save time and memory, the calculations were performed in characteristic 31991. The holes in the tables show that there are limits to what we could calculate, even on a machine with 2 GB of memory.
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# 1 Introduction
## 1 Introduction
We may need to go a long way before we understand how the confinement mechanism works within nonabelian gauge theory. A lot of physicists now consider the Abelian projection as the best way to understand the confinement mechanism.
After Abelian projection one should consider the element of some Abelian subgroup of the gauge group instead of the full gauge group element. Thus, the theory becomes Abelian and one can look at the picture of the confinement phenomenon in a simpler way.
Abelian projections differ from each other by the choice of the subgroup of the gauge group and the projection method. The closeness of a given Abelian projection to the solution of the confinement problem is measured as follows. Suppose that the link gauge group elements $`g_{\mathrm{link}}G`$ are projected onto the elements of the given Abelian subgroup $`e_{\mathrm{link}}EG`$. We consider
$$Z_C=\mathrm{Tr}\underset{\mathrm{link}\mathrm{C}}{}e_{\mathrm{link}}$$
(1)
instead of the Wilson loop and extract the potential from $`Z_C`$, (the projected potential). If that potential is close to the original confining potential at sufficiently large distances one can say that the projection is suitable for the investigation of the confinement mechanism.
Dealing with the $`SU(2)`$ gauge group, the Cartan subgroup $`U(1)`$ and the center subgroup $`Z_2`$ have been considered. The most popular projections are the Maximal Abelian and the Maximal Center projections. Those projections are achieved by minimization with respect to the gauge transformations of the distance between the given configuration of the link gauge field and the Cartan (center) subgroup of $`SU(2)`$. In both cases the potentials are very close to the $`SU(2)`$ confining potential but unfortunately do not exactly coincide with it. Greensite et al. have argued that only central charges $`q=\pm 1`$ are confined in nonabelian gauge theories.
From the work of Bornyakov et al. it is known that gauge fixing procedures suffer from a gauge copies problem. Several gauges have been used so far in practical computations: the direct and indirect center gauges and the Laplacian center gauge . Only the first two suffers from the occurrence of gauge copies, the last one is free from this difficulty. According to Refs. this problem disappears for large lattices and when the number of Gribov copies considered is increased. Of course, this means that the computational effort must also be increased.
In this work we propose a new center projection which is not connected with partial gauge fixing. Thus, all the objects existing within the projected theory have a gauge invariant nature. We call this procedure the Simple Center Projection (SCP). Based on numerical simulations we make the conjecture that the projected potential coincides with the $`SU(2)`$ potential up to the mass renormalization term at sufficiently large distances.
Within the SCP we can construct the center vortices and the center monopoles (also known as nexuses, see Refs. and ). The properties of the SCP monopoles found in our investigations lead us to believe that those monopoles are the objects which play the role of the Cooper pairs in the dual superconductor.
## 2 Simple Center Projection
We consider $`SU(2)`$ gluodynamics with the Wilson action
$$S(U)=\beta \underset{\mathrm{plaq}}{}(11/2\mathrm{Tr}U_{\mathrm{plaq}}).$$
(2)
The sum runs over all the plaquettes of the lattice. The plaquette action $`U_{\mathrm{plaq}}`$ is defined in the standard way.
First we consider the plaquette variable
$`z_{\mathrm{plaq}}=1,`$ $`\mathrm{if}`$ $`\mathrm{Tr}U_{\mathrm{plaq}}<0,`$
$`z_{\mathrm{plaq}}=0,`$ $`\mathrm{if}`$ $`\mathrm{Tr}U_{\mathrm{plaq}}>0.`$ (3)
We can represent $`z`$ as the sum of a closed form $`dN`$ <sup>1</sup><sup>1</sup>1We use the formulation of differential forms on the lattice, as described for instance in Ref. . for $`N\{0,1\}`$ and the form $`2m+q`$. Here $`N=N_{\mathrm{link}}`$, $`q\{0,1\}`$, and $`m𝖹𝖹`$.
$$z=dN+2m+q.$$
(4)
The physical variables depending upon $`z`$ could be expressed through
$$\mathrm{sign}\mathrm{Tr}U_{\mathrm{plaq}}=\mathrm{cos}(\pi (dN+q)).$$
(5)
We shall say that $`N_{\mathrm{link}}`$ is the center projected link variable. There are many different ways to make this projection. The Maximal Center Projection (MCP) uses the gauge ambiguity to make all link matrices $`U`$ as close as possible to $`e^{i\pi N}`$. Thus the 1-form $`N`$ is fixed for every gauge configuration.
There exist several ways to achieve this . One way, the direct way, is to minimize the quantity
$$R=\underset{x}{}\underset{\mu }{}\mathrm{Tr}[U_\mu (x)]\mathrm{Tr}[U_\mu ^{}(x)].$$
(6)
In the indirect way one minimizes the quantity
$$R^{}=\underset{x}{}\underset{\mu }{}\mathrm{Tr}[U_\mu (x)\sigma _3U_\mu ^{}(x)\sigma _3],$$
(7)
and extracts from $`U_\mu (x)`$ the diagonal part $`A_\mu =\mathrm{exp}[i\theta _\mu (x)\sigma _3]`$ thus fixing the gauge. Finally the remnant Abelian symmetry is used to bring $`A_\mu `$ as close to an element of $`Z_2`$ as possible by maximizing
$$R^{\prime \prime }=\underset{x}{}\underset{\mu }{}\mathrm{cos}^2\theta _\mu (x).$$
(8)
It is clear that both procedures are complicated and time consuming, as the number of variables involved in the case of e.g. $`SU(2)`$ is three. Whatever method is used, the 1-form $`N`$ is fixed for every gauge configuration.
Now we choose a simpler and more natural procedure. Imagine the surface $`\mathrm{\Sigma }`$ formed by the plaquettes dual to the ”negative” plaquettes (for which $`z_{\mathrm{plaq}}=1`$). This surface has a boundary. We enlarge the surface by adding a surface $`\mathrm{\Sigma }_{\mathrm{add}}`$ so that:
1. the resulting surface $`\mathrm{\Sigma }^1=\mathrm{\Sigma }+\mathrm{\Sigma }_{\mathrm{add}}`$ will be closed;
2. when we eliminate from the surface $`\mathrm{\Sigma }_{\mathrm{add}}`$ the plaquettes carrying even numbers $`z_{\mathrm{plaq}}=\mathrm{},4,2,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}4},\mathrm{}`$, the area of the remnant surface will be minimal for the given boundary.
So $`\mathrm{\Sigma }^1`$ can be represented by the closed form $`dN`$ on the original lattice for the integer link variable $`N\{0,1\}`$. And $`N`$ is the required center projected link variable.
Numerically this procedure lookes as follows. For the given variable $`z`$ we should choose such a $`Z_2`$ variable $`N`$ that $`[dN]\mathrm{mod}\mathrm{\hspace{0.17em}2}`$ is as close as possible to $`z`$. It means that we minimize the functional
$$Q=\underset{plaq}{}|(zdN)\mathrm{mod}\mathrm{\hspace{0.17em}2}|$$
(9)
with respect to $`N`$.
The procedure works in the following way. We consider a given link $`L`$ and the sum over plaquettes
$$Q_{link}=\underset{L\mathrm{plaq}}{}|(zdN)\mathrm{mod}\mathrm{\hspace{0.17em}2}|$$
(10)
We minimize this sum with respect to the one link value $`N`$. All links are treated in this way and the procedure is iterated until a global minimum is found.
We call this unique and simple procedure the Simple Center Projection. Our projection method also finds local minima, Gribov copies, but as the procedure is much simpler and faster that the methods use until now, we can afford to repeat our calculations to include several Gribov copies.
The physical meaning of the projected variables becomes clear after considering the continuum limit. Naively, the considered surfaces disappear as the field strength on them tends to infinity. Nevertheless this fact should be investigated more carefully. In any case for resonable sizes of lattices finite volume effects occur and the scaling window ends at some value of $`\beta `$. Thus the direct drop into the continuum limit is impossible and within the scaling window our surfaces $`\mathrm{\Sigma }`$ carry large but not infinite field strength. The results of the next section give us the reason to believe that these surfaces are those which play the crucial role in the confinement mechanism.
Let us also mention that it is reasonable to consider the analogous construction for which we change the plaquettes into loops of sizes $`2\times 2`$, $`3\times 3,\mathrm{}`$. These extended projections solve the problem of the positive plaquette model, for which our ”$`1\times 1`$$`\mathrm{\Sigma }`$ is absent.
## 3 Numerical results
In the calculations we report on here, we used as our standard lattice one with dimension $`24^4`$ and for investigations at finite temperature one with size $`24^3\times 4`$. For reasons of comparison we have occasionally used smaller lattices, of sizes $`12^3\times 4`$ and $`16^3\times 4`$. Some results were checked on the larger lattice $`32^3\times 4`$. We have mentioned the gauge copy problem. We checked that for lattices with linear dimensions $`L24`$, which we have used, 15 copies are sufficient and so we used everywhere 15 Gribov copies.
### 3.1 Center Dominance
We consider the following definition of the projected Wilson loop
$$W_C^{\mathrm{SCP}}=\frac{1}{8}Z_C(3\pi /4)^{𝒫(C)/4},$$
(11)
Where $`𝒫(C)`$ is the perimeter of the loop $`C`$.
It may be interesting for the reader that this expression is very similar to the first term in the character expansion from . One can get the above expression from that formula substituting the factor $`1/8`$ instead of $`1/4`$ and the power $`𝒫(C)/4`$ instead of $`𝒫`$. However, as the authors of Refs. derive their expression for lattices without gauge fixing, they can derive their results with local operators, using the characters of the gauge group to obtain dominance of the fundamental representation. In our case, we perform gauge fixing, which results in a nonlocal operator. Consequently, we cannot use the same derivation as Refs. . As we have not been able to find a rigorous derivation of Eq. (11), one must consider our results up till now as empirical ones. It should be stressed that these results were checked for Wilson loops of sizes up till $`6\times 6`$. It might be useful to check whether Eq. (11) continues to give good results if larger Wilson loops are considered and the statistics is improved.
The reader can recognize from Fig. 1 that $`W_C`$ and $`W_C^{\mathrm{SCP}}`$ exactely coincide with each other for large enough sizes of the loop (we represent here $`\mathrm{log}W_C/\mathrm{Area}(C)`$ and $`\mathrm{log}W_C^{\mathrm{SCP}}/\mathrm{Area}(C)`$ as a function of the loop area for the values $`\beta `$ = 2.3 and 2.4).
$`W_C^{\mathrm{SCP}}`$ differs from $`Z_C`$ by the perimeter factor. Thus we understand that the projected potential (extracted from $`Z_C`$) differs from the full potential (extracted from $`W_C`$) only by the mass renormalization term at low energies. Other terms (including the confining linear term) are the same. It demonstrates the Center Dominance.
### 3.2 The center vortices
We construct the closed two-dimensional center vortices as in
$$\sigma =^{}dN.$$
(12)
First, we can express $`Z_C`$ as
$$Z_C=\mathrm{exp}(i\pi LL(\sigma ,C))$$
(13)
where $`LL`$ is the linking number .
Thus, according to the previous subsection the Aharonov - Bohm interaction between the center vortices and the charged particle leads to the confinement of the fundamental charge .
Also we investigate other properties of the center vortices for the finite temperature theory (the nonsymmetric lattice $`24^3\times 4`$).
The density of the vortices $`\rho `$ is shown in Fig. 2.
The fractal dimension, defined as $`D=1+2A/L`$, where $`A`$ is the number of plaquettes and $`L`$ is the number of links of the vortices, is shown in Fig. 3. A line $`L`$ is counted as belonging to the vortex if at least one of the faces of a cube dual to that link contains a plaquette with charge 1.
### 3.3 The center monopoles (nexuses)
Following we construct the center monopoles (nexuses)
$$j=\frac{1}{2}^{}d[dN]\mathrm{mod}\mathrm{\hspace{0.17em}2}.$$
(14)
We show the density of the nexuses as a function of $`\beta `$ in Fig. 2. It is important to make sure that the monopole lines are closed. It follows from the equation $`\delta j=\frac{1}{2}{}_{}{}^{}d{}_{}{}^{}{}_{}{}^{}d[dN]\mathrm{mod}\mathrm{\hspace{0.33em}2}=0`$.
We investigated the percolation properties of $`j`$ and considered the probability of two points to be connected by a monopole worldline on a constant-time hypersurface. The dependence of that probability upon $`\beta `$ for a nonsymmetric lattice is shown in Fig. 4. The percolation probability is dependent on the lattice size. This is obvious for very small lattices, but it is seen to persists at larger sizes.
Here we would like to remark on the role the monopole condensate plays as an order parameter. Ivanenko et al. have shown for abelian monopoles, obtained after Maximal Abelian projection, that the monopole condensate can be used as an order parameter: it vanishes in the deconfined phase, while it takes a finite value in the confining phase. The authors of Ref. made the interesting observation that the behavior of the condensate near the point of the phase transition depends on the “thickness” of the monopole lines. It was demonstrated in Ref. that Abelian monopoles are strongly connected with vortices. For vortex lines with a thickness of one lattice spacing, the condensate vanishes smoothly near the critical point, whereas for lines with a thickness of two lattice units, the variation of the condensate near the critical point is very steep.
We see that the monopoles are condensed in the confinement phase and not condensed in the deconfinement phase, but the phase transition is rather smooth as is the case for thin Abelian monopoles. If we would investigate, along the lines of Ref. nexuses of larger sizes like $`2^3`$, $`3^3`$ and so forth, we think that we shall obtain a better sensitivity of the condensates $`C^{(2)}`$, $`C^{(3)}`$ etc. to the phase transition. (Here we use the obvious notation $`C^{(n)}`$ for the condensate of vortices of size $`n^3`$.
We expect the center monopoles to be the monopoles that are present in the dual superconductor picture of confinement. As the phase transition is not very pronounced for the thin vortices we considered here, our results may be taken as a “proof of principle”. They must be substantiated by considering thick vortices which are supposed to be more strongly connected to confinement .
The analytical connection of the monopole condensation and the formation of the dual superconductor was considered in for the case of $`SU(3)`$ symmetry. Of course, the results of that paper obtain also for the $`SU(2)`$ theory. It follows from that in the case both condensation of nexuses and center dominance occur, the picture of the dual superconductor in which the nexuses play the role of Cooper pairs and the quarks play the role of the monopoles becomes clear.
In particular, one can rewrite the fundamental Wilson loop as follows
$`<W(C)>={\displaystyle _\pi ^\pi }DH{\displaystyle }D_{\mathrm{\Phi }C}\mathrm{\Phi }\mathrm{exp}(Q(dH+\pi ^{}A[C])`$
$`{\displaystyle \underset{xy}{}}\mathrm{\Phi }_xe^{2iH_{xy}}\mathrm{\Phi }_y^+V(|\mathrm{\Phi }|))`$ (15)
Here $`H`$ is the electromagnetic field, $`A[C]`$ is the area of the surface spanned on the quark loop, $`\mathrm{\Phi }`$ is the nexus field, $`Q`$ is nonlocal effective action, $`V`$ is an infinitely deep potential, supporting the infinite value of the nexus condensate.
Thus we have indeed the nonlocal relativistic superconductor theory, in which the nexuses are the Cooper pairs and the quarks are the monopoles. The condensation of nexuses gives rise to the formation of the quark - antiquark string appearing as the Abrikosov vortex.
## 4 Conclusions
In this work we construct the gauge invariant Center projection and show that center dominance takes place. We investigate the properties of the topological defects in the center projected theory, which are shown to be closely connected to the confinement picture. Particulary it occurs that the center monopoles from this projection are good candidates for the Cooper pairs in the dual superconductor. The simple numerical nature of the Simple Center Projection, the exact Center Dominance and the properties of the topological defects existing in the center projected theory give us the reason to propose SCP as the basic abelian projection for the considering of the confinement picture.
## Acknowledgments
We are grateful to V.G. Bornyakov, M.N. Chernodub, and M.I. Polikarpov for useful discussions. A.I.V. and M.A.Z. kindly acknowledge the hospitality of the Department of Physics and Astronomy of the Vrije Universiteit, where part of this work was done. This work was partly supported by the grants RFBR 99-01-01230, INTAS 96-370.
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# BCS superfluidity in ultracold gases with unequal atomic populations
## Abstract
We consider the existence of a BCS superfluid phase in <sup>6</sup>Li due to the pairing of two hyperfine states with unequal number of atoms. We show that the domain of existence for this phase will be increased to a very large extent in the vicinity of the threshold for collapse. This is due to the presence of a new phase with anisotropic order parameter. This phase is induced by the anisotropic part of the scattering, linked to the indirect interaction between atoms through density fluctuation exchange.
Much progress has been made recently toward the experimental investigation of degenerate ultracold Fermi gases. Indeed the degenerate regime has already been reached in <sup>40</sup>K by sympathetic cooling of two hyperfine states whereas the possibility of cooling a mixture of <sup>6</sup>Li and <sup>7</sup>Li atoms in order to reach this range has been demonstrated . In addition to a number of very interesting physical effects linked to phase space restriction one of the most fascinating prospect for experiment is the possibility to reach the onset of Cooper pairs condensation and observe the resulting BCS superfluid, the equivalent of Bose Einstein condensates for Fermi gases. As we will see this phase is likely to present quite unusual features.
In addition to the problems raised by the use of evaporative cooling for Fermi gases a possible stumbling block in the path to the BCS superfluid is the need to achieve a near equality between the number of atoms in the two hyperfine states assumed to form pairs. Indeed in the various schemes under investigation there is no fast relaxation mechanism which equalizes these two populations and the number of atoms in each state will be essentially conserved. Therefore in these ultracold atomic gases, one will necessarily have to deal with situations where the two populations of particles assumed to form pairs are not equal. Now the difference in chemical potential between these two populations acts as an effective field which tends to break pairs and the superfluid phase is destroyed when this difference is of order of the critical temperature. Hence this pair breaking effect becomes more important a problem if the critical temperature is low and, since there is presently some uncertainty in the value of $`T_c`$, it is clearly of interest to investigate this question in detail. Moreover, independently of this possible experimental problem, this difference in chemical potential is an additional control parameter in the system quite interesting to play with and it is worth studying its effect, all the more since it is not trivial.
This problem of pairing with two unequal populations has already been considered a long time ago in the context of superconductivity. Here this is an applied magnetic field which tends to make the two spin populations unequal. Usually the critical field is limited by orbital effects. However the question of the limitation of the superconducting phase, when these orbital effects are small, has been considered very early and it was pointed out by Clogston and Chandrasekhar that the standard BCS phase could, at most, resist to a difference in chemical potential between the two populations of the order of the critical temperature (the so-called paramagnetic limit). Not long after, Fulde and Ferrell (FF), and independently Larkin and Ovchinnikov (LO), showed that pairing could somewhat adjust to the difference in chemical potential, instead of just resisting, by letting the pairs have a common momentum $`𝐊`$ instead of having it equal to zero as for a standard superconductor. However the effect was actually found rather small. Indeed at zero temperature the standard BCS phase goes to the normal state by a first order transition when the chemical potential difference $`2\overline{\mu }`$ is equal to $`\sqrt{2}\mathrm{\Delta }_0`$ , where $`\mathrm{\Delta }_0=1.76T_c`$ is the zero temperature gap for equal populations. The FFLO phase goes to the normal state by a second order transition for $`\overline{\mu }=0.754\mathrm{\Delta }_0`$, which is not much beyond. Actually there is to date no undisputed observation of this phase in standard superconductors , most likely because it is quite sensitive to impurities. It is the purpose of the present paper to show that pairing can adjust even better than in the FFLO phase, and that this new phase can in some instances lead to quite a strong increase of the existence domain for the superfluid phase.
One can view the appearance of the FFLO phase in the following way. When the chemical potentials for the two spin populations are different, it becomes more costly in terms of kinetic energy to form $`(𝐤,𝐤)`$ pairs in the standard BCS way because one can not pick the two particles very near the Fermi surfaces, since these surfaces do not match due to their size difference. In the FFLO phase this problem is remedied by taking a nonzero total momentum which amounts to shifting, in momentum space, one of the Fermi surface with respect to the other so that they almost match on some region. But naturally matching gets worst on the opposite side. Nevertheless the total balance is barely favourable. Now one can think to improve this situation by making pairing stronger on the side where there is energy gain and weaker where there is energy loss. This means one looks for anisotropic pairing to find a lower energy ground state. Naturally this can not work if scattering is isotropic, as it will be essentially the case for a weakly interacting ultracold gas since p-wave scattering is negligible, and this will be indeed the situation in the dilute regime where the coupling constant $`\lambda =2k_F|a|/\pi `$ is small. However, even if the bare scattering is isotropic, the renormalized interaction is not because of the existence of the Fermi surface. This effect is at the basis of the Kohn and Luttinger paper, where they showed that, even with a repulsive interaction, Cooper pairs would necessarily form in high angular momentum. For our purpose the case of <sup>6</sup>Li is of particular interest since the high density regime is bounded by an instability occuring for $`\lambda 1`$ . In the vicinity of this instability the effective interaction will be strongly anisotropic, leading to a marked increase in the existence domain of the superfluid phase, compared to the FFLO phase. Note that most experiments are likely to be done in this range since it corresponds to higher critical temperature. We remark also that a similar phase with anisotropic order parameter should exist in clean superconductors with Pauli limited upper critical field and anisotropic interaction.
Before going into the effect of anisotropic pairing, it is of interest to show that this is the best possible choice for pairing with unequal particle number. Indeed assume that we look for the most general kind of pairing, where $`𝐤`$ is paired with $`f(𝐤)`$. On one hand we want pairing to be stable against scattering, which means that scattering must send this pair into another pair $`𝐤^{}`$ , $`f(𝐤^{})`$. On the other hand scattering conserves momentum which leads, for any $`𝐤`$ and $`𝐤^{}`$, to $`𝐤+f(𝐤)=𝐤^{}+f(𝐤^{})=𝐊`$ where $`𝐊`$ is a constant. This shows that $`f(𝐤)=𝐊𝐤`$, that is pairs with a total momentum $`𝐊`$ is the most general solution. Note that this argument assumes translational invariance. This will not hold for a finite sample such as those obtained with trapped gases, which will induce surface effects. However these surface effects should be small if we do not want to have $`T_c`$ drastically reduced (more precisely we want the pair size to be small compared to the sample size). So it is reasonable to neglect size effects in a first step. Now the common momentum $`𝐊`$ produces a breaking of rotational symmetry, and the standard symmetry analysis leading to pairs with a given angular momentum (s-wave pairs in the case of <sup>6</sup>Li) is no longer valid. The general order parameter $`\mathrm{\Delta }_𝐤`$ will depend on the wavevector $`𝐤`$. However we still have rotational invariance around $`𝐊`$ and we can classify the solutions by their angular momentum $`m`$. We will assume that the most stable pairing corresponds to $`m=0`$, as it is likely to be so for a standard interaction. However other values of $`m`$, corresponding to a breaking of the rotational symmetry around $`𝐊`$, do not seem to be excluded from first principles and would be undoubtedly a very interesting situation.
We will explore now quantitatively the possibility offered by an anisotropic order parameter. However our purpose is more to demonstrate the importance of the effect than to perform an exact calculation. This last goal would require a perfect knowledge of the effective interaction, which is not available. We will consider specifically the case of <sup>6</sup>Li and we restrict ourselves to the determination of the T = 0 critical difference in chemical potential $`2\overline{\mu }\mu \mu `$ above which superfluidity disappears. Since we want to go in the high density regime, we need to have an expression for the effective interaction in this range. For this purpose we take the paramagnon model which we have already used for an evaluation of $`T_c`$ in the high density regime . Actually we will not retain the full interaction of this model because most of the terms lead to a moderate anisotropy and they would produce a more complex calculation without adding much effect. The essential contribution of those terms is already included in the value $`T_c^0`$ of the critical temperature for equal populations. We will only retain the explicit attractive part coming from density fluctuations which produces near the instability a strong contribution for low momentum transfer and in this way lead to a strongly anisotropic interaction. Moreover we will for simplicity perform a weak coupling calculation, omitting all the frequency dependence and taking the zero frequency value of the interaction. To be coherent with this weak coupling approach we will also omit self-energy effects. This leads to the total effective attractive interaction $`V(𝐤,𝐤^{})`$ given by :
$`N_fV(𝐤,𝐤^{})=\lambda +{\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda ^2\overline{\chi }_0(𝐪)}{1\lambda \overline{\chi }_0(𝐪)}}`$ (1)
with $`𝐪=𝐤𝐤^{}`$. Here $`N_f=mk_F/2\pi ^2`$ is the density of states at the Fermi surface for equal population and $`\overline{\chi }_0(𝐪)`$ is the reduced elementary bubble at zero frequency:
$`2\overline{\chi }_0(𝐪)=1+{\displaystyle \frac{1}{y}}(1{\displaystyle \frac{y^2}{4}})\mathrm{ln}{\displaystyle \frac{2+y}{2y}}`$ (2)
with $`y=q/k_F`$. The first term in Eq.(1) is the direct term and the last one is the indirect interaction due to density fluctuation exchange.
The difference in chemical potential $`2\overline{\mu }`$ has the effect of shifting the kinetic energies of the particle (measured from chemical potential) by $`\pm \overline{\mu }`$. For the pair propagator this is just equivalent to shift the frequency $`\omega `$ by $`\overline{\mu }`$. Similarly taking the momenta of the pair members to be $`\pm 𝐤+𝐊/2`$ instead of $`\pm 𝐤`$ produces a shift by $`\pm 𝐤.𝐊/2m`$, leading to an overall shift in frequency by $`\overline{\mu }_k=\overline{\mu }𝐤.𝐊/2m`$. At finite temperature this means replacing the Matsubara frequency $`\omega _n=\pi T(2n+1)`$ by $`\omega _ni\overline{\mu }_k`$. On the other hand we expect the modification of the effective interaction caused by $`\overline{\mu }`$ to be small, of order $`\overline{\mu }/E_F`$, and we can neglect it. This leads us finally to the following gap equation at finite temperature :
$`\mathrm{\Delta }_𝐤^{}={\displaystyle \frac{\lambda T}{N_f}}{\displaystyle \underset{n}{}}{\displaystyle \frac{d𝐤}{(2\pi )^3}}`$ (3)
$`{\displaystyle \frac{\mathrm{\Delta }_𝐤}{\xi _𝐤^2+\mathrm{\Delta }_𝐤^2+(\omega _ni\overline{\mu }_k)^2}}(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda \overline{\chi }_0(𝐪)}{1\lambda \overline{\chi }_0(𝐪)}})`$ (4)
where $`\xi _𝐤`$ is the kinetic energy measured from the Fermi surface for $`\overline{\mu }=0`$. Since we are looking for the critical $`\overline{\mu }`$ and expect a second ordre phase transition as for the FFLO phase, we let $`\mathrm{\Delta }_𝐤`$ go to zero in the denominator. Then, in order to get rid of the cut-off in the frequency summation, it is convenient to make use of $`\pi T(1/|\omega _n|)=\lambda ^1+\mathrm{ln}(T_c^0/T)`$, where $`T_c^0`$ is the critical temperature for $`\overline{\mu }=0`$, $`𝐊=0`$ and in the absence of the indirect term in the interaction. When we specialize to T $`0`$, the frequency summation can be performed by $`\pi T\mathrm{sgn}(\omega _\mathrm{n})[(1/(\omega _\mathrm{n}\dot{\mathrm{ı}}\overline{\mu }_\mathrm{k})1/\omega _\mathrm{n}]=\mathrm{ln}(0.88\mathrm{T}/|\overline{\mu }_\mathrm{k}|)`$. In this way we obtain the following equation for the critical difference in chemical potential $`\overline{\mu }`$ :
$`\mathrm{\Delta }_𝐤^{}={\displaystyle \frac{d\mathrm{\Omega }_k}{4\pi }\mathrm{\Delta }_𝐤(1+\lambda \mathrm{ln}\frac{\mathrm{\Delta }_0}{2|\overline{\mu }_k|})(1+\frac{1}{2}\frac{\lambda \overline{\chi }_0(𝐪)}{1\lambda \overline{\chi }_0(𝐪)})}`$ (5)
where $`\mathrm{\Delta }_0`$ is the zero temperature gap for $`\overline{\mu }=0`$, $`𝐊=0`$ and in the absence of the indirect term in the interaction. When the indirect interaction is not present, $`\mathrm{\Delta }_𝐤`$ is independent of $`𝐤`$. Then the angular integration is easily performed, leading to $`\mathrm{ln}(\mathrm{\Delta }_0/2\overline{\mu })=(1/2)f_0(x)`$, with $`x=Kk_F/(2m\overline{\mu })`$. Here $`f_0(x)=_1^1𝑑u\mathrm{ln}|1+xu|`$ which is negative and minimum for $`x1.200`$ ($`x=\mathrm{coth}x`$) leading to the standard FFLO result $`\overline{\mu }=0.754\mathrm{\Delta }_0`$ .
We consider now the effect of the indirect interaction. Since both $`𝐤`$ and $`𝐤^{}`$ are on the Fermi surface, $`q`$ goes from 0 to $`2k_F`$ and $`\overline{\chi }_0(𝐪)`$ goes from 1 to 1/2 . Although it is not a problem to use numerically the exact expression for $`\overline{\chi }_0(𝐪)`$, one can see that $`\overline{\chi }_0(𝐪)=1a+a\mathrm{cos}\theta `$, with $`a=0.25`$ and $`q=2k_F\mathrm{sin}(\theta /2)`$ is a quite good approximation. Indeed it gives properly the two limits $`q=0`$ and $`q=2k_F`$. Moreover it corresponds to keep only the first two terms in a Legendre polynomials expansion with slightly modified coefficients (the exact $`l=0`$ coefficient is $`(1+2\mathrm{ln}2)/3=0.755`$ and the exact coefficient of $`\mathrm{cos}\theta `$ is 0.232 , the higher order terms being fairly small). This approximation allows to perform the azimuthal integration analytically, and since our calculation is anyway a model calculation it is quite reasonable to make this simplification, despite its slight inaccuracy in the vicinity of $`\theta =0`$.
However, before proceeding to the solution of the resulting equation, it is interesting to consider first a slightly simplified version which allows to carry out the calculation completely explicitely. We merely replace the interaction term by the first two terms of its Legendre polynomial expansion $`V_0+V_1\mathrm{cos}\theta `$, with $`V_0=0.5+0.25(a\lambda )^1\mathrm{log}[1+2a\lambda /(1\lambda )]`$ and $`V_1=1.5(a\lambda )^1+0.75(1\lambda +a\lambda )(a\lambda )^2\mathrm{log}[1+2a\lambda /(1\lambda )]`$. In this case the solution of Eq.(4) has the form $`\mathrm{\Delta }_𝐤=1+\delta _1\mathrm{cos}\alpha `$ with $`\mathrm{cos}\alpha =\widehat{𝐤}.\widehat{𝐊}`$. Then $`L\mathrm{ln}(2\overline{\mu }/\mathrm{\Delta }_0)`$ is solution of the second order equation $`(L1/\lambda +1/\lambda V_0+0.5f_0)(L1/\lambda +3/\lambda V_1+1.5f_2)=0.75f_1^2`$. We have set $`f_n(x)=_1^1𝑑uu^n\mathrm{ln}|1+xu|`$. $`f_1(x)`$ and $`f_2(x)`$ are both found to be largest for $`x`$ in the range 1.1 - 1.2, and the largest $`L`$ is also found for the same range of $`x`$ for all values of $`\lambda `$.
Naturally the indirect interaction produces a trivial effect, namely the renormalization of the coupling constant $`\lambda `$ into $`\lambda V_0`$, producing a corresponding change of the critical temperature and of the gap for equal population, which goes from $`\mathrm{\Delta }_0`$ to $`\overline{\mathrm{\Delta }}_0=\mathrm{\Delta }_0\mathrm{exp}(1/\lambda 1/\lambda V_0)`$. This effect is simply found by looking at the solution $`L_0`$ for $`x=0`$, that is without FFLO phase, since in this case only the isotropic part of the interaction is relevant. In the present case this is given by $`L_0=1/\lambda 1/\lambda V_0`$. We are only interested in $`LL_0`$ which gives the increase in the superfluid domain due to the existence of our FFLO phase, in units of $`\overline{\mathrm{\Delta }}_0`$. In Fig.1 we have plotted as the dashed line $`\overline{\mu }/\overline{\mathrm{\Delta }}_0`$. We see that, roughly for $`\lambda <0.6`$, there is essentially no change with respect to the standard FFLO result. And indeed the gap remains essentially isotropic in this range. On the other hand, for $`\lambda >0.6`$, $`\overline{\mu }`$ increases rapidly and for $`\lambda =0.9`$ it is almost twice the standard FFLO result. It is also quite interesting to consider the anisotropy linked to $`\delta _1`$. As soon as $`\overline{\mu }`$ starts to grow with respect to the standard FFLO result, the order parameter gets a sizeable anisotropy. For $`\lambda =0.94`$ we obtain a node on the Fermi surface, and for larger $`\lambda `$ a change of sign over the Fermi surface with a nodal line.
Let us turn now to the results of the numerical solution of Eq.(4) . They are given as the full line on Fig.1 for $`\overline{\mu }/\overline{\mathrm{\Delta }}_0`$ as a function of $`\lambda `$ (as above we have taken into account the renormalization of $`\mathrm{\Delta }_0`$ into $`\overline{\mathrm{\Delta }}_0`$ by calculating $`LL_0`$). For small and intermediate $`\lambda `$ the result is essentially identical to the result of our simplified version. However when $`\lambda `$ approaches 1, $`\overline{\mu }`$ increases more rapidly. This is easy to understand by looking at Eq.(4). In this regime the effective interaction becomes very large for small $`𝐪`$ and the log term favors wavevectors nearly along $`𝐊`$ so it is better to have the gap function $`\mathrm{\Delta }_𝐤`$ peaked for wavevector along $`𝐊`$ (see Fig.2). This can be seen more explicitely by looking at the simplified form taken by this equation when $`\lambda 1`$. Taking into account the numerical evidence that in this limit the optimal $`\overline{\mu }`$ is obtained for $`x=1`$, and making the change of variable $`\widehat{𝐤}.\widehat{𝐊}=u(1\lambda )/a`$, we obtain for $`F(u)\mathrm{\Delta }_𝐤`$:
$`F(v)={\displaystyle \frac{1}{4a}}{\displaystyle _0^{\mathrm{}}}𝑑uF(u){\displaystyle \frac{M+\mathrm{log}u}{\sqrt{1+2(u+v)+(uv)^2}}}`$ (6)
with $`M=L+\mathrm{log}((1\lambda )/a)1`$. This equation can be solved numerically. However the solution we are looking for behaves approximately as $`\mathrm{exp}(\mu u)`$. Inserting this form into Eq.(5) and requiring that $`F(v)`$ is zero for large $`v`$ gives $`\mathrm{log}\mu =MC`$ (for $`a=0.25`$), where C is the Euler constant. On the other hand requiring $`F(0)=1`$ and making an asymptotic evaluation of the resulting integral gives $`\mu =\mathrm{exp}(\sqrt{2})0.24`$ which is (surprisingly) not much different from the numerical result $`\mu 0.3`$. This leads finally to the asymptotic evaluation $`L1\mathrm{log}(1\lambda )`$, which, together with $`L_011/\mathrm{log}(0.5/(1\lambda ))`$, is in reasonable agreement with our direct solution of Eq.(4) found in Fig.1 . Naturally the divergence of $`\overline{\mu }`$ itself is not to be taken seriously since our calculation requires $`\overline{\mu }/E_F`$ to be small anyway.
Coming back to Fig.1 , we see that beyond $`\lambda 0.8`$ the results from the full Eq.(4) get larger than those of the simplified equation. Hence we find a very large increase of the domain for our FFLO phase. When we take into account that the critical temperature itself will increase rapidly in this range due to the indirect interaction, we see that looking in this region seems quite promising experimentally since the overall domain for the superfluid phase will be much increased. On the other hand the physical properties of this phase will be to a large extent quite different from those found for equal population. A first reason is that this phase is gapless, just as the standard FFLO phase . Next we find an important order parameter anisotropy. Indeed we have plotted in Fig.2, for various values of $`\lambda `$, the angular dependence of $`\mathrm{\Delta }_𝐤`$ with respect to $`𝐊`$. As we have mentionned already, $`\mathrm{\Delta }_𝐤`$ gets more concentrated along the $`𝐊`$ direction when $`\lambda `$ increases. The value of its minimum compared to its maximum is also plotted in Fig.2 . We see that even for moderate values of $`\lambda `$ such as $`\lambda 0.4`$ the anisotropy is quite sizeable. However when $`\lambda `$ is further increased the anisotropy becomes ultimately huge to the point that $`\mathrm{\Delta }_𝐤0`$ for a very large fraction of the Fermi surface. This is a highly unconventional situation and certainly quite unique among BCS superfluids. Another more well-known feature will further complicate the matter. We have in our system a degeneracy with respect to the direction of $`𝐊`$, which will be in general lifted by a texture leading to a spatial inhomogeneity, that is yet another symmetry breaking, as investigated by LO and more recently in Ref. . Together with experimental inhomogeneity due to the trap this will lead to a remarkably complex physical situation.
In conclusion we have shown that in <sup>6</sup>Li the indirect interaction due to density fluctuations exchange will lead to the appearance of a new BCS phase with anisotropic order parameter. Near the instability threshold, this results in a large increase of the superfluid domain as a function of the difference between the atom numbers in the two hyperfine states forming the Cooper pairs.
We have much benefited from discussions with A. J. Leggett. We are grateful to Y. Castin, J. Dalibard, X. Leyronas, C. Mora and C. Salomon for very stimulating conversations.
* Laboratoire associé au Centre National de la Recherche Scientifique et aux Universités Paris 6 et Paris 7.
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# Effective Hamiltonian for striped and paired states at the half-filled Landau level
## I Introduction
The fractional quantum Hall effect (FQHE) is observed around rational filling factors with odd denominators in the lowest Landau level. This effect is caused by the formation of the incompressible liquid state . Recent experiments showed that the compressible state around the half-filled lowest Landau level is similar to the Fermi liquid. In the composite fermion theory, this Fermi-liquid-like state is described by the composite of an electron and two flux quanta, which behaves like a free electron with no magnetic field in the mean field approximation.
In the second Landau level, the situation dramatically changes; the FQHE is observed around the half-filling. The origin of this quantum Hall state is still mysterious. Recent numerical works seem to support the Pfaffian state , which has the spin-polarized p-wave pairing potential, rather than the spin-singlet paired state. However the microscopic explanation for this pairing mechanism is not known yet. More quantitative search is, therefore, needed on the basis of the microscopic BCS-like model. In the present paper, we study the paired state in the mean field theory starting from a microscopic Hamiltonian.
It was shown that the Fermi-liquid-like state can be constructed at an arbitrary filling factor in the von Neumann lattice formalism. In a strong magnetic field, the free kinetic energy is quenched and the kinetic energy is generated from the Coulomb interaction. If the translational symmetry on the von Neumann lattice and U(1) symmetry are unbroken, the Fermi surface is formed in the magnetic Brillouin zone in the mean field theory. In the Hartree-Fock approximation, this state corresponds to the striped state which is observed at half-filled second and higher Landau levels. In Ref. , the first order transition from striped state to paired state is obtained in the numerical calculation of small systems. We regard this striped state as the normal state and examine the U(1) symmetry breaking mechanism. The hopping potential $`\epsilon (p)`$ and gap potential $`\mathrm{\Delta }(p)`$ are determined self-consistently in the Hartree-Fock-Bogoliubov approximation.
Naively it seems that the formation of the paired state by only the repulsive force is impossible in the mean field theory. However, screening effect could make the potential attractive for the pairing. In fact, the quantum effect of fermion loop screens the Coulomb potential and the pseudopotential in the Landau level space becomes attractive. Depending on strength of the screening effect, there are two phases, that is, stripe phase and pairing phase. On the basis of this observation, we construct an effective Hamiltonian which shows a transition from the striped state to the paired state. The transition is continuous and very smooth. The energy spectrum of the quasiparticle is obtained and energy gap is calculated numerically. Furthermore, we find a crossover from the pairing phase to the gap-dominant pairing phase in which the low-energy excitation occurs around zeros of the gap potential.
The paper is organized as follows. In Sec. II, a mean field theory on the von Neumann lattice is presented and self-consistency equations are obtained. In Sec. III, we analyze the self-consistency equations and a screening effect. It is shown that there are a stripe phase and pairing phase depending on strength of the screening effect. In Sec. IV, we present an effective Hamiltonian which shows a transition from the striped state to the paired state and calculate the energy spectrum of the quasiparticle. Summary and discussion are given in Sec. V.
## II Hartree-Fock-Bogoliubov approximation
We consider two-dimensional electron systems in the presence of a perpendicular uniform magnetic field $`B`$. The free particle energy is quenched to the Landau level energy $`E_l=\mathrm{}\omega _c(l+1/2)`$, $`l=0,1,2\mathrm{}`$, where $`\omega _c=eB/m_e`$ (the cyclotron frequency). We suppose that the electrons are spin-polarized and ignore the spin degree of freedom. In this system, the translation is generated by the magnetic translation operators $`T(x,y)`$ with displacement vector $`(x,y)`$. $`T(x,y)`$ commutes with the Hamiltonian. However, magnetic translation operators are non-commutative, that is,
$$T(x_1,y_1)T(x_2,y_2)=e^{i\frac{eB}{\mathrm{}}(x_1y_2x_2y_1)}T(x_2,y_2)T(x_1,y_1).$$
(1)
Therefore, the translation becomes commutative if we restrict the displacement vectors on a two-dimensional lattice which has a unit cell with an area $`2\pi \mathrm{}/eB`$. This lattice is called the von Neumann lattice. We consider lattice sites $`a(m,n)`$, where $`m`$, $`n`$ are integers and lattice spacing $`a=\sqrt{2\pi \mathrm{}/eB}`$. Then $`T(ma,na)`$’s commute with each other and have simultaneous eigenstates as
$$T(ma,na)u_{l,p}(𝐫)=e^{i(p_xm+p_yn)a/\mathrm{}}u_{l,p}(𝐫)$$
(2)
in the $`l`$ th Landau level. We call $`𝐩=(p_x,p_y)`$ the momentum and $`u_{l,p}(𝐫)`$ the Bloch wave on the von Neumann lattice, which is defined on the Brillouin zone $`|p_x|`$, $`|p_y|<\pi \mathrm{}/a`$. It was proved that the eigenstates form an orthogonal complete set . For simplicity, we set $`a=\mathrm{}=c=1`$.
The Bloch wave $`u_{l,p}(𝐫)`$ is constructed by using the coherent states of the guiding center coordinates $`(X,Y)`$ defined by
$`X`$ $`=`$ $`x\xi ,`$ (3)
$`Y`$ $`=`$ $`y\eta ,`$ (4)
$`\xi `$ $`=`$ $`{\displaystyle \frac{1}{eB}}(i{\displaystyle \frac{}{y}}+eA_y),`$ (5)
$`\eta `$ $`=`$ $`{\displaystyle \frac{1}{eB}}(i{\displaystyle \frac{}{x}}+eA_x),`$ (6)
where $`\times A=B`$ and $`(\xi ,\eta )`$ are called the relative coordinates. These operators satisfy the following commutation relations,
$`[\xi ,\eta ]=[X,Y]={\displaystyle \frac{i}{eB}},`$ (7)
$`[\xi ,X]=[\xi ,Y]=[\eta ,X]=[\eta ,Y]=0.`$ (8)
The coherent states of $`(X,Y)`$ are defined by
$`(X+iY)f_{lmn}(𝐫)=(m+in)f_{lmn}(𝐫),`$ (9)
$`{\displaystyle \frac{m_e\omega _c^2}{2}}(\xi ^2+\eta ^2)f_{lmn}(𝐫)=E_lf_{lmn}(𝐫).`$ (10)
Note that these coherent states are non-orthogonal complete set as
$$d^2rf_{l,m^{}+m,n^{}+n}(𝐫)f_{l^{},m^{},n^{}}(𝐫)=\delta _{ll^{}}e^{i\pi (m+n+mn)\frac{\pi }{2}(m^2+n^2)}.$$
(11)
Using $`f_{lmn}(𝐫)`$, the orthogonal basis $`u_{l,p}(𝐫)`$ is given by
$$u_{l,p}(𝐫)=\frac{1}{\beta (p)}\underset{mn}{}e^{ip_xm+ip_yn}f_{lmn}(𝐫),$$
(12)
where $`\beta (p)`$ is a normalization factor. In this basis, two-dimensional momentum is a good quantum number and the Fermi surface is formed in the mean field theory. It was shown that the mean field state could explain the anisotropic compressible state observed at the half-filled second and higher Landau levels. Therefore the basis on the von Neumann lattice is useful tool to investigate the physics at the half-filled second and higher Landau levels.
Fourier transforming the Bloch wave, we can obtain the orthogonal localized basis
$$w_{l,𝐗}(𝐫)=_{\mathrm{BZ}}\frac{d^2p}{(2\pi )^2}u_{l,p}(𝐫)e^{i(p_xm+p_yn)+i\lambda (p)}$$
(13)
where $`𝐗=(m,n)`$, BZ stands for the Brillouin zone, and $`\lambda (p)`$ represents the gauge degree of freedom which is taken to satisfy the periodic condition in the Brillouin zone. We call this basis the Wannier basis. This basis is localized on a position of the lattice site. The explicit forms of these bases are given in Ref. .
Using the Bloch wave basis and Wannier basis, we can expand the electron field operator $`\psi (𝐫)`$ as
$`\psi (𝐫)`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle _{\mathrm{BZ}}}{\displaystyle \frac{d^2p}{(2\pi )^2}}a_l(p)u_{l,p}(𝐫)`$ (14)
$`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{𝐗}{}}b_l(𝐗)w_{l,𝐗}(𝐫),`$ (15)
where $`a_l(p)`$ and $`b_l(𝐗)`$ are anti-commuting annihilation operators in the momentum space and lattice space, respectively. The total Hamiltonian of the present system is
$$H=d^2r\psi ^{}(𝐫)\frac{(i+eA)^2}{2m_e}\psi (𝐫)+\frac{1}{2}d^2rd^2r^{}:(\rho (𝐫)\rho _0)V(𝐫𝐫^{})(\rho (𝐫^{})\rho _0):,$$
(16)
where colons mean the normal ordering, $`\rho (𝐫)=\psi ^{}(𝐫)\psi (𝐫)`$, $`V(𝐫)=q^2/r`$, and $`\rho _0`$ is the uniform background density.
Let us project the system to the $`l`$ th Landau level space. Up to the Landau level energy $`E_l`$, the projected Hamiltonian $`H^{(l)}`$ is written as
$$H^{(l)}=\frac{1}{2}\underset{X_1X_{1}^{}{}_{}{}^{}X_2X_{2}^{}{}_{}{}^{}}{}:b_l^{}(𝐗_1)b_l(𝐗_1^{})V_l(𝐗,𝐘,𝐙)b_l^{}(𝐗_2)b_l(𝐗_2^{}):,$$
(17)
where $`𝐗=𝐗_1𝐗_1^{}`$, $`𝐘=𝐗_2𝐗_2^{}`$, $`𝐙=𝐗_1𝐗_2^{}`$, and $`V_l(𝐗,𝐘,𝐙)`$ is given as
$`V_l(𝐗,𝐘,𝐙)`$ $`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}_{\mathrm{BZ}}\frac{d^2p_1}{(2\pi )^2}\frac{d^2p_2}{(2\pi )^2}\stackrel{~}{v}_l(k)e^{i[p_1X+p_2Y+kZ+f(p_1,p_1+k)+f(p_2+k,p_2)]}},`$ (18)
$`f(p+k,p)`$ $`=`$ $`{\displaystyle _p^{p+k}}(\alpha (p^{})+_p^{}\lambda (p^{}))𝑑p^{}.`$ (19)
Here $`\alpha (p)`$ is a vector potential, $`_p=(\frac{}{p_x},\frac{}{p_y})`$, $`_p\times \alpha (p)=1/2\pi `$ which is a unit flux penetrating the momentum space, and the line integral is along a straight line. The effective Coulomb potential in the $`l`$ th Landau level $`\stackrel{~}{v}_l`$ is given by
$$\stackrel{~}{v}_l(k)=\{L_l(\frac{k^2}{4\pi })\}^2e^{\frac{k^2}{4\pi }}\frac{2\pi q^2}{k},$$
(20)
and $`\stackrel{~}{v}_l(0)=0`$ due to the charge neutrality condition, where $`L_l`$ is the Laguerre polynomial.
The system is translationally invariant in the lattice space and in the momentum space with the uniform magnetic field $`\times \alpha (p)`$. The translational symmetry in the momentum space is called the K-invariance. We consider the case that the K-invariance is spontaneously broken and the kinetic term is induced through the correlation effect.
We apply the mean field approximation to the projected Hamiltonian (17). Let us consider the following mean fields, which are translationally invariant on the von Neumann lattice,
$`U_l(𝐗𝐗^{})`$ $`=`$ $`b_l^{}(𝐗^{})b_l(𝐗),`$ (21)
$`U_l^{(+)}(𝐗𝐗^{})`$ $`=`$ $`b_l^{}(𝐗^{})b_l^{}(𝐗),`$ (22)
$`U_l^{()}(𝐗𝐗^{})`$ $`=`$ $`b_l(𝐗^{})b_l(𝐗).`$ (23)
These mean fields break the K-invariance and U(1) symmetry. $`U_l^{(\pm )}`$ satisfies $`U_l^{()}(𝐗)=U_l^{(+)}(𝐗)`$ and $`U_l^{()}(𝐗)=U_l^{()}(𝐗)`$. Using these mean fields, we approximate $`H^{(l)}`$ as
$`H_{\mathrm{mean}}^{(l)}`$ $`=`$ $`{\displaystyle \underset{XX^{}}{}}U_l(𝐗𝐗^{})\{v_l^H(𝐗^{}𝐗^{})v_l^F(𝐗^{}𝐗)\}b^{}(𝐗)b(𝐗^{})`$ (26)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{XX^{}Y}{}}\{U_l^{(+)}(𝐘)V_l^B(𝐗𝐗^{},𝐘)b_l(𝐗)b_l(𝐗^{})`$
$`+U_l^{()}(𝐘)V_l^B(𝐘,𝐗^{}𝐗)b_l^{}(𝐗)b_l^{}(𝐗^{})\}.`$
Here the Hartree potential $`v_l^H`$ and Fock potential $`v_l^F`$ are given by the following sum rules
$`{\displaystyle \underset{Z}{}}V_l(𝐗,𝐘,𝐙)`$ $`=`$ $`v_l^H(𝐗)\delta _𝐘^𝐗,`$ (27)
$`{\displaystyle \underset{X}{}}V_l(𝐗,𝐘𝐗,𝐙)`$ $`=`$ $`v_l^F(𝐙)\delta _0^𝐘,`$ (28)
where
$`v_l^H(𝐗)`$ $`=`$ $`\stackrel{~}{v}_l(2\pi 𝐗),`$ (29)
$`v_l^F(𝐗)`$ $`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}\stackrel{~}{v}_l(k)e^{ikX}}.`$ (30)
We define the Hartree-Fock potential as $`v_l^{\mathrm{HF}}(𝐗)=v_l^H(𝐗)v_l^F(𝐗)`$ and its Fourier transform as $`\stackrel{~}{v}_l^{\mathrm{HF}}(p)`$. The Bogoliubov potential $`V_l^B`$ is written as
$`V_l^B(𝐗,𝐘)`$ $`=`$ $`{\displaystyle \underset{Z}{}}V_l(𝐗+𝐙,𝐘+𝐙,𝐙)`$ (31)
$`=`$ $`{\displaystyle _{\mathrm{BZ}}}{\displaystyle \frac{d^2p_1}{(2\pi )^2}}{\displaystyle \frac{d^2p_2}{(2\pi )^2}}\stackrel{~}{V}_l^B(p_1,p_2)e^{ip_1X+ip_2Y},`$ (32)
and $`\stackrel{~}{V}_l^B(p_1,p_2)`$ is given by
$$\stackrel{~}{V}_l^B(p_1,p_2)=\underset{N}{}\stackrel{~}{v}_l(p_1+p_2+2\pi N)e^{i[f(p_1,p_22\pi N)+f(p_1,p_2+2\pi N)]},$$
(33)
where $`N=(N_x,N_y)`$, $`N_x`$, $`N_y`$ are integers. Then we introduce the hopping potential $`\epsilon `$ and gap potential $`\mathrm{\Delta }`$ as
$`\epsilon _{l,𝐗}`$ $`=`$ $`U_l(𝐗)v_l^{HF}(𝐗),`$ (34)
$`\mathrm{\Delta }_{l,𝐗}`$ $`=`$ $`{\displaystyle \underset{Y}{}}U_l^{()}(𝐘)V_l^B(𝐗,𝐘).`$ (35)
We define these Fourier transforms as $`\epsilon _l(p)`$ and $`\mathrm{\Delta }_l(p)e^{i(\lambda (p)+\lambda (p))}`$, respectively. The phase factor of the gap potential is the same in Eq. (13). $`V_l^B`$ satisfies $`V_l^B(𝐗,𝐘)=V_l^B(𝐗,𝐘)`$ and we have $`\mathrm{\Delta }_{l,X}=\mathrm{\Delta }_{l,X}`$. Using these potentials, the mean field Hamiltonian reads
$$H_{\mathrm{mean}}^{(l)}=\underset{XX^{}}{}[\epsilon _{l,𝐗𝐗^{}}b_l^{}(𝐗)b_l(𝐗^{})+\frac{1}{2}\mathrm{\Delta }_{l,𝐗^{}𝐗}^{}b(𝐗)b(𝐗^{})+\frac{1}{2}\mathrm{\Delta }_{l,𝐗^{}𝐗}b^{}(𝐗^{})b^{}(𝐗)].$$
(36)
This Hamiltonian is diagonalized by the Bogoliubov transformation in the momentum space. We fix the gauge of vector potential $`\alpha `$ in the momentum space as $`\alpha (p)=(p_y/2\pi ,0)`$. In this gauge, the boundary conditions in the Brillouin zone are given by
$`a_l(p+2\pi N)`$ $`=`$ $`e^{i\pi (N_x+N_y)iN_yp_x}a_l(p),`$ (37)
$`\mathrm{\Delta }_l(p+2\pi N)`$ $`=`$ $`e^{2iN_yp_x}\mathrm{\Delta }_l(p).`$ (38)
Equation (38) and the relation $`\mathrm{\Delta }_l(p)=\mathrm{\Delta }_l(p)`$ make the gap potential $`\mathrm{\Delta }_l(p)`$ have four zeros at $`𝐩=(0,0)`$, $`(0,\pi )`$, $`(\pi ,0)`$, and $`(\pi ,\pi )`$. We assume $`\epsilon _l(p)=\epsilon _l(p)`$. $`U_l`$ and $`U_l^{()}`$ are calculated by using $`H_{\mathrm{mean}}^{(l)}`$. Then, the hopping potential and gap potential in Eqs. (34) and (35) are determined by the self-consistency equations,
$`\xi _l(p)`$ $`=`$ $`{\displaystyle _{\mathrm{𝐁𝐙}}}{\displaystyle \frac{d^2p_1}{(2\pi )^2}}{\displaystyle \frac{E_l(p_1)\xi _l(p_1)}{2E_l(p_1)}}\stackrel{~}{v}_l^{\mathrm{HF}}(p_1p)\mu ,`$ (39)
$`\mathrm{\Delta }_l(p)`$ $`=`$ $`{\displaystyle \frac{d^2p_1}{(2\pi )^2}\frac{\mathrm{\Delta }_l(p_1)}{2E_l(p_1)}\stackrel{~}{v}_l(p_1p)e^{i_p^{p_1}2\alpha (k)𝑑k}},`$ (40)
where $`\mu `$ is the chemical potential and $`E_l(p)`$ is the spectrum of the quasiparticle which is defined by
$`E_l(p)`$ $`=`$ $`\sqrt{\xi _l(p)^2+|\mathrm{\Delta }_l(p)|^2},`$ (41)
$`\xi _l(p)`$ $`=`$ $`\epsilon _l(p)\mu .`$ (42)
These self-consistency equations are equivalent to sum up infinite diagrams of the Fermion self-energy part. Note that the integral region of Eq. (40) are infinite. Note also that the gauge field appears in the gap equation (40). The gauge field represents two unit flux on the Brillouin zone. This corresponds to the flux carried by the Cooper pair $`a(p)a(p)`$ which is separated between $`p`$ and $`p`$ in the Brillouin zone. The mean field $`a^{}(p)a(p)`$, on the other hand, is placed on the same point. Therefore no gauge field appears in Eq. (39).
The filling factor of the $`l`$ th Landau level is denoted by $`\nu _l`$ and total filling factor is given by $`\nu =l+\nu _l`$. Using the mean fields, the filling factor is given by
$$\nu _l=\frac{1}{2}_{\mathrm{𝐁𝐙}}\frac{d^2p}{(2\pi )^2}\frac{\xi _l(p)}{2E_l(p)}.$$
(43)
In the limit of $`\mathrm{\Delta }0`$, $`\nu _l`$ is equal to $`_{\mathrm{𝐁𝐙}}\theta (\mu \epsilon _l(p))d^2p/(2\pi )^2`$.
## III Stripe phase and pairing phase
In this section, we analyze the self-consistency equations (39) and (40). It is shown that the screening effect plays important roles for the pairing mechanism. For simplicity, we omit the Landau level index $`l`$ and use $`q^2/a`$ as the unit of energy.
It is convenient for solving Eq. (39) to use the following mode expansions,
$`\xi (p)`$ $`=`$ $`{\displaystyle \underset{X}{}}\xi _𝐗e^{iXp},`$ (44)
$`{\displaystyle \frac{\xi (p)}{E(p)}}`$ $`=`$ $`{\displaystyle \underset{X}{}}\eta _𝐗e^{iXp},`$ (45)
where $`\xi _𝐗=\xi _𝐗`$, $`\eta _𝐗=\eta _𝐗`$, and $`\xi _𝐗`$, $`\eta _𝐗`$ are real numbers. Then Eq. (39) becomes
$`\xi _𝐗`$ $`=`$ $`{\displaystyle \frac{v^{\mathrm{HF}}(𝐗)}{2}}\eta _𝐗,\mathrm{for}𝐗0,`$ (46)
$`\xi _\mathrm{𝟎}`$ $`=`$ $`{\displaystyle \frac{v^{\mathrm{HF}}(\mathrm{𝟎})}{2}}(\eta _\mathrm{𝟎}1)\mu .`$ (47)
The zero mode $`\eta _\mathrm{𝟎}`$ is related to the filling factor as $`\nu _l=(1\eta _\mathrm{𝟎})/2`$ by Eq. (43). Therefore the zero mode $`\xi _\mathrm{𝟎}`$ is determined by $`\mu `$ and $`\nu _l`$. In particular, $`\eta _\mathrm{𝟎}=0`$ for the half-filling case ($`\nu _l=1/2`$). We introduce hopping parameters $`t_𝐗=2\xi _𝐗`$ and regard $`\xi _𝐗`$ as a functional of $`t_𝐗`$’s and $`\mathrm{\Delta }(p)`$. Then the self-consistency for $`t_𝐗`$ is written as
$$t_𝐗=2\xi _𝐗[\{t_{(m,n)}\},\mathrm{\Delta }(p)].$$
(48)
Using $`t_{(m,n)}`$, the hopping term is written by $`_{m,n0}t_{(m,n)}\mathrm{cos}(p_xm+p_yn)a^{}(p)a(p)`$.
Next we discuss the self-consistency for the gap potential $`\mathrm{\Delta }(p)`$. Eq. (40) is rewritten as
$$\mathrm{\Delta }(p)=\frac{d^2k}{(2\pi )^2}\stackrel{~}{v}(k)e^{ikD}\frac{\mathrm{\Delta }(p)}{2E(p)},$$
(49)
where $`𝐃=(i\frac{}{p_x}+2\alpha _x,i\frac{}{p_y}+2\alpha _y)`$. It is convenient to introduce eigenfunctions of the operator $`𝐃^2`$, that is, $`𝐃^2\psi _n(p)=e_n\psi _n(p)`$ with $`e_n=(2n+1)/\pi `$, $`n=0,1,2,\mathrm{}`$. The index $`n`$ labels the $`n`$ th Landau level in the momentum space. The eigenfunctions which obey the same boundary condition as the gap function are doubly degenerate and are given by
$`\psi _n^{(2)}(p)`$ $`=`$ $`N_n{\displaystyle \underset{N_y}{}}H_n({\displaystyle \frac{p_y+(2N_y+1)\pi }{\sqrt{2\pi }}})e^{i(2N_y+1)p_x\frac{1}{2\pi }(p_y+(2N_y+1)\pi )^2}`$ (50)
$`=`$ $`{\displaystyle \frac{2\pi }{\sqrt{n!}}}({\displaystyle \frac{\pi }{2}})^{n/2}(D_xiD_y)^n\psi _0^{(2)}(p),`$ (51)
$`\psi _n^{(3)}(p)`$ $`=`$ $`N_n{\displaystyle \underset{N_y}{}}H_n({\displaystyle \frac{p_y+2N_y\pi }{\sqrt{2\pi }}})e^{i2N_yp_x\frac{1}{2\pi }(p_y+2N_y\pi )^2}`$ (52)
$`=`$ $`{\displaystyle \frac{2\pi }{\sqrt{n!}}}({\displaystyle \frac{\pi }{2}})^{n/2}(D_xiD_y)^n\psi _0^{(3)}(p),`$ (53)
where $`H_n`$ is the Hermite polynomial, $`N_n=1/\sqrt{2^{n1}n!}`$, and $`\psi _0^{(2)}`$, $`\psi _0^{(3)}`$ are written by theta functions as
$`\psi _0^{(2)}(p)`$ $`=`$ $`N_0e^{\frac{p_y^2}{2\pi }}\vartheta _2({\displaystyle \frac{p_x+ip_y}{\pi }}|2i),`$ (54)
$`\psi _0^{(3)}(p)`$ $`=`$ $`N_0e^{\frac{p_y^2}{2\pi }}\vartheta _3({\displaystyle \frac{p_x+ip_y}{\pi }}|2i).`$ (55)
These eigenfunctions have the parity symmetry, $`\psi _n(p)=()^n\psi _n(p)`$. Therefore the gap potential can be expanded by $`\psi _{2n+1}(p)`$. Furthermore the eigenfunctions satisfy the following relations
$`\psi _n^{(2)}(p_x+\pi ,p_y)`$ $`=`$ $`\psi _n^{(2)}(p_x,p_y),`$ (56)
$`\psi _n^{(3)}(p_x+\pi ,p_y)`$ $`=`$ $`\psi _n^{(3)}(p_x,p_y),`$ (57)
$`\psi _n^{(2)}(p_x,p_y+\pi )`$ $`=`$ $`e^{ip_x}\psi _n^{(3)}(p_x,p_y),`$ (58)
$`\psi _n^{(3)}(p_x,p_y+\pi )`$ $`=`$ $`e^{ip_x}\psi _n^{(2)}(p_x,p_y),`$ (59)
and $`\psi _n(\pi m,\pi n)=0`$. Moreover $`\psi _n^{(2)}(\pi n/2,\pi )=0`$ and $`\psi _n^{(3)}(\pi n/2,0)=0`$. We can expand the gap potential as
$`\mathrm{\Delta }(p)`$ $`=`$ $`{\displaystyle \underset{n1,i=2,3}{}}c_n^{(i)}\psi _{2n1}^{(i)}(p),`$ (60)
$`{\displaystyle \frac{\mathrm{\Delta }(p)}{E(p)}}`$ $`=`$ $`{\displaystyle \underset{n1,i=2,3}{}}d_n^{(i)}\psi _{2n1}^{(i)}(p).`$ (61)
Then Eq. (40) becomes
$`c_n^{(i)}`$ $`=`$ $`{\displaystyle \frac{F_{2n+1}}{2}}d_n^{(i)},`$ (62)
$`F_n`$ $`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}\stackrel{~}{v}(k)L_n(\frac{k^2}{2\pi })e^{\frac{k^2}{4\pi }}}.`$ (63)
Note that although $`F_n`$’s have the same form as the Haldane’s pseudopotentials, the physical meaning is different. Haldane’s one is the 2-body interaction projected to the relative angular momentum in the single Landau level. It is not screened by definition. The potential $`F_n`$, on the other hand, is a potential appearing in the gap equation which is renormalized by higher order corrections. Actually the potential $`F_n`$ is screened by the polarization $`\mathrm{\Pi }(p)`$ due to the Fermion loop diagrams. We approximate the screened potential as
$`\stackrel{~}{v}(p,m_{\mathrm{TF}})`$ $`=`$ $`1/(\stackrel{~}{v}(p)^1+m_{\mathrm{TF}}),`$ (64)
$`m_{\mathrm{TF}}`$ $`=`$ $`\mathrm{\Pi }(0),`$ (65)
where $`m_{\mathrm{TF}}`$ is the Thomas-Fermi mass. Calculating the one-loop diagram, $`m_{\mathrm{TF}}`$ is given by
$$m_{\mathrm{TF}}=_{\mathrm{𝐁𝐙}}\frac{d^2p}{(2\pi )^2}\frac{|\mathrm{\Delta }(p)|^2}{2E(p)^3}$$
(66)
We use the screened potential $`\stackrel{~}{v}(p,m_{\mathrm{TF}})`$ in the Fock potential (30) and Bogoliubov potential (33). In the zero limit of the gap potential $`\mathrm{\Delta }`$, $`m_{\mathrm{TF}}`$ becomes the density of states on the Fermi surface. It should be noted that the Hartree potential $`v^H`$ is not screened because the polarization effect is automatically included in the self-consistency equation.
The $`m_{\mathrm{TF}}`$ dependences of the hopping strength $`v^{\mathrm{HF}}(𝐗)`$ and pseudopotentials $`F_n`$ are plotted in Figs. (1)-(6) for $`l=`$0, 1, 2. As seen in these figures, potentials change their signs at $`m_{\mathrm{TF}}1.02.0`$. Especially for $`l=`$1 and 2, all potentials change their signs together at large $`m_{\mathrm{TF}}`$. This observation leads us to a strong statement for $`l=1`$ and 2. Using Eqs. (44)-(47) and Eqs. (60)- (63), the following inequalities are concluded at $`\nu _l=1/2`$,
$`{\displaystyle \frac{\xi (p)^2}{E(p)}\frac{d^2k}{(2\pi )^2}}`$ $`=`$ $`{\displaystyle \underset{X0}{}}{\displaystyle \frac{v^{\mathrm{HF}}(𝐗)}{2}}\eta _𝐗^20,`$ (67)
$`{\displaystyle \frac{|\mathrm{\Delta }(p)|^2}{E(p)}\frac{d^2k}{(2\pi )^2}}`$ $`=`$ $`{\displaystyle \underset{n1,i=2,3}{}}{\displaystyle \frac{F_n}{2}}|d_n^{(i)}|^20.`$ (68)
Therefore, $`d_n^{(i)}=0`$ for all $`n`$ and $`i`$ at small $`m_{\mathrm{TF}}`$ and we obtain U(1) symmetric phase ($`\mathrm{\Delta }(p)=0`$). This phase was studied in the Hartree-Fock approximation. It was shown that the compressible striped state is favored in this phase. At large $`m_{\mathrm{TF}}`$ for $`l=`$1 and 2, $`\eta _𝐗=0`$ for $`𝐗0`$ and we obtain a pairing phase. Thus, in the Hartree-Fock-Bogoliubov approximation, the stripe phase makes transition to the pairing phase for $`m_{\mathrm{TF}}>1.02.0`$. In this pairing phase, $`\xi =0`$ and the energy spectrum of the quasiparticle becomes $`|\mathrm{\Delta }(p)|`$. Hence, there are at least four gapless points in the Brillouin zone for the obtained paired state. This peculiar result conflicts with the experiment and numerical calculations, in which the paired state has an energy gap. Furthermore we find that the self-consistency for $`m_{\mathrm{TF}}`$ (66) is not satisfied, which goes to infinity as iterating the numerical calculations. This means that the fluctuation is too large and the mean field solution is not stable. In the next section, we change the summation of infinite diagrams of the Fermion self-energy part and obtain a self-consistent paired state with an energy gap.
## IV Effective Hamiltonian for the striped and paired state
In this section, we present an effective Hamiltonian which has a striped state as the normal state and paired state as the U(1) symmetry breaking state. We focus our argument on the half-filled second Landau level space, that is $`l=1`$, $`\nu =1+1/2`$. We use a tricky technique to find a self-consistent solution for the paired state with an energy gap. The most hardest problem in the quantum Hall system is to derive a low-energy effective theory from the microscopic Hamiltonian, because we do not know a small parameter used in the perturbation expansion a priori. Therefore we have to choose a starting point by consideration of symmetry and physical intuition based on experiments. In the present case, we maintain the translational symmetry on the von Neumann lattice and choose the striped state as the normal state.
Let us divide the mean field Hamiltonian (36) into two parts as
$`H_{\mathrm{mean}}`$ $`=`$ $`H_0+H_1,`$ (69)
$`H_1`$ $`=`$ $`{\displaystyle _{\mathrm{BZ}}}{\displaystyle \frac{d^2p}{(2\pi )^2}}t\mathrm{cos}p_ya^{}(p)a(p),`$ (70)
and $`H_1`$ is treated as perturbation to $`H_0=H_1+H_{\mathrm{mean}}`$. Self-consistency is imposed for $`H_0`$ first and correction from $`H_1`$ is included next. $`H_1`$ represents the anisotropy in the stripe state or normal state in which the uniform direction is chosen in the $`y`$ direction. In the striped state at the half-filling, the Fermi sea is formed at $`|p_y|<\pi /2`$, which is shown in Ref. In this case the Fermi surface (line) is formed at $`p_y=\pm \pi /2`$. It is expected that the paired state possesses the same anisotropy as the normal state. In $`H_0`$, the nearest-neighbor (N-N) hopping parameter becomes $`t2\xi _{(0,1)}=t+t_{(0,1)}`$. Then the self-consistent hopping parameter satisfies
$$t_𝐗=2\xi _𝐗[\{t\delta _{m,n}^{0,1}+t_{(m,n)}\},\mathrm{\Delta }(p)].$$
(71)
The N-N hopping parameter is renormalized as $`t^{}=t+t_{(0,1)}`$ in $`H_0`$. Including the first order correction of $`H_1`$, the effective hopping parameter becomes $`t_{\mathrm{eff}}=t^{}t=t_{(0,1)}`$. The parameter $`t`$ controls the effective N-N hopping strength dynamically. The magnitude of $`t_{\mathrm{eff}}`$ corresponds to the strength of the stripe order.
We truncate the expansions of Eqs. (44) and (60) and calculate the self-consistent solution by iteration of numerical calculations until we obtain convergence. Considering only the lowest and next relevant terms, the effective Hamiltonian for the quasiparticle in the striped and paired state is given by
$$H_{\mathrm{eff}}=_{\mathrm{BZ}}\frac{d^2p}{(2\pi )^2}[\epsilon _{\mathrm{eff}}(p)a^{}(p)a(p)+\frac{1}{2}\mathrm{\Delta }_{\mathrm{eff}}(p)a^{}(p)a^{}(p)+\frac{1}{2}\mathrm{\Delta }_{\mathrm{eff}}^{}(p)a(p)a(p)],$$
(72)
where the effective hopping potential and effective gap potential are given by
$`\epsilon _{\mathrm{eff}}(p)`$ $`=`$ $`t_{\mathrm{eff}}\mathrm{cos}p_yt_{(0,3)}\mathrm{cos}3p_y,`$ (73)
$`\mathrm{\Delta }_{\mathrm{eff}}(p)`$ $`=`$ $`c_1^{(3)}\psi _1^{(3)}(p)+c_3^{(3)}\psi _3^{(3)}(p).`$ (74)
In using Eqs. (73) and (74), $`\psi _n^{(2)}`$ and $`\mathrm{cos}((2m+1)p_x+np_y)`$ terms are irrelevant and $`\mathrm{cos}(2mp_x+2np_y)`$ terms induced in $`\epsilon _{\mathrm{eff}}`$ give a negligible correction to the numerical results. The $`t`$ dependence of $`t_{\mathrm{eff}}`$ is determined self-consistently by $`t_{\mathrm{eff}}=2\xi _{(0,1)}[t+t_{\mathrm{eff}},t_{(0,3)},c_1^{(3)},c_3^{(3)}]`$. The $`t`$ dependence is implicit in the effective theory. The hopping parameters $`t_{(0,n)}`$ and gap potential $`\mathrm{\Delta }_{\mathrm{eff}}(p)`$ depend on $`t_{\mathrm{eff}}`$ explicitly. We use the following approximation
$$m_{\mathrm{TF}}=\frac{1}{\pi t_{\mathrm{eff}}},$$
(75)
which is satisfied in the stripe phase with only the N-N hopping term. We checked that Eq. (75) was good approximation using the obtained convergent solutions.
Chemical potential $`\mu `$ is determined so that the filling factor $`\nu _1`$ is equal to $`1/2`$. Note that $`\mu `$ includes the on-site term in $`H_{\mathrm{mean}}`$. We find that $`\mu `$ is negative small number on the order of $`10^4`$ at most. The gap potential $`\mathrm{\Delta }_{\mathrm{eff}}(p)`$ depends on $`t_{\mathrm{eff}}`$. Figure (7) shows the $`t_{\mathrm{eff}}`$ dependence of the energy gap $`\mathrm{\Delta }E=\mathrm{min}(2E(p))`$ at $`\nu =1+1/2`$. The maximum value of the energy gap is $`0.027`$ at $`t_{\mathrm{eff}}=0.03`$. In this case, $`t_{(0,3)}=0.0003`$, $`c_1^{(3)}=0.0104`$, $`c_3^{(3)}=0.0018`$, and $`\mu =0.0006`$. The energy spectrum of the quasiparticle at $`t_{\mathrm{eff}}=0.03`$ is shown in Fig. (8). The absolute value of the gap potential at $`t_{\mathrm{eff}}=0.03`$ is shown in Fig. (9). The excitation energy $`E(p)`$ becomes small around $`p_y=\pm \pi /2`$, which is the Fermi surface of the striped state. The transition to the stripe phase is continuous and very smooth. Near the transition point, the behavior of energy gap is approximated by $`t_{\mathrm{eff}}e^{2\pi t_{\mathrm{eff}}/|F_1|}`$, whose non-perturbative dependence on the coupling is well-known in the BCS theory. At $`m_{\mathrm{TF}}=m_c1.4`$, $`F_1`$ behaves as $`\alpha (m_cm_{\mathrm{TF}})`$ as seen in Fig. (4), and the energy gap approaches to zero as
$$\mathrm{\Delta }Et_{\mathrm{eff}}e^{\frac{2\pi t_{\mathrm{eff}}^2}{\alpha m_c(t_ct_{\mathrm{eff}})}},\mathrm{for}t_{\mathrm{eff}}<t_c,$$
(76)
where $`t_c1/\pi m_c0.2`$. The energy gap is extremely small at $`0.1<t_{\mathrm{eff}}<t_c`$. At $`t_{\mathrm{eff}}>t_c`$, the gap potential vanishes and the compressible striped state is realized. The spectrum of the striped state is uniform in the $`p_x`$ direction. This state is regarded as a collection of the one-dimensional lattice Fermion systems and leads to the anisotropy of the magnetoresistance .
Inspecting the energy spectrum of quasiparticle in the pairing phase, we find a crossover phenomenon at $`t_{\mathrm{eff}}0.01`$, that is, the minimum excitation energy $`\mathrm{min}(E(p))`$ is placed around $`p_y=\pm \pi /2`$ at $`0.01<t_{\mathrm{eff}}<t_c`$, whereas at $`0<t_{\mathrm{eff}}<0.01`$, placed around $`p_y=0`$ and $`\pi `$. We call the latter case the gap-dominant pairing phase. In this phase, the low-energy excitation occurs around zeros of the gap potential and the the energy gap is given by $`2\xi (0)`$. Therefore the $`t_{\mathrm{eff}}`$ dependence of energy gap $`\mathrm{\Delta }E`$ becomes linear at $`0<t_{\mathrm{eff}}<0.01`$ The energy spectrum at $`t_{\mathrm{eff}}0.01`$ is shown in Fig. (10). As seen in this figure, the spectrum is close to the flatband and the stripe order is weakened compared with Fig. (8).
## V Summary and discussion
We propose an effective Hamiltonian which describes the striped and paired state at the half-filled Landau levels. The gap potential and screening mass are determined self-consistently in the Hartree-Fock-Bogoliubov approximation scheme. The energy spectrum of the quasiparticle and energy gap are calculated numerically.
The very smooth transition from pairing to stripe phase occurs at $`t_{\mathrm{eff}}=t_c`$. In other words, the stripe order strongly suppresses the pairing order at $`0.1<t_{\mathrm{eff}}<t_c`$. The energy gap becomes maximum at $`t_{\mathrm{eff}}=0.03`$ and the maximum value is $`0.027q^2/a=0.01q^2/l_B`$ ($`l_B=\sqrt{\mathrm{}/eB}`$) which is the same order as Morf’s value. The experimental value at $`\nu =5/2`$ is smaller than the theoretical one by an order of magnitude. This difference could be explained by effects of disorder and finite thickness of 2D layer, which reduce the energy gap . The regions of small $`t_{\mathrm{eff}}`$ and $`0.1<t_{\mathrm{eff}}<t_c`$ have another possibility to explain the experiment. In our model, the crossover is obtained for small $`t_{\mathrm{eff}}`$. This suggests that the stripe order is destroyed by the pairing order and the system becomes isotropic in the gap-dominant pairing phase.
The electric charge of the quasiparticle is a half of the electron charge and the Hall conductance is quantized as $`e^2/2h`$. The results are consistent with the recent experiment. However, we do not know whether correction to the quantized Hall conductance is finite or zero. Related to this subject, $`U(1)`$ breaking effect in the $`PT`$ violating system is studied in . Whether the Hall conductance in the paired state is topological invariant or not is an interesting future problem.
It is possible to consider a pairing mechanism based on the composite fermion picture. However, it is unclear that the composite fermion picture is applicable to the second and higher Landau levels. It rather seems that the compressible stripe state which has an anisotropic Fermi surface is the normal state at the half-filled second and higher Landau levels in the Hartree-Fock approximation and numerical calculations of small systems. In fact, the transition from the paired state to striped state is observed in the presence of an inplane magnetic field. Furthermore Ref. suggests that the Chern-Simons gauge fluctuations are strongly pair breaking. Controversially, a superfluid state at the half-filling Hall state is proposed on the basis of the dipolar liquid picture of the composite fermion.
Relation between our model and the Pfaffian state is not understood yet. Our effective Hamiltonian for the striped and paired state resembles the effective Hamiltonian proposed by Read and Green, which has the Pfaffian state in the weak pairing phase. However their normal state has the isotropic Fermi surface, which is reminiscent of the composite Fermion theory, and is different from our anisotropic one. The p-wave behavior seen in Eqs. (54) and (55), on the other hand, is in common with the Pfaffian state. In our case, $`\mathrm{\Delta }(p)`$ becomes $`c_xp_x+ic_yp_y`$, $`c_yc_x`$ in the long wave length limit $`p0`$.
More experimental and theoretical studies are needed to decide the symmetry of the gap potential and to understand the underlying physics at the half-filled second Landau level. To compare the theory with experiments quantitatively, effects of the Landau level mixing, finite temperature, finite thickness of 2D layer, and inplane magnetic field need to be included in the calculation. Furthermore, Josephson effect, vortex excitations, and edge states in our model are important future subjects.
###### Acknowledgements.
I would like to thank K. Ishikawa, T. Aoyama and especially J. Goryo for useful discussions. Part of this work was done when I took part in Tenth International Conference on Recent Progress in Many-Body Theories, Seattle. This work was partially supported by the special Grant-in-Aid for Promotion of Education and Science in Hokkaido University provided by the Ministry of Education, Science, Sport, and Culture, the Grant-in-Aid for Scientific Research on Priority area (Physics of CP violation) (Grant No. 12014201), and the Grant-in aid for International Science Research (Joint Research 10044043) from the Ministry of Education, Science, Sports, and Culture, Japan.
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# Dynamics of the chiral phase transition at finite chemical potential
## Abstract
We study the dynamics of the chiral phase transition at finite chemical potential in the Gross-Neveu model in the leading order in large-$`N`$ approximation. We consider evolutions starting in local thermal and chemical equilibrium in the massless unbroken phase for conditions pertaining to traversing a first or second order phase transition. We assume boost invariant kinematics and determine the evolution of the order parameter $`\sigma `$, the energy density and pressure as well as the effective temperature, chemical potential and interpolating number densities as a function of $`\tau `$.
2
The phase structure of QCD at non-zero temperature and baryon density is important for the physics of neutron stars and relativistic heavy ion collisions. The phase structure for two massless quarks reveals a rich structure. At low temperature and chemical potential, the ground state has broken chiral symmetry. At higher chemical potential one finds a superconducting phase. The transition out of the chirally broken phase as one increases the temperature is second order at low chemical potential and then changes to first order as we increase the chemical potential .
Recently we found a simple model which has a similar phase structure to that described above, i.e. both chiral and superconducting transitions as well as asymptotic freedom. Here we consider a special limit without a superconducting phase, where the model reduces to the Gross-Neveu (GN) model whose Lagrangian is
$$=i\overline{\mathrm{\Psi }}_i\gamma ^\mu _\mu \mathrm{\Psi }_i\frac{1}{2}g^2\left(\overline{\mathrm{\Psi }}_i\mathrm{\Psi }^i\right)^2,$$
(1)
which is invariant under the discrete chiral group: $`\mathrm{\Psi }_i\gamma _5\mathrm{\Psi }_i`$. In leading order in large $`N`$ the effective action is
$$S_{eff}=d^2x\left[i\overline{\mathrm{\Psi }}_i\left(\partial ̸+\sigma \right)\mathrm{\Psi }^i\frac{\sigma ^2}{2g^2}\right]+\mathrm{trln}S^1[\sigma ],$$
(2)
where $`S^1(x,y)[\sigma ]=\left(\gamma ^\mu _\mu +\sigma \right)\delta (xy)`$.
The phase structure of the GN model at finite temperature and chemical potential in this approximation has been known for a long time and is summarized in Fig. 1.
The phase structure is determined from the renormalized effective potential
$`V_{eff}(\sigma ^2,T,\mu )={\displaystyle \frac{\sigma ^2}{4\pi }}[\mathrm{ln}{\displaystyle \frac{\sigma ^2}{m_f^2}}1]`$ (3)
$`{\displaystyle \frac{2}{\beta }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk}{2\pi }}[\mathrm{ln}(1+e^{\beta (E\mu )})+\mathrm{ln}(1+e^{\beta (E+\mu )})]`$ (4)
Here $`m_f`$ is the physical mass of the fermion in the vacuum sector. The tricritical point occurs at $`\frac{\mu _c}{m_f}=.608,`$ $`\frac{T_c}{m_f}=.318.`$ We have chosen to renormalize the effective potential so its value at $`T=0`$ in the false vacuum $`\sigma =0`$ is zero. In the true vacuum $`\sigma =m_f`$ the energy density has the value $`ϵ/m_f^2=\frac{1}{4\pi }.`$
Following a heavy ion collision, the ensuing plasma expands and cools traversing the chiral phase transition. In hydrodynamic simulations of these collisions, a reasonable approximation is to treat the expansion as a 1+1 dimensional boost invariant expansion along the beam ($`z`$) axis. In this approximation, the fluid velocity scales as $`z/t`$. In terms of the variables fluid rapidity $`\eta =\frac{1}{2}\mathrm{ln}\left(\frac{t+z}{tz}\right)`$ and fluid proper time $`\tau =(t^2z^2)^{1/2}`$, physical quantities such as $`\sigma ,ϵ`$ become independent of $`\eta `$, as discussed in refs. and applied to the problem of disoriented chiral condensates in ref.. We note that related nonequilibrium techniques have also been developed in ref. and applied to the problem of disoriented chiral condensates in ref. . Although the effective mass $`\sigma `$ is a function solely of $`\tau `$, two-point correlation functions depend on $`\eta `$ as well.
We shall use the metric convention $`(,+)`$. In our approximation, the dynamics are described by the Dirac equation with self-consistently determined mass term. Rescaling the fermion field, $`\psi (x)=\frac{1}{\sqrt{\tau }}\mathrm{\Phi }(x),`$ and introducing conformal time $`u`$ via $`\tau =\frac{e^u}{m}`$, we obtain
$$\left[\gamma ^0_u+\gamma ^3_\eta +\stackrel{~}{\sigma }(u)\right]\mathrm{\Phi }(x)=0,$$
(5)
where $`\stackrel{~}{\sigma }(u)=\sigma \tau =\frac{\sigma }{m}e^u`$ and $`\stackrel{~}{m_f}=m_f\tau .`$
Further letting $`g^2=\lambda /2N`$ we have the gap equation
$$\sigma =i\frac{\lambda }{2N}[\mathrm{\Psi }_i^{},\gamma ^0\mathrm{\Psi }_i]i\frac{\lambda }{2}[\psi ^{},\gamma ^0\psi ],$$
(6)
where we have assumed $`N`$ identical $`\mathrm{\Psi }_i=\psi `$.
These equations are to be solved subject to initial conditions at $`\tau =\tau _0`$. It is sufficient to describe the initial state of the charged fermion field by the initial particle and anti-particle number densities, which we take to be Fermi-Dirac distributions described by $`\mu _0`$ and $`T_0`$.
Expanding the fermion fields $`\mathrm{\Phi }`$ in terms of Fourier modes at fixed conformal time $`u`$,
$$\mathrm{\Phi }(x)=\frac{dk_\eta }{2\pi }\{b(k)\varphi _k^+(u)e^{ik_\eta \eta }+d^{}(k)\varphi _k^{}(u)e^{ik_\eta \eta }\},$$
(7)
the $`\varphi _k^\pm `$ then obey
$$\left[\gamma ^0\frac{d}{du}+i\gamma ^3k_\eta +\stackrel{~}{\sigma }(u)\right]\varphi _k^\pm (u)=0.$$
(8)
The superscript $`\pm `$ refers to positive- or negative-energy solutions. Introducing mode functions $`\varphi _k^\pm (u)`$ via
$$\varphi _k^\pm (u)=\left[\gamma ^0\frac{d}{du}i\gamma ^3k_\eta +\stackrel{~}{\sigma }(u)\right]f_k^\pm (\tau )\chi ^\pm ,$$
(9)
where the momentum independent spinors $`\chi ^\pm `$ are chosen to be the orthornomal $`\pm 1`$ eigenstates of $`i\gamma ^0`$, we obtain the second order equations:
$$\left(\frac{d^2}{du^2}\stackrel{~}{\omega }_k^2\pm i\frac{d\stackrel{~}{\sigma }}{du}\right)f_k^\pm (u)=0,$$
(10)
where $`\stackrel{~}{\omega }_k^2=k_\eta ^2+\stackrel{~}{\sigma }^2(u).`$ We parameterize the positive-energy solutions $`f_k^+`$ in a similar manner to Eq. (3.1) of Ref. :
$`f_k^+(u)={\displaystyle \frac{N_k}{\sqrt{2\stackrel{~}{\mathrm{\Omega }}_k(u)}}}\mathrm{exp}\left\{{\displaystyle _0^u}\left(i\stackrel{~}{\mathrm{\Omega }}_k(u^{}){\displaystyle \frac{\dot{\stackrel{~}{\sigma }}(u^{})}{2\stackrel{~}{\mathrm{\Omega }}_k(u^{})}}\right)𝑑u^{}\right\}.`$ (11)
Using eqs.(7, 9)and the definitions: $`b^{}(k)b(q)=2\pi \delta (kq)N_+(q)`$ and $`d^{}(k)d(q)=2\pi \delta (kq)N_{}(q)`$, we obtain for the gap equation
$`\stackrel{~}{\sigma }`$ $`=\lambda {\displaystyle \frac{dk_\eta }{2\pi }(1N_+(k)N_{}(k))R_k(u)}`$ (12)
where $`R_k(u)=12k_\eta ^2|f_k^+(u)|^2.`$ and
$`\lambda ^1={\displaystyle \frac{dk_\eta }{2\pi }\frac{1}{\sqrt{k_\eta ^2+\stackrel{~}{m}_f^2}}}={\displaystyle \frac{dk}{2\pi }\frac{1}{\sqrt{k^2+m_f^2}}}.`$
This equation is solved simultaneously with eq. (10).
We take our initial state to be in local equilibrium so that $`N_\pm (k,\mu ,T)=[e^{(\omega _k(0)\mu )/T}+1]^1`$ where $`\omega _k(0)=E`$ $`=\sqrt{k^2+\sigma ^2(0)}`$ $`=\frac{\stackrel{~}{\omega }_k(0)}{\tau _0}.`$ Since we start our simulation in the unbroken mode, $`\stackrel{~}{\sigma }(0)=0.`$ We choose the initial $`\tau _0=\frac{1}{m_f}`$ and measure the proper time in these units. We use adiabatic initial conditions on the mode functions $`f`$, i.e. $`f_k(0)=\frac{N_k}{\sqrt{2\stackrel{~}{\omega }_k}}`$, $`\dot{f}_k^+(0)=i\stackrel{~}{\omega }_kf_k^+(0)`$ and $`N_k^2=[\stackrel{~}{\omega }_k(0)+\stackrel{~}{\sigma }(0)]^1.`$ To obtain non-trivial dynamics in this mean field approximation at high temperatures, it is necessary to explicitly break the chiral symmetry by giving $`\dot{\stackrel{~}{\sigma }}`$ a small initial value which we choose to be $`\dot{\stackrel{~}{\sigma }}(0)=10^3`$.
We have studied three separate starting points on the phase diagram of Fig. 1 in our numerical simulations. We determined the energy density and the pressure from the expectation value of the energy momentum tensor as described in . In the $`\eta `$, $`\tau `$ coordinate system $`T_{\mu \nu }`$ is diagonal which allows us to read off the comoving pressure and energy density. After renormalization we obtain
$`ϵ(\tau )\tau ^2={\displaystyle _0^{\stackrel{~}{\mathrm{\Lambda }}}}{\displaystyle \frac{dk_\eta }{2\pi }}[{\displaystyle \frac{\stackrel{~}{\sigma }^2}{\sqrt{k_\eta ^2+\stackrel{~}{m}_f^2}}}+4\mathrm{\Omega }_k(\stackrel{~}{\sigma }^2\omega _k^2)|f_k|^2.`$ (13)
$`(N_++N_{}).[2\stackrel{~}{\sigma }+4\mathrm{\Omega }_k(\omega _k^2\stackrel{~}{\sigma }^2)|f_k|^2+2(k_\eta \stackrel{~}{\sigma })]],`$ (14)
$`p\tau ^2`$ $`={\displaystyle _0^{\stackrel{~}{\mathrm{\Lambda }}}}{\displaystyle \frac{dk_\eta }{2\pi }}[(1N_+N_{})4(\stackrel{~}{\sigma }+\mathrm{\Omega }_k)(\stackrel{~}{\sigma }^2\omega _k^2)|f_k|^2.`$ (16)
$`+2{\displaystyle \frac{k_\eta ^2}{\sqrt{k_\eta ^2+\stackrel{~}{\sigma }^2}}}+2\sqrt{k_\eta ^2+\stackrel{~}{\sigma }^2}2k_\eta {\displaystyle \frac{\sigma ^2}{\sqrt{k_\eta ^2+\stackrel{~}{m}_f^2}}}].`$
The integrations involve a moving cutoff $`\stackrel{~}{\mathrm{\Lambda }}=\mathrm{\Lambda }\tau `$ when the mode functions are truncated at physical $`k_z=\mathrm{\Lambda }`$. In the massless phase, one finds that the exact equation of state is $`p=ϵ.`$ To compare our field theory calculation with a local equilibrium hydrodynamical model we assume
$$T^{\alpha \beta }=pg^{\alpha \beta }+(ϵ+p)u^\alpha u^\beta $$
(17)
The conservation law of energy and momentum $`T^{\alpha \beta }{}_{;\beta }{}^{}=0,`$ combined with scaling law $`v=z/t`$ and $`p=ϵ`$ yields $`\frac{ϵ}{ϵ_0}=(\frac{\tau _0}{\tau })^2`$,$`\frac{T}{T_0}=(\frac{\tau _0}{\tau })`$. From Eq. (14) we can also determine $`p(\mu ,T)`$ and $`ϵ(\mu ,T)`$. Assuming $`T/T_0=\tau _0/\tau `$ and $`\mu /\mu _0=\tau _0/\tau `$ we find that the local equilibrium expressions for $`ϵ`$ and $`p`$ evolve identically to the numerically determined field theory evolution before the phase transition. (We note that in thermodynamic equilibrium $`dT/T=d\mu /\mu `$ and so close to equilibrium we expect the temperature and chemical potential to have a similar falloff with time. In fact, a different falloff for the two quantities as a function of time can be viewed as a departure from local thermal and chemical equilibrium.) With the same assumptions we find the distributions for $`N_\pm `$ plotted against $`k_\eta `$ are independent of $`\tau `$. This also agrees with the exact evolution before the phase transition.
We want to understand how the particle number distributions evolve in time. In relativistic quantum mechanics, particle number is not conserved. However in a mean field approximation one can define an interpolating number operator which at late times becomes the outstate number operator. By fitting the interpolating number densities for both fermions and antifermions to Fermi-Dirac distributions we extract the best value of $`\mu `$ and $`T`$ for that value of the proper time. To define the interpolating number operator we use a set of orthonormal mode functions $`y_k`$ which are the adiabatic approximation to the exact mode functions: $`y_k^+=u_ke^{i{\scriptscriptstyle \stackrel{~}{\omega }_k𝑑u}}`$; $`y_k^{}=v_ke^{i{\scriptscriptstyle \stackrel{~}{\omega }_k𝑑u}}`$ with $`u_k=\frac{i\gamma ^\mu k_\mu +\stackrel{~}{\sigma }}{\sqrt{2\stackrel{~}{\omega }_k(\stackrel{~}{\omega }_k+\stackrel{~}{\sigma })}}\chi ^+`$; $`v_k=\frac{i\gamma ^\mu k_\mu +\stackrel{~}{\sigma }}{\sqrt{2\stackrel{~}{\omega }_k(\stackrel{~}{\omega }_k+\stackrel{~}{\sigma })}}\chi ^{}.`$ The creation and annihilation operators then become time dependent and the expansion of the quantum field becomes
$`\mathrm{\Phi }(x)={\displaystyle \frac{dk_\eta }{2\pi }[a(k,u)y_k^+(u)+c^{}(k,u)y_k^{}(u)]e^{ik_\eta \eta }}.`$
This is an alternative expansion to that found in Eq.(7). The two sets of creation and annihilation operators are related by a Bogoliubov transformation $`a(k,u)=\alpha _k(u)b(k)`$ +$`\beta _k^{}d^{}(k)`$; $`c^{}(k,u)=`$\- $`\beta _k(u)b(k)`$\+ $`\alpha _k^{}d^{}(k).`$ To ensure that at $`u=0`$ the two number operators match, one chooses adiabatic initial conditions: $`\varphi _k=y_k`$, so that $`\alpha _k(0)=1`$;$`\beta _k(0)=0.`$ The interpolating number operators for fermions and anti-fermions are defined by $`N^+(k,u)=a^{}(k,u)a(k,u);`$ $`N^{}(k,u)=c^{}(k,u)c(k,u).`$ With $`\mathrm{\Delta }_k=\frac{\dot{\stackrel{~}{\mathrm{\Omega }}}_k+\dot{\stackrel{~}{\sigma }}}{2\stackrel{~}{\mathrm{\Omega }}_k}`$ we have explicitly
$`|\beta _k|^2=k_\eta ^2{\displaystyle \frac{(\stackrel{~}{\mathrm{\Omega }}_k\stackrel{~}{\omega }_k)^2+\mathrm{\Delta }_k^2}{2\stackrel{~}{\omega }_k(\stackrel{~}{\omega }_k+\stackrel{~}{\sigma })[\stackrel{~}{\mathrm{\Omega }}_k^2+\stackrel{~}{\omega }_k^2+2\stackrel{~}{\mathrm{\Omega }}_k\stackrel{~}{\sigma }+\mathrm{\Delta }_k^2]}},`$
$`N^\pm (k,u)=N^\pm (k)+[1N^+(k)N^{}(k)]|\beta _k(u)|^2.`$
We have solved the simultaneous equations Eqs. (10, 12) numerically. Comparing $`N^\pm (k,u)`$ with an equilibrium parameterization we have determined $`T(k,u)`$ and $`\mu (k,u)`$ as a function of $`k`$. When these quantities are independent of $`k_\eta =k\tau `$ this defines a time evolving temperature and chemical potential. We found that $`T`$ and $`\mu `$ are independent of $`k`$ except at high momentum before the chiral phase transition.
From Fig. 2 we see that for both the 1st and 2nd order transitions, $`\sigma (\tau )`$ shows a sharp transition during evolution from the unbroken mode to the broken symmetry mode. Before the phase transition the temperature falls consistent with the equation of state $`p=ϵ`$. For the 2nd order transition, the chemical potential follows the temperature and falls as $`\tau ^1`$. After the phase transition, there is now a mass scale $`m_f`$ which leads to oscillations of $`\sigma `$. For the 1st order transition the chemical potential falls faster than $`\tau ^1`$. If one traverses the tricritical regime one finds results for $`\mu `$ intermediate between the two cases displayed.
The order of the transition has a more noticeable effect on the spectrum of particles and antiparticles. If the system evolves in local thermal equilibrium with $`\sigma =0`$, then when $`N^\pm (k,u)`$ is plotted vs. $`k_\eta =k\tau `$ it is independent of $`u`$. A deviation from this result is an indication of the system going out of equilibrium. We find because of the “latent heat” released during a first order transition that the distortion of the spectrum is greatest in that case. (see Fig. 3). If one traverses the tricritical regime one finds results intermediate between the two cases displayed.
In local equilibrium with $`\sigma =0`$, $`ϵ=p\tau ^2`$ . Simulations, shown in Fig. 4 agree with this before the phase transition occurs. After the phase transition we find that the energy density oscillates around the true broken symmetry value discussed earlier, namely $`ϵ_0=1/4\pi `$. These oscillations would be damped if we included hard scattering effects . The details of this calculation as well as a discussion of correlation functions and the effects of a bare mass will be presented elsewhere.
We would like to thank Emil Mottola, Salman Habib and Dan Boyanovsky for discussions.
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# Untitled Document
Resolution of a long standing discrepancy in R with spin zero quarks
S. Pacetti and Y. Srivastava
1. Dipartimento di Fisica & INFN, Università di Perugia, Perugia, Italy
2. Physics Department, Northeastern University, Boston, MASS, USA
Abstract
A previously successful dispersive method has been applied to understand different values for $`R(\sqrt{s}=5÷7.5GeV)`$ obtained by MARK I and Crystal Ball collaborations. We compute $`R`$ in the disputed region with data from outside this region and asymptotic behavior given by the standard model with 5 quark flavors, and find agreement with the Crystal Ball result. On the other hand, the MARK I data are reproduced if we augment the asymptotic behavior with contributions from a single spin zero quark of charge ($`1/3`$). The visible hadronic fragments from such scalar quarks are not likely to produce predominantly pure 2-jet events at such low energies. Hence, such decay modes may have been removed by the Crystal Ball energy imbalance cuts in their definition of hadronic events but not in MARK I events, thus accounting for the discrepancy in the two results. Upper bounds on spin zero quark production at LEP through $`Z`$ decay data are used to estimate the mixing angle between $`T_3=1/2`$ and $`T_3=0`$ scalar quarks. Recent negative results about spin zero quarks from CLEO are critically examined. We briefly discuss diquark production hypothesis and find it very unlikely to explain the discrepancy.
The purpose of this work is to shed some light on an ancient and as yet unresolved discrepancy in the measured values of $`R(s)`$ in the energy range $`\sqrt{s}=5÷7.5GeV`$, both measured at SLAC. The MARK I & II collaborations found a higher value $`R_{av}=4.3\pm 0.4`$, whereas the Crystal Ball collaboration found a lower value $`R_{av}=3.44\pm 0.03\pm 0.18`$. There are a few data points in this region from PLUTO and DASP collaborations as well which are in agreement with the MARK I collaboration.
We employ a dispersion-theoretic formalism which has been quite successful in studies of elastic form factors of nucleons and mesons \[6-8\]. For example, using space-like and above the $`N\overline{N}`$ threshold time-like data for the nucleon form factor, we were able to compute the same elastic nucleon form factor in the experimentally inaccessible region between the two pion and two nucleon thresholds, in a model independent way without any bias towards expected resonances. Remarkably, resonance structures with peaks for the $`\rho (770)`$, $`\rho ^{}(1600)`$ and a structure near the $`N\overline{N}`$ were automatically generated. For details about the regularization scheme employed to solve the integral equations, we refer the reader to references \[6-8\].
For the problem at hand, we consider the (Lorentz scalar) vacuum polarization function $`\mathrm{\Pi }_\gamma (s)`$ which is defined via the EM polarization tensor $`\mathrm{\Pi }_{\mu \nu }(s)`$ as:
$`ie^2{\displaystyle d^4xe^{iqx}0\left|Tj_{em}^\mu \left(x\right)j_{em}^\nu \left(0\right)\right|0}=\left(q^2g^{\mu \nu }q^\mu q^\nu \right)\mathrm{\Pi }_\gamma \left(q^2\right),`$
where $`j_{em}^\mu (x)`$ is the electro-magnetic current. $`\mathrm{\Pi }_\gamma (s)`$ is an analytic function in the complex $`s`$ plane with a branch cut on the real positive axis for $`s>s_0`$, where $`s_0=4m_\pi ^2`$. We use the dispersion relation:
$`\mathrm{\Pi }_\gamma \left(t\right)\mathrm{\Pi }_\gamma \left(0\right)={\displaystyle \frac{t}{\pi }}{\displaystyle _{s_0}^{\mathrm{}}}𝑑s{\displaystyle \frac{Im\mathrm{\Pi }_\gamma \left(s\right)}{s\left(stiϵ\right)}},`$ (1)
subtracted at $`t=0`$, with the renormalized photon vacuum polarization: $`\mathrm{\Pi }(t)=\mathrm{\Pi }_\gamma (t)\mathrm{\Pi }_\gamma (0)`$. The optical theorem relates the imaginary part of the vacuum polarization $`\mathrm{\Pi }(s)`$ to the function $`R(s)`$ by means of the identity:
$`Im\mathrm{\Pi }\left(s\right)=\left({\displaystyle \frac{\alpha }{3}}\right)R\left(s\right).`$ (2)
One obtains, therefore, the dispersion integral :
$`\mathrm{\Pi }\left(t\right)={\displaystyle \frac{t\alpha }{3\pi }}{\displaystyle _{s_0}^{\mathrm{}}}𝑑s{\displaystyle \frac{R\left(s\right)}{s\left(st\right)}},`$ (3)
through which $`\mathrm{\Pi }`$ for space-like $`t`$ gets related to $`R`$ for time-like $`s`$.
The idea is to consider this relation as an integral equation to compute $`R(s)`$ in the interval $`[5÷7.5GeV]`$, using the following input: (i) experimental data in the time-like region from threshold up to $`5GeV`$ and from $`7.5÷100GeV`$, (ii) PQCD asymptotic behavior for the rest of the time like region and (iii) in the space-like region only the asymptotic value of $`\mathrm{\Pi }(t)`$ calculated for values of $`t`$ in the interval $`[205GeV^2\overline{t}200GeV^2]`$ through PQCD. Then we have:
$`\mathrm{\Pi }\left(t\right)I_{01}\left(t\right)I_{23}\left(t\right)I_a\left(t\right)={\displaystyle \frac{t\alpha }{3\pi }}{\displaystyle _{s_1}^{s_2}}𝑑s{\displaystyle \frac{R\left(s\right)}{s\left(st\right)}}`$ (4)
where:
$`I_{01}={\displaystyle \frac{t\alpha }{3\pi }}{\displaystyle _{s_0}^{s_1}}𝑑s{\displaystyle \frac{R_{exp}\left(s\right)}{s\left(st\right)}},I_{23}={\displaystyle \frac{t\alpha }{3\pi }}{\displaystyle _{s_2}^{s_3}}𝑑s{\displaystyle \frac{R_{exp}\left(s\right)}{s\left(st\right)}},I_a\left(t\right)={\displaystyle \frac{t\alpha }{3\pi }}{\displaystyle _{s_3}^{\mathrm{}}}𝑑s{\displaystyle \frac{R_a\left(s\right)}{s\left(st\right)}},`$ (5)
with: $`s_1=(5GeV)^2`$, $`s_2=(7.5GeV)^2`$ and $`s_3=(100GeV)^2`$. The function $`R_{exp}(s)`$ used in the first and second integral (5) is a fit of the experimental data, the function $`R_a(s)`$ in the third integral is the PQCD asymptotic behavior at leading order:
$`R_a=N_c{\displaystyle \underset{q}{}}Q_q^2\left[1+{\displaystyle \frac{\alpha _s}{\pi }}\right],`$ (6)
where $`N_c`$ is the colour factor and $`Q_q`$ is the charge of the quark $`q`$. We note that higher order QCD correction terms (that is, beyond the $`\alpha _s/\pi `$ term) are indeed “corrections to a correction” and thus play a negligible role in our analysis of the discrepancy.
The five light quark ($`u,d,s,c,b`$) contributions to $`\mathrm{\Pi }(t)`$ can be safely calculated only at large $`t`$, since at low energies the quark interactions are modified considerably by strong interactions. At high energies, by virtue of asymptotic freedom inherent in QCD, we can treat the quark contribution similar to that for the leptons and, at leading order, we obtain:
$`\mathrm{\Pi }\left(t\right)=N_c{\displaystyle \frac{\alpha }{3\pi }}{\displaystyle \underset{q}{}}Q_q^2\left[\mathrm{ln}\left({\displaystyle \frac{t}{m_q^2}}\right){\displaystyle \frac{5}{3}}+O\left({\displaystyle \frac{m_q^2}{t}}\right)\right].`$ (7)
The difficulty in using this function comes from the masses, which is particularly severe for the light quarks $`u`$, $`d`$ and $`s`$, that are not unambiguously defined. To avoid this problem, but without any further approximation, we consider the derivative in $`t`$ of the dispersion integral (3) so as to obtain a new integral equation:
$`{\displaystyle \frac{d\mathrm{\Pi }\left(t\right)}{dt}}{\displaystyle \frac{dI_{01}\left(t\right)}{dt}}{\displaystyle \frac{dI\left(t\right)}{dt}}{\displaystyle \frac{dI_a\left(t\right)}{dt}}={\displaystyle \frac{\alpha }{3\pi }}{\displaystyle _{s_1}^{s_2}}𝑑s{\displaystyle \frac{R\left(s\right)}{\left(st\right)^2}}.`$ (8)
The asymptotic behavior of $`\frac{d\mathrm{\Pi }(t)}{dt}`$ for large $`t`$ is:
$`\left[{\displaystyle \frac{d\mathrm{\Pi }\left(t\right)}{dt}}\right]_A=N_c{\displaystyle \frac{\alpha }{3\pi }}{\displaystyle \underset{q}{}}Q_q^2{\displaystyle \frac{1}{t}}.`$ (9)
It has not just the virtue of being independent of the quark masses: it requires only a knowledge of the quark charges and the colour factor. Upon solving the integral equation (8) \[6-8\], we can calculate the values of $`R(s)`$ in the disputed region $`[s_1,s_2]`$ using as input the time-like experimental data and the quark charges.
Including only the five standard (spin 1/2) light quarks $`u,d,s,c,b`$, as shown in Fig.(1a), we find a value for $`R(s)`$ in good agreement with the Crystal Ball data .
But, if we include an additional contribution from a single species of spin zero quark of charge (-1/3), which requires simply adding in the summations (6) and (9) another charge (-1/3) squared, multiplied by a factor 1/4 (due to zero spin of the quark), we reproduce the MARK I data, shown in Fig.(1b).
In Fig.(2a) and Fig.(2b), we show a similar comparison but imposing continuity at the boundaries.
As a consistency check, we write a dispersion relation for $`\mathrm{\Delta }R(s)`$ in the above region whose asymptotic value contains contribution from one species of colour triplet, charge $`(1/3)`$, spin zero quark only. Phenomenologically, in the range $`s_3ss_2`$, the standard 5 quark model gives a rather good description of the data. Hence, we obtain the approximate relation
$`{\displaystyle _{s_1}^{s_2}}{\displaystyle \frac{\mathrm{\Delta }R\left(s\right)ds}{\left(s\overline{t}\right)^2}}{\displaystyle \frac{1}{12\overline{t}}}\left[1+{\displaystyle \frac{\overline{t}}{\left(s_3\overline{t}\right)}}\right]`$ (10)
This gives an average value $`\overline{\mathrm{\Delta }R}0.75`$. By way of comparison, the difference $`\mathrm{\Delta }R(s)`$ between the two curves from Fig.(1) are plotted in Fig.(3a). Fig.(3b) shows the difference $`\mathrm{\Delta }R(s)`$ for Fig.(2).
These average values are mutually consistent. Such a “locally averaged” $`\overline{\mathrm{\Delta }R}`$ should not be confused with $`\mathrm{\Delta }R_{asymptotic}`$. The above has been obtained under the assumption that beyond the disputed region, experimental data are adequately described by the standard 5 spin $`1/2`$ quark model. Thus, roughly $`(\mathrm{\Delta }R_{asy}\times (\overline{t})(\overline{\mathrm{\Delta }R})\times (s_2s_1))`$.
The two experiments can be reconciled only under the hypothesis that Crystal Ball experimental cuts may have eliminated signals due to the spin zero quarks. The bound states and resonances with $`J^{PC}=1^{}`$ coupling to $`e^+e^{}`$ are in relative $`P`$\- wave for spin zero quarks. They would be relatively close in energy and more importantly, are expected to decay copiously into a low energy photon ($`E_\gamma 300MeV`$) recoiling against $`S`$\- wave states which would subsequently decay into hadrons via 2 gluons . Such radiative decay events would have a very asymmetric energy pattern. For the model discussed below, a similar pattern may also follow (for the non resonant) scalar quark production. In their definition of hadronic events contributing to $`R`$, Crystal Ball imposed the following kinematic cuts for the energy imbalance between left-right, top-bottom and front-back hemispheres. If any of these fractional energy differences, $`A_{leftright}`$, $`A_{topbottom}`$ or $`A_{frontback}`$ was less than $`0.8`$, the event was classified as being due to beam-gas interactions, $`\gamma \gamma `$ collisions or large missing energy $`\tau `$ decays. Thus, such events were not included in $`R`$ by Crystal Ball.
In Fig.(4), we show a plot of the ratio between $`N_{bg}`$, which is the number of the background events from beam-gas and beam-wall interactions, and the total collisions $`N_{coll}`$ passing the hadron selection criteria, versus $`\sqrt{s}`$ from Table II of ref.(3). That the rejected events have peaks and structures resembling the MARK I data appears significant and supports our hypothesis.
Now we turn to a rough estimate of the mass of such a scalar quark (Y). Taking our cue from the steep rise in the CB rejected events between $`7.2÷7.4GeV`$, we may estimate $`m_Y3.6÷3.7GeV`$ from the production threshold for a pair ($`Y^{}Y`$) of such quarks. The local rise in MARK I data, prior to this threshold, beginning around $`5GeV`$ would however require that a single scalar quark be produced along with a couple of other light (spin $`1/2`$ $`u,d,s,c`$) quarks. One possibility (out of many others) would be a charge $`(1/3)`$ color triplet scalar quark carrying baryon number $`(2/3)`$ to coincide with the quantum numbers of a standard $`\overline{U}\overline{D}`$ (where $`U=u,c`$ and $`D=d,s`$) quark pair. Then, the final state $`YUD`$ would be allowed in $`e^+e^{}`$ reaction, with a threshold around $`5GeV`$. We expect no sharp structures associated with it and the level of production should decrease rapidly with $`s`$ so as to reach its asymptotic parton level. Two points are worthy of note here. First, such a scalar quark would be an “elementary” particle and not a composite “diquark” state of two standard $`\overline{U}`$ and $`\overline{D}`$ quarks. (For this case, there might be non trivial mixing and interference terms). While for the narrow window of pair produced resonances ($`Y^{}Y`$ states), non-relativistic potential models may be a reasonable rough guide, the production dynamics and level of cross section for multi-particle relativistic systems (with exotic quantum numbers such as $`YUD`$) are not easy to compute or estimate directly since non relativistic potential models can not be employed for this purpose. The hadronization of scalar quarks would also be very different from those for the other quarks. Not knowing the nature of the beast, our modest aim here has been to estimate these contributions to R through unitarity, dispersion relations and asymptotic behavior which are consistent with data from outside this region. Secondly, if one were to study the mass distribution of an $`e^+e^{}`$ or $`\mu ^+\mu ^{}`$ in the final state produced say from a $`p\overline{p}`$ initial state, we would expect to see an enhancement visible only around $`7.2÷7.4GeV`$ (due to pair production of scalar quarks) since the probability to find $`YUD`$ in the initial state would be negligible.
An alternative explanation of this rather substantial difference between MARK I and CB data has been offered by Eidelman and Jegerlehner . According to them, this difference arises from QED and QCD radiative corrections. In the final analysis, our explanation of this difference also involves radiative events but its source is different. Ours is generated from an extra scalar quark and MARK I data are “physical”, whereas in the other explanation MARK I data should be divided (“corrected”) by the radiative effects. It should be firmly kept in mind that the ($`g2`$) results discussed in ref.(12) can not be used to discriminate between MARK I and CB data since the results are hardly changed whether one or the other data are used.
The non resonant part of events from scalar quark decays would also not be of the 2-jet type at low energies. Both in the ALEPH and the CLEO collaboration searches for scalar quarks, it has been assumed that each scalar quark decays weakly, viz., into $`c\mu ^{}\stackrel{~}{\nu }_s`$ and $`ce^{}\stackrel{~}{\nu }_s`$, where $`\stackrel{~}{\nu }_s`$ is a scalar neutrino. For $`\sqrt{s}=5÷7.5GeV`$ and scalar quark mass in the $`3÷4GeV`$ range, the decay events would not be of the 2-jet type. On the other hand, at much higher energies, for example $`\sqrt{s}=160÷205GeV`$ as in the ALEPH data, the decay products would be confined to the forward and backward directions similar to a 2-jet profile.
Evidence supporting a low mass spin zero quark must now be confronted with the extremely accurate decay systematics for the $`Z^o`$ for various channels available from LEP \[15-18\]. A computation of the coupling of observable spin zero quarks to the $`Z`$ boson requires a knowledge of their weak iso-spin. Let $`Y_L`$ and $`Y_R`$ denote spin zero quarks with $`T_3=1/2`$ and $`T_3=0`$ respectively. After mixing, let $`Y_{}`$ and $`Y_+`$ denote the low and high mass eigenstates. If the higher of these masses happens to be larger than the $`Z^o`$ mass, kinematically $`Z^o`$ can decay only into $`Y_{}^{}Y_{}`$, with the decay amplitude given by
$`Amplitude\left(Z^o\left(P\right)Y_{}^{}\left(p_1\right)Y_{}\left(p_2\right)\right)=\left({\displaystyle \frac{e}{sin\vartheta _Wcos\vartheta _W}}\right)Kϵ_\mu \left(P\right)\left(p_1p_2\right)^\mu ,`$ (11)
where $`K`$ is given by
$`K=\left({\displaystyle \frac{1}{3}}sin^2\vartheta _W\right)\left|U_R\right|^2+\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{3}}sin^2\vartheta _W\right)\left|U_L\right|^2.`$ (12)
Here $`U_R=cos\delta `$ and $`U_L=sin\delta `$ are the mixing matrix elements to the lower mass state ($``$). Thus, the branching ratio of $`Z^o`$ into a pair of low mass spin zero quarks can be arbitrarily small, even zero for $`\delta ^{}`$
$`\delta ^{}=arcsin\left(\sqrt{{\displaystyle \frac{2sin^2\vartheta _W}{3}}}\right)23^o.`$ (13)
Present measurements of the $`Z^o`$ width and observed branching ratios therefore can serve to place limits on the mixing angle $`\delta `$. For example, if we assume that the branching ratio for the spin zero quark decay channel is less than one part per mille, approximately $`\delta =(23\pm 15)^o`$ (for $`\delta `$ chosen in the first quadrant).
The recent CLEO experiment has given quite stringent bounds regarding a low mass squark partner of the $`b`$ quark. They have looked for and not found an expected high level $`D`$ and $`D^{}`$ signal. Of course, it could be that the low mass spin zero quark is not the partner of the $`b`$ quark or even if it were so, the squark mixing matrix may not be similar to the quark mixing matrix. Also, the estimate of the weak decay signal expected in the CLEO experiment is based on phase space, i.e., a structureless matrix element and an essentially zero mass scalar neutrino. Neither of these hypotheses can be justified theoretically. For example, even a $`300MeV`$ neutrino would lower the signal to practically its background value. We are not aware of any argument which can so restrict the scalar neutrino mass.
Among other presently available beams, high luminosity asymmetric B-factories offer a good avenue for new and independent check of data on $`R`$ in the $`5÷8GeV`$ region, by observing a bremsstrahlung photon of energy $`1.3÷3.8GeV`$. An independent check of scalar quark production at B-factories would be through their weak decays by looking at changes in the number of opposite sign leptons (in the opposite hemispheres) as the beam energy passes through the $`b\overline{b}`$ threshold. Signals from spin zero quarks should be looked for also in $`\mu ^+\mu ^{}`$ production and in the hadronic machines via the gluon induced process $`g+\overline{D}Y+U`$ (for the mechanism discussed earlier). In addition, the next Tevatron run offers the possibility of looking for $`Y`$ jets as well as a means to probe for $`\stackrel{~}{w}`$ (any generic spin $`1/2`$ object recoiling against the $`Y`$) up-to a mass of about $`170GeV`$ through the top quark decay $`tY+\stackrel{~}{w}^+`$.
We now discuss a completely different hypothesis to resolve the discrepancy. If one accepts that the difference between the two experiments is due to one experiment being essentially sensitive only to 2-jet events and the other not so biased, this may be accounted for through the production of a diquark-antidiquark pair, since they would produce predominantly 4 quark final states. As the diquark photon vertex includes a form factor their production would disappear at high energies, e.g., at LEP.
A diquark diagram (inclusive of a form factor) to the self energy of the photon is shown in Fig.(5).
It contributes to the imaginary part
$`\text{Im}\left(\pi _{DQ}\left(s\right)\right)={\displaystyle \frac{1}{4}}{\displaystyle \frac{N_c\alpha Q^2}{3}}\left|F\left(s\right)\right|^2\left(1{\displaystyle \frac{4m^2}{s}}\right)^{\frac{3}{2}}`$ (14)
and thus to the dispersion relation
$`\widehat{\pi }_{DQ}\left(t\right)={\displaystyle \frac{1}{4}}{\displaystyle \frac{N_c\alpha Q^2}{3}}{\displaystyle \frac{t}{\pi }}{\displaystyle _{4m^2}^{\mathrm{}}}\left(1{\displaystyle \frac{4m^2}{s}}\right)^{\frac{3}{2}}{\displaystyle \frac{\left|F\left(s\right)\right|^2}{s\left(st\right)}}𝑑s.`$ (15)
$`F\left(s\right)`$ is normalized to $`1`$ at $`s=0`$ and is expected to go asymptotically as
$`\left|F\left(s\right)\right|{\displaystyle \frac{Q_0^2}{s}},`$ (16)
with $`Q_0^210GeV^2`$ . For a diquark of charge $`q`$ to saturate the discrepancy in $`\overline{\mathrm{\Delta }R}=0.75`$ (as discussed previously), we require
$`\left|F\left(s\right)\right|=\{\begin{array}{cc}\frac{3}{N_cq}\hfill & \text{if }s[s_1,s_2]\hfill \\ N_cq\frac{Q_0^2}{s}\hfill & \text{if }s>s_2\hfill \end{array},`$ (19)
to be fed in
$`{\displaystyle \frac{d\widehat{\pi }_{DQ}\left(t\right)}{dt}}={\displaystyle \frac{\alpha }{3\pi }}{\displaystyle \frac{1}{12}}{\displaystyle _{s_1}^{\mathrm{}}}\left(1{\displaystyle \frac{s_1}{s}}\right)^{\frac{3}{2}}{\displaystyle \frac{\left|F\left(s\right)\right|^2}{\left(st\right)^2}}𝑑s.`$ (20)
Unfortunately, the contribution of such diquarks in the space like region falls short by a factor $`4`$. Hence, the hypothesis of diquarks being responsible for the discrepancy does not appear to be internally self consistent.
In conclusion, our dispersive analysis of the MARK I data (the only set included in PDG, the particle data tables ) appears to require the existence of a charge $`(1/3)`$ scalar quark with a low mass. While the level of the reported cross section appears to be consistent with our analysis, it remains an open dynamical problem how to generate such a large value. Preliminary estimates of contributions to $`R`$ from bound states computed through static potentials appear to fall short by a factor of 3 or more. This problem obviously deserves further study and the expected data from BABAR on a new measurement of $`R`$ in the disputed region suggested here would be most useful. On the other hand, lack of evidence in the present LEP data about the direct production of such low mass scalar quarks in $`Z^o`$ decay only serves to indicate that it is predominantly an iso-singlet. On the other hand, a $`b`$ jet from the $`Z^o`$ decay may sequentially produce such a scalar quark whose signature might be revealed through an analysis of the final $`\mu ^+\mu ^{}`$ mass distribution in the $`7.2÷7.4GeV`$ region.
It is with much gratitude that we thank P. Giromini for directing our attention to this problem and for fruitful discussions at every stage of this work. It is a pleasure to thank S. Glashow, U. Heintz, K. Lane, M. Narain, G. Pancheri, S. Reucroft, J. Swain and A. Widom for much encouragement and many helpful suggestions.
Footnotes and References
1. J. L. Siegrist et al. MARK I Coll., Phys. Rev. D26 (1982) 969.
2. M. W. Coles et al. MARK II Coll., Phys. Rev. D26 (1982) 2190.
3. C. Edwards et al. Crystal Ball Coll., SLAC-PUB-5160 (1990).
4. J. Burmeister et al. PLUTO Coll., Phys. Lett. 66B (1977) 395; Ch. Berger et al., Phys. Lett. 81B (1979) 410.
5. R. Brandelik et al., DASP Coll. it Phys. Lett. 76B (1978) 361; H. Albrecht et al, Phys. Lett. 116B (1982) 383.
6. R. Baldini, S. Dubnička, P. Gauzzi, E. Pasqualucci, S. Pacetti and Y. Srivastava, Euro. Phys. J. C11 (1999) 709.
7. R. Baldini, S. Dubnička, P. Gauzzi, E. Pasqualucci, S. Pacetti and Y. Srivastava, Nuc. Phys. A666&667 (2000) 38c.
8. R. Baldini, S. Dubnička, P. Gauzzi, E. Pasqualucci, S. Pacetti and Y. Srivastava, Frascati Report LNF-98/024(IR).
9. N. Cabbibo and R. Gatto, Phys. Rev. Lett. 4 (1960) 313; Phys. Rev. 124 (1961) 1577.
10. R. M. Barnett, M. Dine and L. McLerran, Phys. Rev. D22 (1980) 594.
11. C. R. Nappi Phys. Rev. D25 (1982) 84.
12. S. Eidelman and F. Jegerlehner, Z. Phys. C67 (1995) 585.
13. G. Taylor, LEPC Presentation, 20 July 2000, http://alephwww.cern.ch.
14. V. Savinov et al. CLEO Coll., hep-ex/0010047.
15. P. Abreu et al. DELPHI Coll. Nucl. Phys. B418 (1994) 403.
16. M. Acciarri et al. L3 Coll. Z. Phys. C62 (1994) 551.
17. R. Akers et al. OPAL Coll. Z. Phys. C61 (1994) 199.
18. D. Buskulic et al. ALEPH Coll. Z. Phys. C62 (1994) 539.
19. M. Anselmino, E. Predazzi, S. Ekelin, S. Fredriksson and D.B. Lichtenberg, Rev. Mod. Phys. 65 (1993) 1199.
20. Particle Data Group, Review of Particle Physics, Eur. Phys. J. C3 (1998) 227.
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# Optimal distillation of a GHZ state
## Abstract
We present the optimal local protocol to distill a Greenberger-Horne-Zeilinger (GHZ) state from a single copy of any pure state of three qubits.
Because of its relevance in quantum information theory, entanglement has been attracting considerable attention in recent years. The initial efforts, mainly devoted to acquire a quantitative description of bipartite quantum correlations, led to the identification of the measures governing pure-state entanglement both in the asymptotic and non-asymptotic regimes. As a result, it is nowadays known how a bipartite system prepared in a pure state can be optimally manipulated under the restriction that only local operations on the subsystems aided with classical communication (LOCC) are allowed.
Once a relatively complete command on pure-state entanglement had been achieved for bipartite systems, efforts have moved quite recently to address the somehow more intrincate case of tripartite entanglement, for which a system consisting of three two-level subsystems —that is, a three-qubit system— is the simplest scenario.
There are two distinct ways in which three qubits can be entangled , in the following sense. Let us identify as equivalent all pure states that can be reversibly interconverted, with some finite probability of success, when the parties are allowed to perform LOCC. Then, whereas all entangled pure states of two qubits are equivalent to the Einstein-Podolsky-Rosen (EPR) state , $`1/\sqrt{2}(|00+|11),`$ it turns out that truly tripartite pure-state entanglement of three qubits is either equivalent to the GHZ state ,
$$|GHZ\frac{1}{\sqrt{2}}(|000+|111),$$
(1)
or else to the W state , $`1/\sqrt{3}(|001+|010+|100),`$ with these two states being completely inaccessible from each other by means of LOCC.
Although it is genuinely tripartite, the entanglement of the $`W`$ state only maximizes two-qubit quantum correlations . As far as three-qubit correlations are concerned, the GHZ state appears as the maximally entangled state : it violates Bell inequalities maximally and it maximizes the mutual information of local measurements. It is also the only state from which an EPR state between any two chosen qubits can be obtained with certainty.
In this Letter we are concerned with the distillation of a GHZ state from an arbitrary pure state $`|\psi `$ of three qubits (see also ). We present the optimal strategy to distill a GHZ state starting from a single copy of $`|\psi `$. That is, we provide the local protocol that allows the three parties to transform the state $`|\psi `$ into a GHZ state with the greatest a priori probability of success compatible with LOCC. Marginally, we also discuss the best approximate transformation of $`|\psi `$ into $`|GHZ`$ by means of LOCC, which turns out to simply consist in a concatenation of local unitary transformations. Obviously, these results may be of practical interest in any situation where three parties wishing to have a GHZ state only share the state $`|\psi `$, and where for some reason —for instance because they are in sufficiently distant locations— the restriction to perform only LOCC is not merely academical. From a more general perspective, several difficulties have been faced which stem from characteristic features of entanglement involving more than two subsystems, such as the non-existence of the Schmidt decomposition (that is, from the fact that not all three-qubit states can be brought, by means of local unitary transformations, into the form $`\alpha |000+\beta |111`$, see ). In this sense, our work explores the nature of entanglement beyond the bipartite case, while providing useful tools for its study.
Let us first briefly review how, in the two-party setting, Alice and Bob can optimally distill a generic pure state of two qubits into an EPR state. The initial state can be transformed, by means of local unitary operations on Alice and Bob’s sides, into a reference state whose Schmidt decomposition is
$$\alpha |00+\beta |11,$$
(2)
where $`\alpha \beta 0`$ and $`\alpha ^2+\beta ^2=1`$. Then, following , the diagonal operator $`M\frac{\beta }{\alpha }|00|+|11|`$ —corresponding to a POVM on, say, Alice’s part— transforms state (2) into an EPR state with probability $`2\beta ^2`$. Finally, that such distillation protocol is optimal can be interpreted in terms of the entanglement monotone $`E_2\beta ^2`$ . Indeed, if a protocol would distill the EPR state with greater probability, this would contradict the non-increasing character of $`E_2`$ under LOCC.
Similar steps will be followed in order to find which is the optimal distillation protocol of a GHZ state in a three-qubit system: ($`i`$) the Schmidt decomposition will be replaced with a convenient two-term product decomposition; ($`ii`$) a subclass of distillation protocols, namely those in which each party is only allowed to perform one positive operator-valued measurement (POVM) and only one of the overall outcomes is a GHZ state, will be discussed and optimized; ($`iii`$) the optimal of such one succesful branch protocols (OSBP) will be proved to be the best distillation protocol by showing that the associated probability of success is an entanglement monotone .
Recall that a necessary and sufficient condition for the distillability of a GHZ from a three-qubit pure state $`|\psi `$ is that the three reduced density matrices $`\rho _A`$Tr$`{}_{BC}{}^{}|\psi \psi |`$, $`\rho _B`$ and $`\rho _C`$ have rank $`2`$, and that the range of $`\rho _{BC}`$Tr$`{}_{A}{}^{}|\psi \psi |`$ contains two product vectors, say $`|b_1c_1`$ and $`|b_2c_2`$ . Only in this case $`|\psi `$ admits a (unique) two-term product decomposition,
$$|\psi =\mu _1|a_1b_1c_1+\mu _2e^{i\phi }|a_2b_2c_2,\mu _1\mu _2>0,$$
(3)
where $`\phi [0,2\pi )`$ and the normalized local vectors $`|\kappa _1,|\kappa _2𝒞^2`$ satisfy $`1>\kappa _1|\kappa _20`$, $`\kappa =a,b,c`$. In its turn, a distillation protocol consists of a series of local POVMs, where the specific POVM to be performed at any time may be choosen depending on previous outcomes, the whole protocol having a tree structure. For a given branch of outcomes the initial state $`|\psi `$ is transformed into a final state
$$ABC|\psi ,$$
(4)
where the local operator $`A`$ collects contributions from all POVMs performed by Alice in that branch, and similarly for operators $`B`$ and $`C`$. Characterizing operator $`A`$ by how it transforms $`|a_1`$ and $`|a_2`$,
$$A=\alpha _1|a_1^{}\stackrel{~}{a}_1|+\alpha _2e^{i\phi _a}|a_2^{}\stackrel{~}{a}_2|,$$
(5)
where $`\alpha _1,\alpha _2,a_1^{}|a_2^{}`$$``$$`0`$ and $`\{|\stackrel{~}{a}_1,|\stackrel{~}{a}_2\}`$ is the biorthonormal basis to $`\{|a_1,|a_2\}`$, i.e. $`\stackrel{~}{a}_i|a_j=\delta _{ij}`$, we observe that the modifications introduced by A in $`|\psi `$, $`A1𝐥_B1𝐥_C|\psi =\alpha _1\mu _1|a_1^{}b_1c_1+\alpha _2\mu _2e^{i(\phi +\phi _a)}|a_2^{}b_2c_2,`$ concern the weights $`\mu _1`$ and $`\mu _2`$, the relative phase $`e^{i\phi }`$ and the local scalar product $`a_1|a_2`$, whereas the other two local scalar products, $`b_1|b_2`$ and $`c_1|c_2`$, remain unchanged. We can therefore readily conclude, from the uniqueness of the product decomposition (3) and the previous observation, that the only way for the resulting state (4) to be proportional to $`|GHZ`$ is that ($`i`$) each local operator transforms the corresponding couple of non-orthogonal, local states into an orthonormal pair, $`\{|\kappa _1,|\kappa _2\}\{|0,|1\},\kappa =a,b,c,`$ that is,
$`A`$ $`=`$ $`\alpha _1|0\stackrel{~}{a}_1|+\alpha _2e^{i\phi _a}|1\stackrel{~}{a}_2|,`$ (6)
$`B`$ $`=`$ $`\beta _1|0\stackrel{~}{b}_1|+\beta _2e^{i\phi _b}|1\stackrel{~}{b}_2|,`$ (7)
$`C`$ $`=`$ $`\gamma _1|0\stackrel{~}{c}_1|+\gamma _2e^{i\phi _c}|1\stackrel{~}{c}_2|,`$ (8)
and that ($`ii`$) their combined effect modifies the weights $`\mu _1`$ and $`\mu _2`$ to be equal, $`\mu _i\alpha _i\beta _i\gamma _i\mu _i=\sqrt{p/2},i=1,2,`$, where
$$p\psi |A^{}AB^{}BC^{}C|\psi =2(\alpha _1\beta _1\gamma _1\mu _1)^2$$
(9)
is the probability that the distillation protocol follows that particular, successful branch.
We have thus seen that the three parties must act on the system in order to distill a GHZ, because each of them must orthonormalize its local couple of states. The other relevant feature of the distillation process, namely making the relative weights equal, may be distributed in several ways among the parties.
Our OSBP for distillation consists in each of the parties performing a unique, two-outcome POVM, and it is built in such a way that after each POVM one of the two possible resulting states contains no three-partite entanglement at all, so that only one branch of the whole protocol succeeds in distilling a GHZ (see figure 1). In mathematical terms this implies that if e.g. $`A`$ in (7) is the successful operator in Alice’s POVM, then the other operator of the local POVM, $`\overline{A}`$, which satisfies $`A^{}A+\overline{A}^{}\overline{A}=1𝐥_A,`$ has rank equal to 1, so that the resulting state $`\overline{A}1𝐥_B1𝐥_C|\psi `$ is product in Alice’s subsystem. Expressing the identity operator as
$$1𝐥_A=\underset{i,j}{}a_i|a_j|\stackrel{~}{a}_i\stackrel{~}{a}_j|,$$
(10)
and requiring that det\[$`1𝐥_AA^{}A]=0`$, we find
$$(1\alpha _1^2)(1\alpha _2^2)=a_1|a_2^2.$$
(11)
Operators $`B`$ and $`C`$ are also accordingly constrained by
$`(1\beta _1^2)(1\beta _2^2)=b_1|b_2^2,`$ (12)
$`(1\gamma _1^2)(1\gamma _2^2)=c_1|c_2^2.`$ (13)
To completely characterize our protocol, we will further impose that it succeeds with the greatest possible probability over all OSBPs. The optimal OSBP probability,
$$P(\psi )\underset{\alpha _1,\beta _1,\gamma _1}{\mathrm{max}}2(\alpha _1\beta _1\gamma _1\mu _1)^2,$$
(14)
—where the constraints $`\alpha _1\beta _1\gamma _1\mu _1=\alpha _2\beta _2\gamma _2\mu _2`$ and (11-13) hold— reads
$`P`$ $`=_{x>0}^{max}{\displaystyle \frac{f_1(x)f_2(x)}{2}}\left(1\sqrt{1{\displaystyle \frac{4(1a_1|a_2^2)}{f_1^2(x)}}}\right)`$ (15)
$`\times `$ $`\left(1\sqrt{1{\displaystyle \frac{4\mu _1^2\mu _2^2(1b_1|b_2^2)(1c_1|c_2^2)}{f_2^2(x)}}}\right),`$ (16)
where
$`f_1(x)(x^2+1)/x,`$ (17)
$`f_2(x)(\mu _2^2x^2+2\mu _1\mu _2b_1|b_2c_1|c_2x+\mu _1^2)/x.`$ (18)
The maximization in (15) involves a polynomial equation of degree 6, which in general requires numerical calculations. For illustrative purposes and later reference, we will next consider two particular situations which can be solved analytically.
Suppose, first, that the parties share the state
$$|\varphi ^1=\mu _1|00c_1+\mu _2|11c_2,\mu _1\mu _2>0$$
(19)
that is, with the local states already orthogonal in Alice’s and Bob’s subsystems. Then the previous optimization over all OSBP leads to a maximal probability
$$P(\varphi ^1)=1\sqrt{14\mu _1^2\mu _2^2(1c_1|c_2^2)}.$$
(20)
This probability is two times the smallest eigenvalue of $`\rho _C`$, which has a decreasing behaviour under LOCC. Monotonicity of this eigenvalue implies therefore that $`P(\varphi ^1)`$ is also the maximal probability for any general distillation protocol. It turns out that Claire —i.e. the party that has to orthonormalize its local states— is the only one that needs to perform a local POVM.
The second particular case concerns states of the form
$$|\varphi ^2=\mu _1|0b_1c_1+\mu _2|1b_2c_2,\mu _1\mu _2>0$$
(21)
that is with orthogonal states in Alice’s subsystem. The maximum probability for OSBP is
$`P`$ $`(\varphi ^2)=(1+2\mu _1\mu _2b_1c_1|b_2c_2)`$ (22)
$`\times `$ $`\left(1\sqrt{1{\displaystyle \frac{4\mu _1^2\mu _2^2(1b_1|b_2^2)(1c_1|c_2^2)}{(1+2\mu _1\mu _2b_1c_1|b_2c_2)^2}}}\right).`$ (23)
Again, only Bob and Claire need to act on the system, with the corresponding POVMs satisfying
$$\frac{\beta _1^2}{\beta _2^2}=\frac{\mu _1}{\mu _2}\frac{\mu _1b_1|b_2+\mu _2c_1|c_2}{\mu _2b_1|b_2+\mu _1c_1|c_2},\frac{\gamma _1^2}{\gamma _2^2}=\frac{\beta _2^2}{\beta _1^2}\frac{\mu _1^2}{\mu _2^2}.$$
(24)
So far we have analyzed the optimal distillation protocols under the constraint that only one branch leads to a GHZ. In what follows we will show that no distillation protocol can succeed with probability greater than that for OSBP, $`P(\psi )`$. In order to do so, we will study the behavior of $`P(\psi )`$ under LOCC, to conclude that it is a decreasing entanglement monotone. That is, given the state $`|\psi `$ and a sequence of local quantum operations that transform it into $`|\psi _i`$ with probability $`p_i`$, we will show that
$$P(\psi )\underset{i}{}p_iP(\psi _i),$$
(25)
which means that the average probability to obtain a GHZ state from $`|\psi `$ using several branches is not greater than when using just one branch.
Although the set of transformations LOCC is very large and one should in principle check (25) for any local protocol, we can use the fact that any such protocol decomposes into individual POVMs. Indeed, we need only prove the monotonic character of $`P(\psi )`$ under the most general local POVM on each subsystem. But, as a matter of fact, due to the symmetry of the problem it suffices to consider a general local POVM performed by one of the parties, say Alice. Notice, furthermore, that any POVM can be decomposed into a sequence of two-outcome POVMs . Let us then consider a two-outcome POVM with operators $`\{N_1,N_2\}`$, $`N_1^{}N_1+N_2^{}N_2=1𝐥_A`$, applied by Alice. With some probability $`p_i`$ the resulting state will be $`|\psi _ip_i^{1/2}N_i1𝐥_B1𝐥_C|\psi `$, and then the parties can apply the optimal OSBP, with $`A_iB_iC_i`$ being the corresponding successful branch (see figure 2(i)). We want to show that
$$P(\psi )p_1P(\psi _1)+p_2P(\psi _2).$$
(26)
Notice that the global action on Alice’s side in order to distill a GHZ state can be reproduced by means of a single four-outcome POVM, namely by $`\{A_1N_1,\overline{A}_1N_1,A_2N_2,\overline{A}_2N_2\}`$, where $`\overline{A}_i`$ are the operators satisfying $`A_i^{}A_i+\overline{A}_i^{}\overline{A}_i=1𝐥_A`$. The second and fourth operators disentangle the state. They are irrelevant for distillation, and we can join them together into some other operator $`R`$ (figure 2(ii)). Both operators $`A_1N_1`$ and $`A_2N_2`$, being the last transformation Alice applies to its subsystem, have to leave the local states orthonormal and therefore can be written as
$$A_iN_i=\alpha _{1,i}|0\stackrel{~}{a_1}|+\alpha _{2,i}|1\stackrel{~}{a_2}|.$$
(27)
Consequently, this three-outcome POVM can be replaced with a two-outcome POVM with operators $`R`$ and
$$Q\sqrt{\alpha _{1,1}^2+\alpha _{1,2}^2}|0\stackrel{~}{a}_1|+\sqrt{\alpha _{2,1}^2+\alpha _{2,2}^2}|1\stackrel{~}{a}_2|,$$
(28)
followed —if the outcome corresponds to operator $`Q`$— by a diagonal, two-outcome POVM (figure 2(iii)),
$$D_i=\frac{\alpha _{1,i}}{\sqrt{\alpha _{1,1}^2+\alpha _{1,2}^2}}|00|+\frac{\alpha _{2,i}}{\sqrt{\alpha _{2,1}^2+\alpha _{2,2}^2}}|11|.$$
(29)
Let $`|\varphi ^2`$ in (21) be the normalized vector after applying $`Q`$ to $`|\psi `$. Then we can translate condition (26) into
$$P(\varphi ^2)q_1P(\varphi _1^2)+q_2P(\varphi _2^2),$$
(30)
where $`q_i`$ is the probability of the outcome related to $`D_i`$ starting from $`|\varphi ^2`$, and $`|\varphi _i^2`$ is the corresponding final state.
In order to finally check that (30) holds, we further notice that the action of any two-outcome POVM on a given state $`|\psi `$ can be reproduced by a conditional series of balanced, two-outcome POVMs, namely POVMs such that, for each of them, the probability of the two outcomes is exactly $`1/2`$ . That is, we only need to analyze a diagonal POVM where the square of the diagonal elements are $`x/(2\mu _1^2)`$ and $`(1x)/(2\mu _2^2)`$ for the first operator $`D_1`$, where $`x[2\mu _1^21,2\mu _1^2]`$, and their completion to 1 for the second. Then (30) is an inequality for $`x`$ that exhaustive numerical calculations have shown to saturate only for $`x=\mu _1^2`$, that is, when Alice acts trivially on her subsystem with $`D_i=1/\sqrt{2}1𝐥_A`$. Thus we conclude that the optimal probability for OSBP cannot be improved by any distillation protocol.
Finally, we have also considered the optimal local approximate transformation of $`|\psi `$ into a GHZ state. That is, given the state $`|\psi `$ and an arbitrary local protocol which transforms it into the (possibly mixed) state $`\rho _i`$ with probability $`p_i`$, we have looked for the maximal averaged fidelity between the output states $`\rho _i`$ and a GHZ state,
$$F_{opt}\underset{LOCC}{\mathrm{max}}\underset{i}{}p_iGHZ|\rho _i|GHZ.$$
(31)
As we argued in the bipartite setting , we only need to consider final pure states $`\{\psi _i,p_i\}`$ due to the manifest linearity of (31) in $`\rho _i`$ and the possibility of purifying any protocol that produces mixed states. Exhaustive numerical search shows then that the optimal approximate transformation consists in a deterministic transformation of $`|\psi `$ into $`|\psi ^{}`$ by means of local unitary transformations, exactly as in the bipartite case.
We have addressed the optimal distillation of a GHZ state starting from a pure state of a three-qubit system. We have shown that, contrary to what happens in the bipartite case, tripartite distillation requires that each party performs at least one local POVM. This feature can be associated to the lack of Schmidt decomposition for tripartite pure states. An alternative decomposition, based on non-orthogonal product states, has proven very useful to indicate which modifications must be introduced locally in a pure state in order to distill it into a GHZ state, and to optimize such distillation.
We acknowledge financial support by the Austrian SF, by the Spanish MEC (AP98 and AP99) and by the European Community (ESF; TMR network ERB-FMRX-CT96-0087; project EQUIP; HPMF-CT-1999-00200). This work was concluded during the 2000 session of the Benasque Center for Science, Spain.
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# The multiphase nature of the intra–cluster medium of some clusters of galaxies
## 1 Introduction
Studies of X-ray emission from the hot intracluster medium (ICM), using the latest data and analysis methods, led many to infer that the ICM may contain various gas components considerably cooler than the Virial temperature (see, e.g., Allen 2000 and Buote et al. 1999). The discovery of excess EUV and soft X-ray emission from clusters of galaxies (Lieu et al. 1996a,b; Mittaz, Lieu, & Lockman 1998) was originally interpreted as due to substantial amounts of warm gas at temperatures of $`10^6`$ K. In this paper we present evidence for the widespread existence of even cooler gases, which bolsters the possibility of a warm intermediate phase and the notion of a multiphase ICM with sub-Virial temperature gases playing an important role in a proper understanding of clusters.
Our results are obtained by a search for signatures of spatially resolved absorption in the C-band images of the ROSAT PSPC (defined here and after as the passband between PI channels 20 and 41, or approximately 0.2 – 0.4 keV). This band is ideally suited to the detection of intervening cold ICM gas: any depletion in the background cluster emission cannot be attributed to such other possibilities as, e.g., warm absorbers. The C-band also responds to small amounts of HI: for incident radiation from a plasma of 0.5 solar abundance and temperatures kT = 0.1 and 5.0 keV, reminiscent of the warm and hot ICM, an optical depth in the C-band corresponds to cold gas of HI column density N$`{}_{H}{}^{}1.3\times 10^{20}`$ and $`1.7\times 10^{20}`$ cm<sup>-2</sup>, respectively. Moreover, at high Galactic latitudes the absorption of C-band extragalactic emission by our interstellar medium is not very severe, so that data of good statistical quality are available for the brighter clusters.
## 2 Method of analysis
The method we adopted is to evaluate the spatial smoothness of the soft excess from two clusters, using PSPC C-band images. The first step consists of estimating the hot gas contribution to the C-band fluxes. To accomplish this, the cluster diffuse emission was divided into concentric annuli and spectra in the R47 band (a terminology often used to denote PI channels 42-201, or energies 0.4 – 2.0 keV) was modelled with a photoelectrically absorbed thin plasma emissivity code (MEKAL in XSPEC), where temperature (T) and elemental abundances (A) were fitted to the data (see Table 1). In the R47 band, in fact, any contributions from a ‘warm’ gas are minimal, and this modelling returns a spatially resolved measurement of the hot ICM parameters. At each radius $`r`$, the best-fit model defines a ratio $`f(r)`$ of the C-band to R47 count rates; removal of the hot ICM contribution from the C-band is obtained via the equation
$$c_{res}=cf(r)\times h,$$
(1)
where $`c`$ and $`h`$ are respectively the background subtracted C-band and R47 images.<sup>1</sup><sup>1</sup>1Since azimuthal symmetry is assumed, $`f(r)`$ is a radially symmetric ‘image’ of the C-to-R47 conversion factor described in the text. The distribution of the C-band residuals, $`c_{res}`$, reveals regions of soft excess (positive values) as well as regions of absorption (negative values). For the Coma cluster, a similar technique was adopted, whereby PSPC spectra were succesfully modelled with a photoelectrically absorbed MEKAL code with T and A fixed at the best-fit Ginga measurements (T=8.21 keV, A=0.21 solar, Hughes et al. 1993). Given that the central region of the cluster is sufficiently isothermal, the ratio f(r) is a constant function of radius.
After removal of the hot ICM contribution, a box of size 2.25$`{}_{}{}^{}\times `$ 2.25 scans the residual image area ($`c_{res}`$ in the notation of Eq. 1), stepping its center by 0.25 arcmin at a time. The box size was chosen to ensure sufficient enclosed counts for normal distribution of Poisson fluctuations. Then, if a background subtracted image is smooth, its mean brightness should be zero and deviations about the mean should follow a fixed gaussian. This, a consequence of the central limit theorem, was confirmed by our own simulations as well as C-band images of blank fields. As an example, we show in Figure 1 the results for various annuli of a blank field within the central PSPC area; the data were acquired during a 30 ksec pointing to a celestial direction where the Galactic N<sub>H</sub> is $`1.4\times 10^{21}`$ cm<sup>-2</sup> (RA=50.29<sup>o</sup>, DEC=40.74<sup>o</sup>, RP number 800034a01). The C-band sky background is low, due to its anti–correlation with HI, yet the forementioned box size still yielded an average of $``$ 40 counts per box, so that gaussian statistics apply.
When we apply the method to a cluster field, the X-ray emission is often not smooth to begin with. Deviations from the azimuthally averaged surface brightness for a R47 image of Coma are shown in Fig. 2a (see later for details on the observation): the image clearly reveals substructures. Nonetheless our method of analysis remains meaningful, because such regions of anisotropy do not correspond to a different spectral hardness when compared with the surroundings (i.e. the same conversion factor $`f`$ applies when subtracting the hot ICM contribution to the C-band flux). Further, the lack of any resemblance between the soft excess and X-ray images, as we demonstrate by Fig. 2b, confirms that any spatial structures in the former are due to emission or absorption of the soft cluster flux. For the Virgo cluster, similar results hold, with the exception of three strong absorption features detected towards the cluster center, which are positionally coincident with enhancements of the X-ray (0.4-2 keV) emission. These regions are the subject of a separate spectral analysis in the following section.
## 3 The Virgo and Coma clusters
The central region of the Virgo cluster is dominated by the giant galaxy M87; near the center of Coma are likewise located the cluster’s two brightest supergiant elliptical galaxies (NGC 4874 and NGC 4889) and a few discrete X–ray sources. We therefore excluded from analysis the innermost portions of the two clusters (angular radii $``$ 3 arcmin). We show in Figure 3 the C-band soft–excess image smoothness evaluation for the region between radii of 3 and 15 arcmin, centered at M87, and divided into 4 annuli. The dataset is from the ROSAT archive (observation number RP8000365). Significant departures from gaussian behavior, clearly noticeable within a radius of 9 arcmin, are in the form of an asymmetric extension to the left, which is symptomatic of absorption features resolved by the scanning box<sup>2</sup><sup>2</sup>2 Emission features, including any residual radial surface brightness gradients not removed by the subtraction of C-band radiation from the hot ICM, would have been manifested as right extension tails, which are evidently absent.. Indeed the features within 6 arcmin radius are obviously identifiable in the image.
A simple ‘partial covering’ model to interpret the observations assumes that for each annulus inwards of 12 arcmin, a certain fraction of its area (which can be broken into smaller, disjointed areas) has emission uniformly shadowed by foreground intrinsic HI, as illustrated in Fig. 4. Then the distribution of soft emission is the sum of two gaussians: their peak positions contain information about the mean brightness in the absence of intrinsic absorption and the HI column for the absorbed areas. Moreover, the relative normalization of the gaussians converts to a partial covering factor. In Figure 5 we show the performance of this model for one of the Virgo annuli, and in Table 2 the best-fit parameters of the model are listed for all the regions. The decreasing trend of the covering factor points to a radially declining influence of absorption by clouds within the cluster. We emphasize that despite its success as indicated by the $`\mathrm{\Delta }\chi ^2`$ values when compared with a single gaussian fit (Table 2), the model may still be over-simplistic because the absolute $`\chi ^2`$ values do not imply acceptability. Indeed, additional spatially unresolved absorption may well be present, and the effect of foreground cluster emission may well have caused an underestimate of the HI towards these central radii. It is also possible that the column density of HI clouds is not uniform; nonetheless the quality of present data does not warrant more sophisticated models, for this purpose one must await observations by the XMM/EPIC instrument.
In Figure 6 we show the three deepest absorption features which exist in the innermost area; they are positionally coincident with prominent radio lobes (Harris et al. 1999) and with enhancements in the X–ray (0.5–2 keV) emission, see Böhringer et al. (1995), whose study led to the proposition of a lower temperature for the hot ICM in the radio lobes than in the surrounding regions. The forementioned enhancement was then interpreted as due to a Fe line feature at $``$ 1 keV. We accordingly examined spectra for two regions, one including one feature to the east and the other to the south-west of M87. Since similar results apply to both regions, here we focus on the eastern knot: Figure 7a shows that a photo-absorbed MEKAL code with T=1.3 and A=0.45 (as in Böhringer et al. 1995) can in fact model the observed spectrum ($`\chi ^2`$=188 for 155 degrees of freedom, with a null probability of 3.6 %). It is however not physically obvious why radio features (symptomatic of relativistic particles) are positionally coincident with a temperature decrease in the hot gas. An alternative explanation of the C-band flux reduction would invoke cold gas (HI): we therefore attempted to model the two regions using the best-fit hot ICM parameters for their corresponding annulus (T=1.47 keV, A=0.46 solar, see Table 1): the eastern knot shows a -7$`\sigma `$ decrement from the model in the 0.2-0.4 keV band (fig. 7b), yet if the total absorption was allowed to vary we could obtain a good fit ($`\chi ^2`$=165 for 154 degrees of freedom, null probability of 25 %). The extra absorption required above the Galactic line-of-sight column is $`N_H`$ 4.5$`\times 10^{19}`$ cm<sup>-2</sup> for a region $``$ 2 square arcmin in size. A possible interpretation is that the high pressure of relativistic particles can lead to compression of the hot gas trapped inside the radio lobes to high density, with consequent rapid cooling to very low temperatures. The evidence for depleted emission in other regions of the cluster (Fig. 3), as well as in the Coma cluster (see below), strengthens our present scenario, because none of those are caused by the subtraction of spatially corresponding enhancements in the hot ICM radiation.
A similar analysis of PSPC C–band data of the Coma cluster (archival identification RP800005) proved equally fruitful. The spatial distribution of the emission is clearly not smooth, see Fig. 8. The use of a two–gaussian ‘partial covering’ model results in significant improvement of the goodness–of–fit (see Table 3), where the decrease of the covering factor with radius points again to a centrally peaked distribution of absorbers; the outermost annulus (12–15 arcmin) in fact does not require a second gaussian. As an example, the performance of the two–gaussian model for the 6–9 arcmin region is shown in Fig. 9. This result is of particular importance as Coma is the first non–cooling flow cluster which exhibits evidence for cold gas clouds.
## 4 Discussion and conclusions
The adopted method of analysis affords us estimates of column densities and HI mass for the two clusters, see Table 4. Could this absorption be caused by line–of–sight cluster galaxies? Rich clusters of galaxies, such as Virgo and Coma, are known to have a large fraction of gas–poor elliptical and S0 galaxies, and only a small fraction of late type galaxies (such as spirals). Moreover, cluster galaxies show an HI deficit when compared to field galaxies of same morphology, especially toward the central regions (e.g. Huchtmeier and Richter 1986; Dickey 1997). In the case of Virgo a simple galaxy count for the 3–12 arcmin region from M87 (using the NED database) revealed only three known, early type member galaxies along the line of sight, which is too few in number to account for the absorption features of Fig. 6. Likewise, a recent 21-cm line investigation of the Coma cluster with the VLA array (Bravo–Alfaro et al. 2000) detected no galaxies with HI content above $`10^8`$ M, the typical detection limit of that survey, in the central 10 arcminutes. Within the context of a galaxy origin of the cold gas, this is again inconsistent with our HI estimates in Table 4.
It is far more plausible that the reported effect is due to intracluster HI which might have, at least in part, been released to the intergalactic space by ram pressure stripping. Cold gas masses of the order 10$`{}_{}{}^{911}M_{}^{}`$ contribute to only a small fraction of the cluster’s total baryonic mass, the latter of order a few $`\times 10^{12}M_{}`$ for Virgo (see, e.g., Bahcall and Sarazin 1979) and a few $`\times 10^{14}M_{}`$ for Coma (Briel, Henry and Böhringer 1992).
We note that the phenomenon of the soft excess radial trend (meaning rising importance of the soft component with radius), known to exist in the clusters A1795 (Mittaz, Lieu and Lockman 1998) and A2199 (Lieu et al. 1999), was recently interpreted as indicative of centrally peaked intrinsic cluster absorption (Lieu, Bonamente and Mittaz 2000). As shown in Figure 10, this effect is also present in the C–band excess of Virgo, and may indeed be due to the gradual disappearance of absorption with increasing radius reported in Table 4.
The co–existence of a cold phase and hot ICM components raises questions concerning if and where an intermediate warm phase is present. The evidence for cold HI presented in this paper reinforces the thermal interpretation of the cluster soft–excess syndrome: some warm gas can certainly be generated at the interface between the cold and hot gases, e.g. by the ‘mixing layer’ mechanism (Fabian 1997). The emerging scenario is a three–phase ICM, where the hot gas emits the bulk of the X–ray emission, and the warm gas is responsible for the soft–excess emission. The cooler gas (T $`10^5`$ K), found by our spatial smoothness test, causes substantial reduction of EUV and soft X–ray fluxes in some locations, and the soft excess radial trend detected in the C band of Virgo (see Fig. 10) could then be due to the absorption effect of a centrally peaked distribution of HI.
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Figure 1: Spatial distribution of events for the inner annuli of a blank field centered at PSPC boresight. A 2.25$`{}_{}{}^{}\times 2.25^{}`$ scanning box detects, for each detector position, the deviation z between the number of enclosed counts (N) and the mean background level ($`\mu `$), in units of the standard deviation $`\sigma `$ (where z=(N-$`\mu `$)/$`\sigma `$). The y–axis gives total number of boxes at a given deviation z and the dashed line is the best–fit profile expected from a smooth distribution, viz., a gaussian of fixed width $`\sigma `$=1. The gaussian mean was varied to account for any possible systematics in the background subtraction.
Fig. 2: (a) At each detector position, a 2.25$`{}_{}{}^{}\times 2.25^{}`$ scanning box detects the difference between the number of enclosed R47 counts and the azimuthally averaged value at that radius, in units of the standard deviation (see caption of Fig. 1). (b) C-band residual image of Coma: same scanning box detects, at each detector position, the deviation between the number of enclosed counts and 0, the mean expected value in the absence of absorption or soft excess emission. Unit of measure is the standard deviation $`\sigma `$. The $`\sigma `$-scale of 2(b) is set to reveal the absorption features responsible for the left tail of Figure 3; these features are clearly not positionally coincident with any emission features of the hot ICM, as is evident from a comparison of 2(a) with 2(b). Regions containing obvious point sources were excluded from the analysis (gray boxes).
Figure 3: Spatial distribution of events for the central region of the Virgo cluster. As in Fig. 1, dashed line gives the best–fit model for a smooth distribution (i.e. one–gaussian) except that it fails to account for the data here. Annular radii are measured from the position of M87, which is at PSPC boresight.
Figure 4: The scanning box (light grey) can detect silhouettes (dark grey) of the surface brightness which are responsible for the bimodal behavior of the spatial distribution of events (see e.g. Fig. 4 and Fig. 7).
Figure 5: Fitting the distribution of events in the 6–9 arcmin annulus of Virgo with the two–gaussian ‘partial covering’ model. A significant improvement over one gaussian is achieved (see also Table 1). The positive mean of the right (main) gaussian revelas soft–excess emission.
Figure 6: C–band soft excess image of Virgo, in units of statistical significance $`\sigma `$ of signals within a 2.25$`{}_{}{}^{}\times 2.25^{}`$ scanning box centered at each pixel position. The deep absorption features to the south of M87 (cross) are coincident with the location of radio lobes (Harris et al. 1999).
Figure 7: (a) PSPC spectrum of a $``$ 2 square arcmin region encompassing the ‘absorption feature’ to the east of M87 (cross), fitted to a MEKAL code with the parameters of Böhringer et al. (1995). (b) same PSPC spectrum when only PI channels 42-201 ($`0.42`$ keV) are fitted to a MEKAL code with the best-fit parameters of the 0-3 arcmin annulus (see Table 1).
Figure 8: Event distribution for the Coma cluster (see caption of Fig. 3); the dashed line is the best–fit single gaussian model, as before. The center of our annular system is at the X–ray centroid of Coma, which is near boresight. The positive mean of the right gaussian reveals soft excess emission.
Figure 9: Use of the ‘partial covering’ model for the 6–9 annulus of Coma brings considerable improvement to the fit (see also Table 2).
Figure 10: The ‘soft-excess radial trend’ effect of the Virgo cluster, illustrated by a plot against cluster radius of the soft X-ray fractional excess $`\eta `$, defined as $`\eta =(pq)/q`$, where for a given annulus $`p`$ is the observed C–band flux after subtracting the sky background, and $`q`$ is the expected flux from the hot ICM as determined by fitting the PI channels 50 – 200 ($``$ 0.5 – 2.0 keV) using the MEKAL thin plasma emission code (Mewe, Lemen and van den Oord 1986; Mewe, Gronenshild and van den Oord 1985; Kaastra 1992) and Galactic absorption according to Morrison & McCammon (1983).
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# Tunneling of the Closed Friedmann Universe with Generation of Scalar Waves
## 1 Statement of the problem
Here we present the basic equations of the problem.
In the space-time of the closed cosmological Friedmann models with space-time metric
$$ds^2=dt^2a^2(t)[d\chi ^2+\mathrm{sin}^2\chi (d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)],$$
(1)
where $`0\chi ,\theta \pi ,0\phi 2\pi `$ , the action is assumed to be
$$S=\left[\frac{R}{16\pi G}\stackrel{}{\mathrm{\Lambda }}+\frac{1}{2}\phi ^{,i}\phi _{,i}\frac{1}{2}m^2\phi ^2\right]\sqrt{g}𝑑\mathrm{\Omega },$$
(2)
Here $`R`$ is a scalar curvature, $`\phi `$ is a potential of the scalar quantum field with a mass $`m`$, ($`\mathrm{}=c=1`$), $`\stackrel{}{\mathrm{\Lambda }}`$ is a cosmological $`\mathrm{\Lambda }`$ \- term.
The equation for $`\phi `$ has the form $`\mathrm{}\phi +m^2\phi =0`$. Below let us set the scalar field as radial standing waves with time-dependent amplitudes
$$\phi _k=\mathrm{A}_k(t)f_k(\chi );f_k(\chi )=\frac{\mathrm{sin}k\chi }{\mathrm{sin}\chi }$$
(3)
Note that the functions $`f_k(\chi )`$ do not have singularities on the poles $`(\chi =0,\pi )`$ for integer numbers $`k`$, and that they are orthonormal on the coordinate $`\chi `$ in the measure defined by metric (1), i.e.,
$$\mathrm{sin}^2\chi f_kf_nd\chi =\delta _{kn}.$$
(4)
Now the action (2) can be integrated over the angular variables, and we rewrite it after replacement of $`R`$ by its value through $`a`$ as
$$S=\frac{3\pi }{4G}\left[\underset{k}{}\left(a^3\underset{k}{\overset{2}{\stackrel{.}{\mathrm{A}}}}k^2a\mathrm{A}_k^2m^2a^3\mathrm{A}_k^2\right)\mathrm{\Lambda }a^3\stackrel{2}{\stackrel{.}{a}}a+a\right]𝑑t,$$
(5)
where
$$\mathrm{A}_k=\sqrt{8G/3}\mathrm{A}_k(t);\mathrm{\Lambda }=8\pi G\stackrel{}{\mathrm{\Lambda }}/3.$$
These redefinitions are made to get $`G`$ and $`\pi `$ only in front of the integral, and the integrand to be free of both constants.
In this formula, the point denotes the derivative with respect to $`t`$, and one supposes that all $`k1`$, i.e., we assume that waves with wavelength comparable to the diameter of the Universe are negligible for the dynamics.
The waves packet $`(_k)`$ ensures the homogeneity of the model during expansion with a high accuracy. After introduction of generalized momenta the Hamiltonian function can be presented as
$$H=\frac{P_a^2}{4a}a+\mathrm{\Lambda }a^3+\underset{k}{}\left(\frac{q_k^2}{4a^3}+k^2\mathrm{A}_k^2a+m^2a^3\mathrm{A}_k^2\right)=0.$$
(6)
Here $`P_a,a`$ are a generalized momentum and a coordinate for the Universe as entire; $`q_k`$ and $`\mathrm{A}_k`$ are the same for each radial mode of the field. Corresponding to eq. (6) the Wheeler-DeWitt equation for the wave function of the Universe (WF) has the form
$$\frac{^2\mathrm{\Psi }}{a^2}\frac{1}{a^2}\underset{k}{}\frac{^2\mathrm{\Psi }}{\mathrm{A}_k^2}V\mathrm{\Psi }=0,$$
(7)
where the superpotential $`V`$ is given as
$$V=4\left[a^2\frac{p}{8a^2}\left(1\frac{p}{2}\right)\underset{k}{}k^2\mathrm{A}_k^2a^2\mathrm{\Lambda }a^4\right].$$
(8)
Here we are neglecting $`_km^2\mathrm{A}_k^2a^4`$ in comparison with the sum residual in eq. (8) since $`k1`$. This is justified because in the short-wavelength approximation we are applying here, the mass of the particles does not have much influence.
The factor ordering parameter $`p`$ is also introduced in eq. (8).
## 2 Semiclassical approximation
We are going to search a solution of equation (7) in semiclassical approximation
$$\mathrm{\Psi }_c=\mathrm{exp}(iS_c),$$
(9)
where $`S_c`$ is a classical action. In this approximation we have
$$\left(\frac{S_c}{a}\right)^2\frac{1}{a^2}\underset{k}{}\left(\frac{S_c}{\mathrm{A}_k}\right)^2+V=0,$$
(10)
This nonlinear equation of the first order is similar to the Hamilton-Jacobi equation of analytical mechanics. For its solution one can take advantage of the complete integral . It can be found from the system of characteristic equations . As the eq. (10) supposes the solution for $`(S/a)`$, the scale factor becomes an argument along the characteristics. From here we have
$$\frac{dS_c}{da}=\frac{V}{F};\frac{d\mathrm{A}_k}{da}=\frac{q_k}{a^2F};\frac{dq_k}{da}=\frac{4k^2\mathrm{A}_ka^2}{F}.$$
(11)
Here $`V`$ is indicated by (8), and
$$F=\sqrt{(\underset{k}{}q_k^2)/a^2V}.$$
Now the evolution of the WF $`\mathrm{\Psi }`$ goes along a characteristic. Turning points are the roots of $`V(a)=0`$. They separate classically available and forbidden regions along the characteristic eq. (11) . A similar approach for the strongly anisotropic WF in a Bianchi type I model is presented in ref. , and for a closed Friedmann model in ref. .
In the region under the barrier the semiclassical solution can be found by an ansatz of the form
$$\mathrm{\Psi }=\mathrm{exp}(|S_e|).$$
(12)
The equation for $`S_e`$ can be obtained from eq. (10) by changing $`VV`$. Thus, as distinct from eq. (11) the system of the characteristics has the form
$$\frac{dS_e}{da}=\frac{V}{F_e};\frac{d\mathrm{A}_k}{da}=\frac{q_k}{a^2F_e};\frac{dq_k}{da}=\frac{4k^2\mathrm{A}_ka^2}{F_e};$$
(13)
$$F_e=\sqrt{(\underset{k}{}q_k^2)/a^2+V}.$$
## 3 The evolution of the scalar field modes in the classically allowed region
Let us investigate how amplitudes of short wavelength change at the above mentioned evolution of the WF. Equations for $`q_k`$ and $`\mathrm{A}_k`$ follow from (11)
$$\frac{d\mathrm{A}_k}{da}=\frac{q_k}{a^2F};\frac{dq_k}{da}=\frac{4k^2\mathrm{A}_ka^2}{F},$$
(14)
where
$$F=\sqrt{(\underset{k}{}q_k^2)/a^2V}.$$
This system is an essentially non-linear one since co-factors of $`q_k`$ and $`\mathrm{A}_k`$ in the right hand sides of the equations depend on $`a`$, too. However, the presence of the large parameter $`k^2`$ allows finding a solution.
Following from the system (14), the equation for the generalized momenta $`q_k`$ is
$$\frac{d^2q_k}{da^2}=\frac{4k^2q_k}{F^2}+\left(\frac{2}{a}\frac{F^{}}{F}\right)\frac{dq_k}{da},$$
(15)
where $`{}_{}{}^{}=d/da`$. Here $`F(a)`$ is an unknown function according to eq. (14).
It is easy to show that the following solution
$$q_k=\frac{C_{k0}a}{\sqrt{2k}}\mathrm{cos}\mathrm{\Phi }(a);\mathrm{\Phi }(a)=2k\underset{0}{\overset{a}{}}\frac{da}{F}$$
(16)
satisfies the last equation with the high accuracy of $`\mathrm{O}(1/k^2)`$. And correspondingly,
$$\mathrm{A}_k=\frac{C_{k0}}{2\sqrt{2}ak^{3/2}}\mathrm{sin}\mathrm{\Phi }(a),$$
(17)
where $`C_{k0}`$ is an arbitrary constant, and the phase of the solution is chosen from the requirement of amplitude finiteness at $`a0`$. The mentioned accuracy of solutions requires the assumption that the logarithmic derivative $`F^{}/F`$ does not give a large parameter $`k`$. The last statement follows from the fact that each mode of wave enters into eq. (14) for $`F`$ as a term
$$\frac{q_k^2}{a^2}+4k^2\mathrm{A}_k^2a^2=\frac{C_{k0}^2}{2k}=const.$$
(18)
It is clear from substitution eqs. (16), (17) in (18). With taking into account (18) after rewriting of the amplitude $`C_{k0}/\sqrt{2k}=C_k`$, we have for $`F`$ the following expression
$$F=\sqrt{\underset{k}{}C_k^24\left[a^2\mathrm{\Lambda }a^4\frac{p}{8a^2}\left(1\frac{p}{2}\right)\right]},$$
(19)
Further, we suppose that $`\mathrm{\Lambda }1`$ and we have in the superpotential
$$V4\left[a^2\frac{p}{8a^2}\left(1\frac{p}{2}\right)\frac{1}{4}\underset{k}{}C_k^2\mathrm{sin}^2\mathrm{\Phi }(a)\right].$$
(20)
Values of $`V`$ and $`F`$ define an action of WF $`\mathrm{\Psi }(a)`$ according to the system of characteristics (11). Because of the large number of short waves, the sum in eq. (20) can be taken as fast oscillations average with an adequate accuracy. Then the expression for the superpotential gets the form
$$V=4\left[a^2\frac{p}{8a^2}\left(1\frac{p}{2}\right)\frac{1}{8}\underset{k}{}C_k^2\right],$$
(21)
The existence of an inner classically available region requires a realization of the condition $`V(a)<0`$ in the mentioned interval $`0aa_0`$. (Let us remember — see eq. (6) — that the total energy $`E`$ vanishes for the closed universe.) The potential barrier should begin at $`a=a_0`$ , i. e. $`V(a)>0`$ at $`aa_0`$. Let us notice that the vanishing of the factor ordering parameter ($`p=0`$) allows us to get an inner classically available region only at the expense of the sum $`C_k^2`$. But in that case, a potential barrier necessary for a further development is absent. So, we choose the factor ordering $`p=1`$ in the following.
Owing to the amplitude of short waves begin to increase quickly under the barrier, we assume that their total energy is negligible in the considered region $`V(a)<0`$, and the evolution of the Universe is determined only by the first two terms in the r.h.s. of eq. (21). In this case the dynamics of the amplitudes of the short waves can be given explicitly. The phase in eq. (16) is
$$\mathrm{\Phi }(a)=\frac{k}{2}\mathrm{arcsin}\left(\sqrt{1\left(\frac{a}{a_0}\right)^4}\right).$$
(22)
Here, $`a_0`$ can be obtained from the condition
$$V(a_0)=4\left[a_0^2\frac{1}{16a_0^2}\right]=0,a_0=\frac{1}{2}.$$
(23)
A case when the total energy of the modes is not negligible in the expressions for $`V`$ and $`F`$ can be similarly considered. But then the final expression becomes essentially complicated.
## 4 Modes of the scalar field under the barrier
The superpotential $`V`$ is positive at $`a>a_0`$. This means that the evolution of the WF occurs under the barrier. At this range, the problem is defined by the system of characteristics (13), and we have for $`q_k`$ the following equation
$$\frac{d^2q_k}{da^2}=\frac{4k^2q_k}{F_e^2}+\left(\frac{2}{a}\frac{F_e^{}}{F_e}\right)\frac{dq_k}{da}.$$
(24)
This equation has an inverse sign at first term on the right in comparison with eq. (15) for the classically allowed region. It changes the solution in principle, but for a large parameter $`k^2`$ we can find a solution by a method similar to the one used for eq. (15)
$$q_k=C_ka\mathrm{cosh}\mathrm{\Phi }_e(a);$$
$$\mathrm{A}_k=\frac{C_k}{2ak}\mathrm{sinh}\mathrm{\Phi }_e(a);$$
(25)
$$\mathrm{\Phi }_e(a)=2k\underset{a_0}{\overset{a}{}}\frac{da}{F_e}.$$
Here the conditions of continuity of $`q_k`$ and $`\mathrm{A}_k`$ at the point $`a=a_0`$ were used for the selection of the constants of integration. Both in section 3 and here, the correctness of the solution (25) is determined by the requirement of $`F_e^{}/F_ek`$. This requirement is checked by substituting of the obtained solution in $`F_e`$ of eq. (13). There
$$F_e=\sqrt{\underset{k}{}C_k^2+4a^2}.$$
In the last formula the term containing the factor ordering, that plays a part only close to zero, is ignored.
Because of the smallness of the sum $`_kC_k^2`$ under the barrier this term can be neglected in the expression for $`F_e`$, and we suppose with a necessary accuracy
$$F_e(a)2a.$$
(26)
At that $`a`$ the phase $`\mathrm{\Phi }_e`$ has the form
$$\mathrm{\Phi }_e=k\mathrm{ln}\left(\frac{a}{a_0}\right),$$
and leaving only growing modes in solutions at $`a>a_0`$ we obtain
$$q_k=\frac{C_ka}{2}\left(\frac{a}{a_0}\right)^k,\mathrm{A}_k=\frac{C_k}{4ka}\left(\frac{a}{a_0}\right)^k.$$
(27)
In that solution the amplitude of short waves catastrophically increases with growth of $`k`$ but not under the exponential law. The latter is realized in the case of a slow variation of $`F_e`$ in eq. (25). It is similar to the Rubakov-effect of catastrophic particle creation, see ref. .
A final value of an amplification of the amplitude of short waves is determined by a “time” of the WF stay under the barrier. Let us notice that we consider a self-consistent problem of amplification of short waves (an analogue of particle creation). It means that the energy of the increasing modes under the barrier controls the form of the barrier itself (see Fig. 1).
Let us consider the equation for the action of the WF under the barrier (13)
$$\frac{dS_e}{da}=\frac{V}{F_e};F_e2a;V=4a^2\frac{C_k^2}{4}\left(\frac{a}{a_0}\right)^{2k}.$$
(28)
Only a “potential” energy of increasing modes is present in the superpotential $`V`$, but here again it is not cancelled by their “kinetic” energy as it was realized in (26).
As it follows from $`V(a_0,a_1)=0`$, an external boundary of the tunneling region is determined by the equation
$$V=4a_1^2\frac{C_k^2}{4}\left(\frac{a_1}{a_0}\right)^{2k}.$$
The final value of amplification of the amplitude of the modes is defined by substitution of $`a_1`$ in eq. (27).
Further, the Euclidean action $`S_e(a_1)`$ has a form
$$S_e(a_1)=\frac{1}{16}\underset{k}{}C_k^2\left(\frac{a_1}{a_0}\right)^{2k}\left(\frac{1}{k}1\right).$$
(29)
Subject to $`k1`$ we have $`S_e=a_1^2`$. Hence, a probability of tunneling of the WF through the barrier is given as
$$w=\mathrm{exp}(2|S_e|)=\mathrm{exp}(2a_1^2)$$
(30)
at $`w1`$. The formula
$$w=\frac{\mathrm{exp}(2a_1^2)}{1+\mathrm{exp}(2a_1^2)}$$
is more exact in the case $`w<1`$.
## 5 Conclusions
The self-consistent problem of the evolution of the closed Friedmann Universe with an amplification of wave perturbations of a scalar field under a barrier is considered. Within the framework of a semiclassical approximation of the Wheeler-DeWitt equation it is shown that the increase of the amplitude of short waves under the barrier controls the barrier shape, i.e. leads to both a reduction of the barrier and fast output of WF from under the barrier. In this sense the problem is self-consistent. The process of tunneling and the increase of the amplitudes of short waves is defined here by the assignment of their initial spectrum in the inner classically allowed region. If it is required after tunneling in the Universe to receive a spectrum close to the flat one, the initial spectrum should be almost exponentially suppressed in the area $`k\mathrm{}`$. It is qualitatively equivalent to the presence of an effective temperature in the inner classically allowed region.
A more detailed discussion of these problems obviously requires a consideration of the quantum nature of the problem for short-wave perturbations of a scalar field.
ACKNOWLEDGEMENT
The authors thank A. A. Starobinsky and M. V. Sazhin for information about papers of V. A. Rubakov et al. and H. Kleinert for useful comments. H.-J.S. acknowledges financial support from the HSP III-program. The work is done within the project KR-154 of the International Science and Technology Centre (ISTC).
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# Tumbleweeds and airborne gravitational noise sources for LIGO
## I Introduction
Interferometric detectors such as LIGO and VIRGO rely on exquisite sensitivity to the positions of hanging test masses in order to detect the perturbations of passing gravitational radiation. The sensitivity is so great that the measurements can also be affected by fluctuations in the local Newtonian gravitational field, which create tiny accelerations of the mass. This noise source, known as gravity gradient noise or Newtonian gravitational noise, is caused by the near-field gravity of masses moving near the interferometer, and is not to be confused with the far-field propagating gravitational waves that the instruments are intended to measure.
Gravity gradient noise has the potential to be quite insidious, since it cannot be shielded by improvements to the test-mass vibrational isolation. The only effective way to eliminate gravity gradient noise is to eliminate the moving masses that create the perturbing fields. Fortunately, though, the strongest perturbations to the local gravitational field are at frequencies well below the detectors’ pass-bands. Since the proposed terrestrial interferometric detectors all have sensitivity cutoffs around 3 Hz or higher, we need only worry about noise sources that can perturb the local gravity field on timescales less than about 0.3 seconds. Most of the noise sources that I consider are motivated by the expected sensitivity of advanced LIGO interferometers, which were originally projected to have a low-frequency cutoff around 10 Hz, and instrumental noise of $`S_h2\times 10^{45}\mathrm{Hz}^1`$ in the gravitational-wave signal output at that frequency Abramovici\_A:1992 <sup>1</sup><sup>1</sup>1Specifically, I will be using the noise curve in Fig. 10 of Abramovici\_A:1992 , not Fig. 7, whose suspension thermal noise is a factor of 3 too small.. Although changes in instrumentation technology will modify the ultimate sensitivity goals of LIGO, this “standard” advanced LIGO noise level is a good reference point when considering new noise sources. Also, as pointed out below, gravity gradient noise will make it difficult to push the detector noise much below this level at 10 Hz, regardless of improvements to the interferometers.
Saulson Saulson\_P:1984 was the first to estimate the effect of gravity gradient noise on terrestrial interferometric detectors, considering the effects of seismic waves passing through the earth and of sound waves in the air. In both cases he found the spectral density of noise in the interferometer path-length difference to be less than $`10^{39}\mathrm{m}^2\mathrm{Hz}^1`$ around 10 Hz, corresponding to noise in the gravitational-wave signal at levels less than $`10^{46}\mathrm{Hz}^1`$ for a 4 km interferometer. By comparison, this is significantly less than the noise floor of $`2\times 10^{45}\mathrm{Hz}^1`$ that advanced LIGO interferometers expect to achieve at 10 Hz. More recently, a detailed analysis Hughes\_S:1998 has been made of seismic gravity gradient noise; this study indicated that seismic gravity gradient noise would be within a factor of 2 of the advanced LIGO noise floor at 10 Hz for most times, and could actually *exceed* this noise floor during seismically noisy times, making seismic gravity gradients a significant barrier to improvements in low-frequency sensitivity. It therefore seems prudent to revisit the issue of *atmospheric* gravity gradients as well.
In Sec. II I consider gravity gradients caused by atmospheric pressure perturbations—the same noise source considered by Saulson. Attention is paid, however, to the effects of the ground, and of buildings that reduce the pressure noise in the immediate vicinity of the interferometer test masses. Nonetheless, I find that these tend only to weaken the gravity gradient noise in the pass-bands of interferometric detectors, reinforcing the conclusion that this noise source is not of any great concern.
A much stronger source of high-frequency density perturbations in the atmosphere is the presence of temperature fluctuations, which are advected past a detector by the wind. In Sec. III I analyze this as a potential source of gravity gradient noise. However, I find that while small-scale temperature perturbations can produce high-frequency temperature fluctuations at any given point, they do not produce the same high-frequency fluctuations in the test mass position, since a given pocket of warm or cool air will affect the test mass gravitationally over the entire time that it is in the vicinity of the test mass, which is on the order of seconds. This produces a cutoff in the noise spectrum above a few tenths of a Hz. The presence of turbulent vortices around the interferometer buildings can increase the high-frequency component, but still probably not enough to show up in the gravitational-wave noise spectrum.
In Secs. IV and V I turn away from sources of background noise to consider possible sources of transient gravity gradient signals that might be detected as spurious events in the gravitational-wave instruments. Sec. IV extends the analysis in Sec. II to look at the effects of atmospheric shockwaves, such as might be generated by an explosion or supersonic aircraft. I find that sources such as these can indeed produce detectable signals that might be interpreted spuriously as gravitational-wave events. However, such signals would easily be vetoed using acoustic monitors outside the interferometer buildings.
Sec. V analyzes the gravity gradients produced by individual objects, such as wind-borne debris, moving around outside the interferometer buildings. Saulson considered this effect for the case of objects moving with fairly uniform velocity, but typically, in order to produce significant signal above 3 Hz, an object’s motion must be changing on timescales of less than 0.3 seconds. In particular, I find that objects *colliding* with the interferometer buildings produce much stronger signals than objects simply passing by the buildings. As an example, tumbleweeds at the Hanford LIGO facility will be a steady source of spurious signals in advanced detectors if they are allowed to collide with the end stations. Preventing such signals requires shielding a region of a few tens of meters around the end station, screening any wind-borne debris that masses more than a few hundred grammes.
Sec. VI presents some concluding remarks, including recommendations to the gravitational-wave experimental groups and possible directions for further research.
## II Atmospheric pressure waves
Pressure perturbations are the only type of atmospheric gravity gradient noise considered by Saulson Saulson\_P:1984 . The derivation in this section gives largely the same result as his.
Consider a plane pressure wave with frequency $`f`$ propagating through a homogeneous airspace at some sound speed $`c`$. If the fractional pressure change $`\delta p/p`$ is small, it will induce an adiabatic density change $`\delta \rho /\rho =\delta p/\gamma p`$, where $`\gamma 1.4`$ is the ratio of heat capacities at constant pressure and constant temperature for air at normal temperatures. The gravitational acceleration produced by this wave in the direction of propagation $`𝐞_z`$ is:
$$g_z(t)=\frac{Gz\delta \rho }{r^3}𝑑V=\frac{2G\rho c}{\gamma pf}\delta p(t+1/4f),$$
(1)
where $`\rho 1.3\mathrm{kg}\mathrm{m}^3`$ and $`p10^5\mathrm{N}\mathrm{m}^2`$ are the ambient air density and pressure, and $`\delta p(t)`$ is the pressure perturbation measured at the same point as the acceleration is being measured. By symmetry, there is no acceleration transverse to the wave.
Now consider sound waves in the vicinity of the interferometer. First, since the interferometer is only sensitive to motions of the test mass parallel to the arms, the gravitational acceleration is reduced by a factor $`\mathrm{cos}\theta `$, where $`\theta `$ is the angle between the propagation direction and the interferometer arm.
Second, the interferometer test mass is inside a building, which can in principle be used to suppress noise within a distance $`r_{\mathrm{min}}`$ of the test mass. Roughly speaking, this results in a high-frequency cutoff factor $`C(2\pi fr_{\mathrm{min}}/c)`$, where the function $`C(x)`$ depends on the precise shape of the building, the manner in which it reflects sound waves, and many other factors, but is normally close to 1 for $`x1`$. For instance, if one simply removes from the volume integral in Eq. (1) a cylinder with length and diameter both $`2r_{\mathrm{min}}`$ aligned with the $`z`$-axis, then $`C(x)1`$ for $`x1`$, but oscillates with an amplitude of $`0.3`$ for $`x1`$. This function is shown in Fig. 1.
The constant-amplitude oscillations of $`C(x)`$ for large $`x`$ reflect the assumption that the sound wave has a coherence length much longer than the building size, so the field between the two ends of the excluded cylinder is generated almost entirely by the first half-wavelength beyond each cap. Realistically, the actual high-frequency behaviour of $`C(x)`$ will depend on how the sound waves bend and scatter around the building; however, this should not change the order of magnitude of $`C(x)`$. The behaviour shown in Fig. 1 is therefore probably a good estimate of the true cutoff function for $`x1`$, and a reasonable order of magnitude estimate for $`x1`$. For the LIGO end stations, $`r_{\mathrm{min}}`$ is of order 5 metres, giving $`xf/(10\mathrm{H}\mathrm{z})`$; the factor will not be too far off for the frequencies of greatest interest. More precise estimates would depend on the specific architectural details of a particular facility.
Third, the interferometer is on the ground, not in homogeneous empty space. For simplicity I assume that the waves are almost entirely reflected off the ground; I will justify this assumption in Sec. II.1. In this case the gravity gradient in directions parallel to the ground contributed by the reflected wavefront is the same as if the wavefront were extended below ground, while the pressure perturbations measured by detectors near the ground (much less than a wavelength) will be doubled. The acceleration experienced by an interferometer test mass is therefore:
$$g_z(t)=\frac{G\rho c}{\gamma pf}\mathrm{cos}(\theta )C(2\pi fr_{\mathrm{min}}/c)\delta p(t+1/4f).$$
(2)
The gravitational wave signal $`h(t)`$ in the interferometer is related to the acceleration of one of the test masses by $`\ddot{h}(t)=g(t)/L`$, or in frequency space $`\stackrel{~}{h}(f)=(2\pi f)^2\stackrel{~}{g}(f)/L`$, where $`L`$ is the length of the interferometer arm. Thus:
$$\stackrel{~}{h}(f)=\frac{G\rho c}{4\pi ^2\gamma pLf^3}\mathrm{cos}(\theta )C(2\pi fr_{\mathrm{min}}/c)i\stackrel{~}{\delta p}(f).$$
(3)
Assuming stationary noise, the one-sided spectral density $`S_h(|f|)`$ is given by $`\stackrel{~}{h}(f)\stackrel{~}{h}(f^{})^{}=S_h(|f|)\delta (ff^{})`$, where $`\mathrm{}`$ denotes an expectation over all random phases of all plane wave modes contributing to the noise, and denotes complex conjugation. Taking mode amplitudes and directions to be uncorrelated, this gives:
$$S_h(|f|)=\left[\frac{G\rho c}{4\pi ^2\gamma Lf^3}C(2\pi fr_{\mathrm{min}}/c)\right]^2\mathrm{cos}^2\theta \frac{S_p(|f|)}{p^2},$$
(4)
where $`S_p(|f|)`$ is the acoustic noise spectral density measured outside the building in the vicinity of a particular test mass. Since the two test masses in an arm are many wavelengths apart, their noise will be uncorrelated, and will thus add in noise power. Similarly, the noise from the two arms will add in power. (Actually this is a bit of an overestimate, since the noise in the motion of the test masses at the corner station will be somewhat correlated.) Noting that $`\mathrm{cos}^2\theta =1/3`$, one finds that the total noise in the gravitational wave signal is:
$$S_h(|f|)=\left(\frac{G\rho c}{4\pi ^2\gamma L}\right)^2\frac{1}{3f^6p^2}\underset{i=1}{\overset{4}{}}C(2\pi fr_{\mathrm{min}}^{(i)}/c)^2S_p^{(i)}(|f|),$$
(5)
where $`i`$ denotes a particular test mass in the interferometer, $`r_{\mathrm{min}}^{(i)}`$ is the dead air radius about the $`i`$th test mass, and $`S_p^{(i)}(|f|)`$ is the acoustic noise spectrum measured outside the building enclosing that test mass.
Infrasound noise spectra should be taken at the actual interferometer sites, but one can make estimates based on typical terrestrial atmospheric noise. An empirical study Posmentier\_E:1974 collected 256 power spectra of 1–16 Hz infrasound data from a rural forest 50km from New York City over a period of months, and found that the average noise spectrum $`S_p(f)`$ was relatively flat at 6–16 Hz, though with widely varying amplitude: 25% of the spectra had noise under $`100\mathrm{n}\mathrm{b}\mathrm{a}\mathrm{r}^2/\mathrm{Hz}`$, 50% under $`300\mathrm{n}\mathrm{b}\mathrm{a}\mathrm{r}^2/\mathrm{Hz}`$, and 75% under $`1000\mathrm{n}\mathrm{b}\mathrm{a}\mathrm{r}^2/\mathrm{Hz}`$. I use Eq. (5) to compute the corresponding noise in the LIGO detector. Assuming that the end masses (with $`r_{\mathrm{min}}5\mathrm{m}`$) dominate the contribution to the gravity gradient noise and contribute equally, the noise in the gravitational-wave signal around 10 Hz is:
$`S_h(|f|)`$ $``$ $`(6\times 10^{48}\mathrm{Hz}^1)\left({\displaystyle \frac{f}{10\mathrm{H}\mathrm{z}}}\right)^6\left({\displaystyle \frac{S_p(|f|)}{1000\mathrm{n}\mathrm{b}\mathrm{a}\mathrm{r}^2/\mathrm{Hz}}}\right)`$ (6)
$`\times C(f/10\mathrm{H}\mathrm{z}).`$
The results are plotted in Fig. 2, using infrasound power spectra read off of Fig. 3 of Posmentier\_E:1974 .
Even the third-quartile power spectrum is between two and three orders of magnitude below the expected noise floor of $`2\times 10^{45}\mathrm{Hz}^1`$ at 10 Hz projected for advanced LIGO interferometers. Thus ambient infrasound is probably a negligible effect for determining the noise floor for most interferometric gravitational-wave detectors. Nonetheless the issue cannot be completely resolved without infrasonic noise data from the actual interferometer sites.
### II.1 Ground absorption
In the derivation above I treated the sound waves as being reflected off the ground. Now consider what happens if a sound wave is absorbed by the ground. The energy flux in a traveling compression wave is $`C^{3/2}\rho ^{1/2}(\delta x/x)^2`$, where $`C`$ is the compression modulus of the medium ($`C_{\mathrm{air}}=\gamma p_{\mathrm{air}}`$), and $`(\delta x/x)^2`$ is the average squared dimensionless compression factor over a wave cycle. The gravity gradient noise induced by such a wave goes as $`c^2(\delta \rho )^2`$, where $`(\delta \rho )^2\rho (\delta x/x)^2`$, and $`c=\sqrt{C/\rho }`$ is the wave speed in that medium. Therefore, if a sound wave is completely absorbed by the ground, the resulting ground motions will produce gravity gradient noise contributions in the ratio:
$$\frac{S_h^{(\mathrm{ground})}}{S_h^{(\mathrm{air})}}=\left(\frac{\rho _{\mathrm{ground}}}{\rho _{\mathrm{air}}}\right)^{3/2}\left(\frac{C_{\mathrm{ground}}}{C_{\mathrm{air}}}\right)^{1/2}.$$
(7)
Now for $`C_{\mathrm{air}}=\gamma p_{\mathrm{air}}=1.4\times 10^5\mathrm{N}\mathrm{m}^2`$, $`C_{\mathrm{ground}}3\times 10^8\mathrm{N}\mathrm{m}^2`$, $`\rho _{\mathrm{air}}=1.3\mathrm{kg}\mathrm{m}^3`$, $`\rho _{\mathrm{ground}}1.8\times 10^3\mathrm{kg}\mathrm{m}^3`$, the ratio turns out to be of order $`10^3`$; that is, if sound waves were completely absorbed by the ground, the resulting ground vibrations would produce gravity gradient noise levels about 1000 times greater than the atmospheric gravity gradients. However, it was shown in Hughes\_S:1998 that seismic gravity gradients are only just large enough to worry about. So the only way that atmospheric gravity gradients can be larger or of the same order as seismic gravity gradients is if the waves are mostly reflected off of the ground. This was one of the assumptions used in deriving Eq. (2).
## III Atmospheric temperature perturbations
The largest small-scale atmospheric density perturbations are caused not by pressure waves but by temperature perturbations. As heat is transported up through a convective atmospheric layer, convective turbulence mixes pockets of warm and cool air to form temperature perturbations on all lengthscales down to a few millimetres. On the timescales of interest (less than a second) these perturbations are effectively “frozen” into the airmass, while pressure differences disperse rapidly in the form of sound waves. Perturbations in the air density $`\rho p/T`$ are therefore caused predominantly by the temperature perturbations, which are typically several orders of magnitude larger than the pressure perturbations. Although they are frozen into the airmass, these temperature perturbations can cause rapid time-varying density fluctuations $`\delta \rho =\rho \delta T/T`$ as the wind carries them past a point in space. This is the primary source of “seeing” noise that affects optical astronomy.
The appendix to this chapter gives a rigorous mathematical derivation of the gravity gradient noise spectrum due to these temperature perturbations. This section gives a qualitative derivation that reproduces the final result to order of magnitude.
The gravity gradient signal at some frequency $`f`$ is caused by pockets of warm or cool air with some lengthscale $`l`$ being advected past the interferometer test mass at a speed $`v`$, where $`lv/2\pi f`$. Consider a single such pocket of air with a temperature perturbation $`\delta T`$ away from the ambient temperature $`T`$. The gravitational acceleration produced in the instrument as a function of time $`t`$ is $`g_x(t)=G\rho l^3(\delta T/T)x(t)r^3(t)`$, where $`\rho `$ is the ambient air density, $`r(t)`$ is the distance of the air pocket from the test mass as it is blown past, and $`x(t)`$ is this distance projected onto the axis of the interferometer arm. This geometry is sketched in Fig. 3.
Now in general, the noise power spectral density in any quantity $`a`$ due to a background of independent, uncorrelated events is $`S_a(|f|)=(2/\mathrm{\Delta }t)|\stackrel{~}{a}(f)|^2`$, where $`\stackrel{~}{a}(f)`$ is the Fourier spectrum from a single event and $`\mathrm{\Delta }t`$ is the spacing between events. Assuming uncorrelated pockets of air, then independent pockets of air arrive along any given streamline at intervals $`\mathrm{\Delta }tl/v`$, and streamlines separated by more than $`l`$ add noise incoherently (i.e., add linearly in power). This gives the following noise spectrum:
$$S_g(|f|)\frac{2l}{v}_{\{S\}}\frac{dA}{l^2}\left(\frac{G\rho }{T}\right)^2\delta T^2(l)|\stackrel{~}{G}_S(f)|^2,$$
(8)
where $`_{\{S\}}`$ denotes an integral over a plane crossing all streamlines $`S`$, and $`\stackrel{~}{G}_S(f)`$ is the Fourier transform of the function $`G(t)=x(t)/r^3(t)`$ taken along a given streamline. The quantity $`\delta T^2(l)`$ is the average squared temperature difference between points a distance $`l`$ apart. Turbulent mixing theory, as well as actual micrometeorological measurements, predict a power-law behaviour for small separations: $`\delta T^2(l)c_T^2l^p`$, where $`p`$ is typically $`2/3`$. This applies for horizontal separations $`l`$ up to of order 50 times the height of a given air pocket above the ground Busch\_N:1972 . For streamlines more than a metre or so above ground level, then, this behaviour for $`\delta T^2`$ should be good out to distances $`l50`$m, corresponding to frequencies $`0.2\mathrm{Hz}(v/10\mathrm{m}\mathrm{s}^1)`$.
Eq. (8) gives the gravity gradient noise on a given test mass in the interferometer. Denoting the test masses by the index $`k=1\mathrm{}4`$, and assuming that each test mass contributes independently to the noise in the gravitational-wave signal $`h`$, one has $`S_h(|f|)=(2\pi f)^4L^2_kS_g(|f|)`$. Combining this with Eq. (8) and the relation $`lv/2\pi f`$ yields:
$`S_h(|f|)`$ $``$ $`2\left({\displaystyle \frac{G\rho }{LT}}\right)^2c_T^2(2\pi f)^{(p+7)}v^{p+4}`$ (9)
$`\times {\displaystyle \underset{k}{}}{\displaystyle _{\{S\}}}dA|\stackrel{~}{G}_{S,k}(f)|^2.`$
The more rigorous analysis in the appendix (Sec. A) gives an expression with roughly the same form, but covers the factors of order unity, and also accounts for the fact that wind speed can vary along a streamline and between streamlines. The more accurate formula is:
$`S_h(|f|)`$ $`=`$ $`2\pi ^2\left({\displaystyle \frac{G\rho }{LT}}\right)^2c_T^2(2\pi f)^{(p+7)}\mathrm{sin}(p\pi /2)\mathrm{\Gamma }(p+2)`$ (10)
$`\times {\displaystyle \underset{k}{}}{\displaystyle _{\{S\}}}\stackrel{~}{F}_{S,k}(f)^{}\stackrel{~}{G}_{S,k}(f)wdA,`$
where $`w`$ is the wind speed of the streamline as it crosses the plane of integration, and $`\stackrel{~}{F}_{S,k}(f)`$ and $`\stackrel{~}{G}_{S,k}(f)`$ are Fourier transforms of functions $`F_{S,k}(t)`$ and $`G_{S,k}(t)`$ describing the motion of a point along a streamline $`S`$ past a test mass $`k`$, of the form:
$`F(t)`$ $`=`$ $`{\displaystyle \frac{x(t)}{r(t)^3}}v(t)^{p+3},`$ (11)
$`G(t)`$ $`=`$ $`{\displaystyle \frac{x(t)}{r(t)^3}}.`$ (12)
It is worth noting that the frequency structure of $`S_h(|f|)`$ depends on the time behaviour of the functions $`x(t)`$ and $`r(t)`$ describing the position of a point on a streamline relative to a test mass, which can be some distance away. The minimum distance $`r_{\mathrm{min}}`$ from the test mass to the passing air is thus an additional important scale in the problem: if $`x(t)`$ and $`r(t)`$ change significantly only on timescales $`r_{\mathrm{min}}/v`$, then the noise spectrum will be cut off at frequencies $`v/2\pi r_{\mathrm{min}}`$.
By comparison, the *temperature* noise spectrum $`S_T(|f|)`$ measured at a point depends only on the local properties of the atmosphere at that point. Applying similar order-of-magnitude arguments, one can write $`S_T(|f|)(2l/v)c_T^2l^p|\stackrel{~}{H}(f)|^2`$, where $`H(t)1`$ for $`|t|l/v`$ and 0 otherwise. At high frequencies $`0.2\mathrm{Hz}(v/10\mathrm{m}\mathrm{s}^1)`$ the system involves only one lengthscale $`l`$, giving the spectrum a power-law dependence:
$$S_T(|f|)2c_T^2v^p(2\pi f)^{(p+1)}.$$
(13)
This is the same, to order of magnitude, as the exact result given in Eq. (39).
### III.1 Uniform airflow
The gravity gradient noise is easy to compute from Eq. (10) for the case of uniform airflow parallel to the ground with some constant velocity $`𝐯`$. Placing the reference plane orthogonal to $`𝐯`$ and passing through the test mass, the equations of motion for a given streamline past the test mass become quite simple: $`v(t)=w=v`$, $`x(t)=vt\mathrm{cos}\theta +r_0\mathrm{sin}\theta \mathrm{cos}\psi `$, $`r(t)=(r_0^2+v^2t^2)^{1/2}`$, where $`r_0`$ is the distance from the test mass to the nearest point on the streamline, and $`r_0\mathrm{cos}\psi `$ is the projection of this distance onto the ground. The geometry of the airflow is shown in Fig. 4. It is easy to show that:
$`\stackrel{~}{G}(f)`$ $`=`$ $`{\displaystyle \frac{4\pi f}{v^2}}[i\mathrm{cos}\theta K_0(2\pi fr_0/v)`$ (14)
$`+\mathrm{sin}\theta \mathrm{cos}\psi K_1(2\pi fr_0/v)],`$
$`\stackrel{~}{F}(f)`$ $`=`$ $`v^{p+3}\stackrel{~}{G}(f),`$ (15)
where $`K_0`$ and $`K_1`$ are moddified Bessel functions of the second kind of orders 0 and 1. I perform the integral over the above-ground half of the reference plane, out from some radius $`r_{\mathrm{min}}`$ that is roughly the closest distance that the outside air can approach the test mass. This gives a noise contribution from a single test mass equal to:
$`S_h(|f|)`$ $`=`$ $`8\pi ^2\mathrm{cos}(\pi [p1]/2)\mathrm{\Gamma }(p+2)\left({\displaystyle \frac{G\rho r_{\mathrm{min}}}{TL}}\right)^2`$ (16)
$`\times `$ $`c_T^2\left({\displaystyle \frac{v}{2\pi f}}\right)^p(2\pi f)^5\{\mathrm{cos}^2\theta [K_0^2(x)K_1^2(x)]`$
$`+{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta [K_1^2(x)K_0(x)K_2(x)]\},`$
where $`x=2\pi fr_{\mathrm{min}}/v`$.
For typical values $`v10`$m/s and $`r_{\mathrm{min}}5`$m, at frequencies above 10 Hz or so, one has $`x30`$ or more, well into the exponentially damped regime of the Bessel functions. Even for gale-force winds of 30m/s or so, the argument $`x`$ of the Bessel functions will still be of order 10 or more. The asymptotic expansions of $`K_0`$, $`K_1`$, and $`K_2`$ give $`K_0^2(x)K_1^2(x)K_1^2(x)K_0(x)K_1(x)\pi e^{2x}/(2x^2)`$. Also, I note that the total noise will be dominated by the contribution from the two end stations, which have smaller $`r_{\mathrm{min}}`$ than the corner station. So the total noise in the interferometer for uniform airflow is:
$`S_h(|f|)`$ $`=`$ $`(2\pi )^3\mathrm{cos}(\pi [p1]/2)\mathrm{\Gamma }(p+2)\left({\displaystyle \frac{G\rho }{TL}}\right)^2c_T^2`$ (17)
$`\times \left({\displaystyle \frac{v}{2\pi f}}\right)^p(2\pi f)^5e^{4\pi fr_{\mathrm{min}}/v}.`$
I consider “typical” values of $`p=2/3`$, $`\rho =1.3\mathrm{kg}\mathrm{m}^3`$, $`T=300\mathrm{K}`$, $`L=4000\mathrm{m}`$, and $`r_{\mathrm{min}}=5\mathrm{m}`$. The parameters $`v`$ and $`c_T^2`$ can vary on a minute-by-minute basis, and should really be measured at the site of a given interferometer. However, $`c_T^20.2\mathrm{K}^2\mathrm{m}^{2/3}`$ is a typical daytime peak temperature fluctuation index Coulman\_C:1985 , and $`v20\mathrm{m}/\mathrm{s}`$ might be typical of a fairly windy day. At frequencies around 10 Hz, this gives a noise spectrum of:
$$S_h(|f|)(1.6\times 10^{40}\mathrm{Hz}^1)\left(\frac{f}{10\mathrm{H}\mathrm{z}}\right)^{23/3}10^{14(f/10\mathrm{H}\mathrm{z})}.$$
(18)
The two dotted lines in Fig. 7 show the gravity gradient noise spectra computed from Eq. (18) for wind speeds of 10m/s and 30m/s.
Clearly the exponential cutoff makes this a negligible source of noise for LIGO or similar detectors. Physically this cutoff arises from the fact that the gravity from a particular temperature perturbation passing near the end station will affect the test mass coherently over the second or so that it takes to travel the width of the end station. Thus, even though the temperature noise spectrum has a high-frequency power law tail (reflecting the fact that temperature perturbations exist on all lengthscales), the gravity gradient signal will have this much sharper exponentially cut-off tail.
### III.2 Potential flow near the end station
As described above, uniform airflow is not likely to produce much atmospheric gravity gradient noise in the pass-band of interferometric detectors, since the shortest timescale over which the gradients change is of order the wind crossing time of the interferometer buildings. However, if an air pocket could be made to accelerate over shorter timescales, it might produce a stronger gravity gradient signal at high frequencies.
One possibility is the acceleration of the air as it is forced up and around the wall of an end station: streamlines that approach the ground-wall corner of the end station can have curvature scales much shorter than the building size. Treating the flow as incompressible and vorticity-free, the resulting velocity field near the corner is (p. 27 of Landau\_L:1989 ):
$$v_x=2AX,v_z=2AZ,$$
(19)
where $`A=v/2R`$ is a constant, and $`R`$, $`X`$, and $`Z`$ are measured from the corner. This approximation is clearly only good near the corner, since it gives velocity increasing monotonically with radius; one would expect $`v`$ to approach the free-streaming airspeed $`V`$ at a distances $`R_m`$ of order the building height. Fig. 5 shows a schematic of the airflow near the corner.
It is clear from Eq. (19) that, although the streamlines are sharply curved near the corner of the end station, the advection speed is smaller in direct proportion. The shortest timescale over which the motion can change is of order $`R_m/V`$ for all streamlines. Thus the streamlines close to the corner will contribute no more high-frequency noise than the streamlines further out, at distances of order $`R_m`$, and the spectrum should not differ greatly from the one for uniform flow, Eqs. (17), (18).
### III.3 Vortices
Perhaps the most serious contenders for high-frequency atmospheric gravity gradient noise are circulating vortices of air near the end stations. This is somewhat stretching the assumptions of the formalism I have established, since I had previously separated the effects of the homogeneous turbulence (which establishes the temperature perturbations on various lengthscales), and wind flow (which carries these perturbations past the instrument). However, the results in this section are only expected to be good to order of magnitude anyway, in the absence of detailed hydrodynamic analysis of airflow past a particular interferometer building. I therefore apply the above formalism to a simple model for a turbulent-like flow past an interferometer end station.
For simplicity, I specialize to airflow along the axis of an interferometer arm, since these give the largest gravitational accelerations. A simple model for the turbulent flow is to take a uniform flow and then add a cycloidal motion to it. This gives $`x(t)vtR\mathrm{sin}(vt/R)`$, $`z(t)r_0R\mathrm{cos}(vt/R)`$, $`r(t)=\sqrt{z^2(t)+x^2(t)}`$, where $`R`$ is the radius of the cycloidal motion, and $`r_0r_{\mathrm{min}}`$ is the distance from the test mass to the unperturbed streamline. I treat the speed $`v`$ along the streamline as a constant. If $`R`$ is also a constant, then one would expect the Fourier transform of $`G(t)=x(t)/r^3(t)`$ \[Eq. (12)\] to have a spike at frequency $`f=v/2\pi R`$, with a width of order $`v/r_0`$. However, to give a somewhat more realistic behaviour for $`R`$, I treat it as growing from zero at the leading edge of the end station to some scale value $`R_0`$ over the half-length $`r_{\mathrm{min}}`$ of the end station: $`R(t)=R_0\sqrt{vt/r_{\mathrm{min}}}`$. The square-root dependence mimics the less-than-linear growth of the thickness of a boundary layer. Although quite crude, this model covers the essential features of a turbulent flow past the building: a uniform translation, accompanied by circulating motions over a range of radii with a scale set by $`R_0`$, with both the uniform flow rate and the circulation speed set by the free-streaming airspeed $`v`$. A typical streamline of this type is shown in Fig. 6.
The Fourier transform of $`G(t)`$ is too complicated to perform analytically, but is simple enough to compute numerically using a fast Fourier transform. The resulting $`\stackrel{~}{G}(f)`$ has the usual exponential cutoff with frequency scale $`v/2\pi r_{\mathrm{min}}`$, as for a smooth streamline, but then rises to a second peak value of $`(4/v)\sqrt{R_0^3/r_{\mathrm{min}}^5}`$ at a frequency $`f0.06v/R_0`$ and decreases from there as $`f^3`$. This high-frequency tail is the power law that one would expect from the cusps on the bottom of the cycloid. Since I treat $`v`$ as constant, $`\stackrel{~}{F}(f)=v^{p+3}\stackrel{~}{G}(f)`$. The cross-sectional area of streamlines that contribute significantly around this peak frequency is of order $`2r_{\mathrm{min}}R_0`$. Plugging these into Eq. (10) one gets a noise spectral density of
$$S_{h(\mathrm{max})}1.5\times 10^6\left(\frac{G\rho }{LT}\right)^2c_T^2R_{\mathrm{max}}^{35/3}r_{\mathrm{min}}^4v^5$$
(20)
at the peak frequency of $`f0.06v/R_0`$, where I have assumed $`p=2/3`$.
For wind speeds around 10m/s, a cycloid radius $`R_00.06\mathrm{m}(v/10\mathrm{m}/\mathrm{s})`$ puts the noise peak at the 10 Hz seismic noise wall for advanced LIGO detectors. The atmospheric gravitational noise contribution from a single end station is then:
$`S_{h(\mathrm{max})}`$ $``$ $`(1.3\times 10^{49}\mathrm{Hz}^1)\left({\displaystyle \frac{c_T^2}{0.2\mathrm{Km}^{2/3}}}\right)`$ (21)
$`\times \left({\displaystyle \frac{v}{10\mathrm{m}/\mathrm{s}}}\right)^{20/3}\left({\displaystyle \frac{r_{\mathrm{min}}}{5\mathrm{m}}}\right)^4.`$
This is over five orders of magnitude below the expected advanced LIGO noise floor of $`2\times 10^{45}\mathrm{Hz}^1`$ at 10 Hz. Gale-force winds ($`v30`$m/s) will bring this up to $`2\times 10^{46}\mathrm{Hz}^1`$, still an order of magnitude below the advanced LIGO noise curve. The dashed lines in Fig. 7 show the actual data from the numerical Fourier transforms for these two cases.
Since even the worst-case estimate is still an order of magnitude below the advanced LIGO noise floor (as well as the seismic gravity gradient noise floor in Hughes\_S:1998 ), it seems unlikely that turbulent vortices will be sufficient to raise the atmospheric gravitational noise to significant levels, even given the approximations made in this analysis. However, to settle this matter definitively would require a much more sophisticated numerical analysis of the temperature perturbations and airflow past the buildings of a particular facility.
## IV Shockwaves
Although atmospheric pressure waves are unlikely to be a significant source of gravity gradient noise in interferometric gravitational-wave detectors, the sudden pressure changes caused by atmospheric shockwaves could potentially produce detectable transient signals in the detector, if such shocks occur in the vicinity of the detectors. Shocks are specifically a matter of concern because they can produce significant pressure changes over timescales less than 0.1 s, corresponding to the lower end of the pass-bands of most interferometric detectors.
Consider a shock that produces a sudden jump in air pressure in the vicinity of one of the interferometer test masses: $`\delta p(t)=\mathrm{\Delta }p\mathrm{\Theta }(tt_0)`$. It is a simple matter to take the Fourier transform and apply Eq. (3) to obtain:
$$\stackrel{~}{h}(f)=\frac{G\rho c}{8\pi ^3\gamma Lf^4}\frac{\mathrm{\Delta }p}{p}ie^{2\pi ft_0}\mathrm{cos}(\theta )C(2\pi fr_{\mathrm{min}}/c).$$
(22)
Here, $`\theta `$ is the angle between the interferometer arm and the normal to the shock front. If the shock has a finite rise time $`\mathrm{\Delta }t`$, we can mimic this analytically by convolving in the time domain with a Gaussian of width $`\sigma \mathrm{\Delta }t`$. In the frequency domain this multiplies our amplitude by a Gaussian of width $`\sigma 1/\mathrm{\Delta }t`$, giving an exponential cutoff at frequencies above $`1/\mathrm{\Delta }t`$. Typical shocks from, for instance, supersonic objects have rise times on the order of a few milliseconds Hayes\_W:1971 , corresponding to cutoff frequencies of a few hundred Hz. However, one expects the dominant contribution of the signal to come from much lower frequencies, before the building-size cutoff factor $`C`$ kicks in.
Shocks are transient phenomena that will produce *signals* in the detector, rather than raising the noise floor. What one would like to know is what signal-to-noise ratio the shock will produce. In general this depends on what filters one is using to search for signals, and how well these filters overlap with the signal produced by a shock. However, the signal from a shock is likely to overlap quite well with templates designed to search for generic impulsive phenomena; such templates are likely to be used in advanced interferometers as control over nonstationary instrumental noise improves. Thus it is reasonable to consider the signal-to-noise ratio $`\rho `$ that a matched filter would give Thorne\_K:1999 :
$$\rho ^2=_0^{\mathrm{}}\frac{4\stackrel{~}{h}(f)\stackrel{~}{h}^{}(f)}{S_h(f)}𝑑f=_{\mathrm{}}^{\mathrm{}}\frac{|2f\stackrel{~}{h}(f)|^2}{fS_h(f)}d\mathrm{ln}f.$$
(23)
This shows, roughly speaking, that the relative magnitude of the dimensionless signal amplitude $`2|f\stackrel{~}{h}(f)|`$ and the dimensionless noise amplitude $`\sqrt{fS_h(f)}`$ over a logarithmic frequency interval gives a good indication of the signal-to-noise ratio produced in the detector.
### IV.1 Sonic booms
Sonic booms caused by supersonic bodies are one example of atmospheric shocks that might affect interferometric gravitational-wave detectors. Direct shockwaves from a supersonic aircraft will typically hit the ground in a “carpet” about 15–20 km wide under the aircraft’s flight path. Outside this carpet, the temperature gradient near the ground will completely reflect the shockwave before it touches down. However, shockwaves will also reflect downward off of the temperature inversion in the stratosphere and thermosphere, forming secondary and higher-order “carpets” out to many hundreds of kilometres Balachandran\_N:1977 . The presence or absence of these higher-order waves can depend quite sensitively on conditions in the upper atmosphere.
A detailed prediction of these effects is beyond the scope of this paper. However, to give an indication of their potential seriousness, I will consider what would happen if a supersonic aircraft were actually to overfly the instrument at a height of several kilometres.
The characteristic profile of a sonic boom is a symmetric N-wave, consisting of a shock that increases the pressure by an amount $`\mathrm{\Delta }p`$, followed by a smooth decrease in pressure of $`2\mathrm{\Delta }p`$ over a time $`\mathrm{\Delta }t`$, followed by a second rising shock $`\mathrm{\Delta }p`$ to restore the ambient pressure. According to Eq. (9.78) of Witham\_G:1974 , the strength of the shocks is:
$$\frac{\mathrm{\Delta }p}{p}\frac{2^{1/4}\gamma }{(\gamma +1)^{1/2}}(M^21)^{1/8}\kappa \delta l^{3/4}r^{3/4},$$
(24)
where $`\gamma `$ is the adiabatic coefficient of air (1.4 at normal temperatures), $`M`$ is the Mach number of the aircraft (its speed divided by the sound speed), $`\kappa `$ is a dimensionless form factor that depends on the shape of the aircraft (typically around 1), $`l`$ is the length of the aircraft, $`\delta `$ is the ratio of the aircraft’s typical thickness to its length, and $`r`$ is the closest distance that the aircraft came to the point of measurement. Between the two shocks, the rate of change of pressure in a direction $`𝐞_x`$ parallel to the line of flight is given, in Eq. (9.80) of Witham\_G:1974 , as:
$$\frac{1}{p}\frac{dp}{dx}\frac{\gamma }{\gamma +1}\frac{(M^21)^{1/2}}{M^2}\frac{1}{r}.$$
(25)
The shock fronts move outward at the sound speed $`c`$ in the direction orthogonal to their surface, while the entire cone travels in the $`𝐞_x`$ direction along with the aircraft at a speed $`Mc`$. The total change of pressure between the two shocks is $`2\mathrm{\Delta }p`$. From these facts and Eqs. (24) and (25), one can show that the time between the two shocks is:
$$\mathrm{\Delta }t2^{5/4}(\gamma +1)^{1/2}\frac{M}{(M^21)^{3/8}}\kappa \delta \frac{l^{3/4}r^{1/4}}{c}.$$
(26)
On frequency scales higher than $`1/\mathrm{\Delta }t`$ (typically a few Hz for a supersonic aircraft a few kilometres away), the sonic boom looks like simple Heaviside shocks, giving $`\delta \stackrel{~}{p}(f)1/f`$. At lower frequencies, though, the entire N-wave looks like the derivative of a $`\delta `$-distribution, giving $`\delta \stackrel{~}{p}(f)f`$. Performing the Fourier transform analytically and plugging into Eq. (3), one obtains:
$`\stackrel{~}{h}(f)`$ $`=`$ $`{\displaystyle \frac{G\rho c}{4\pi ^3\gamma L}}{\displaystyle \frac{1}{f^4}}{\displaystyle \frac{\mathrm{\Delta }p}{p}}\mathrm{cos}(\theta )C(2\pi fr_{\mathrm{min}}/c)`$ (27)
$`\times \left[{\displaystyle \frac{\mathrm{sin}(\pi f\mathrm{\Delta }t)}{\pi f\mathrm{\Delta }t}}\mathrm{cos}(\pi f\mathrm{\Delta }t)\right]e^{2\pi ift_0},`$
where $`t_0`$ is the time when the midpoint of the N-wave crosses the detector. As expected, the amplitude goes roughly as $`f^4`$, except for frequencies less than $`1/\mathrm{\Delta }t`$, where it goes as $`f^1`$.
Now let us plug in some typical numbers. The numbers $`G=6.67\times 10^{11}\mathrm{m}^3\mathrm{kg}^1\mathrm{s}^2`$, $`\rho =1.3\mathrm{kg}\mathrm{m}^3`$, $`\gamma =1.4`$, $`c=332\mathrm{m}\mathrm{s}^1`$, and $`L=4000\mathrm{m}`$ can be treated as constant. A supersonic jet aircraft might have a length of $`l=10`$m, a typical diameter of $`\delta l=2`$m, and be traveling at Mach $`M=1.5`$ at a distance of $`r=10`$km or so. Let $`\mathrm{cos}\theta `$ be 1 for an upper limit. Then $`\mathrm{\Delta }t0.2`$s, and for frequencies $`f10`$Hz the dimensionless signal amplitude is:
$`2|f\stackrel{~}{h}(f)|`$ $``$ $`1.4\times 10^{19}(M^21)^{1/8}C(2\pi fr_{\mathrm{min}}/c)\left({\displaystyle \frac{\delta }{0.2}}\right)`$ (28)
$`\times \left({\displaystyle \frac{l}{10\mathrm{m}}}\right)^{3/4}\left({\displaystyle \frac{r}{10\mathrm{k}\mathrm{m}}}\right)^{3/4}\left({\displaystyle \frac{f}{10\mathrm{H}\mathrm{z}}}\right)^3.`$
This is three orders of magnitude above the expected noise floor of $`\sqrt{fS_h(f)}1.4\times 10^{22}`$ at 10 Hz for advanced LIGO interferometers!
By contrast, consider a .30-calibre rifle bullet ($`l0.025`$m, $`\delta 0.3`$) passing at Mach 3 within 10m of an interferometer test mass. (This stretches the assumption of a plane-wave shock front at the test mass, but the order of magnitude should be correct.) The bullet produces a much stronger double shock, but with a time interval $`\mathrm{\Delta }t0.5`$ms, so $`1/\mathrm{\Delta }t2`$kHz. The low-frequency tail of this signal will have dimensionless amplitude:
$`2|f\stackrel{~}{h}(f)|`$ $``$ $`1.8\times 10^{23}{\displaystyle \frac{M^2(M^21)^{5/8}}{2.5}}C(2\pi fr_{\mathrm{min}}/c)`$ (29)
$`\times `$ $`\left({\displaystyle \frac{\delta }{0.3}}\right)^3\left({\displaystyle \frac{l}{0.025\mathrm{m}}}\right)^{9/4}`$
$`\times \left({\displaystyle \frac{r}{10\mathrm{m}}}\right)^{1/4}\left({\displaystyle \frac{f}{10\mathrm{H}\mathrm{z}}}\right)^1.`$
This is nearly an order of magnitude *below* the dimensionless noise amplitude in advanced LIGO, and therefore too small to be of any serious concern.
Fig. 8 shows more complete gravity gradient signal spectra computed using Eq. (27) with the above parameters for a supersonic aircraft and rifle bullet. These are plotted along with the anticipated dimensionless noise amplitude for advanced LIGO detectors.
### IV.2 Vetoing shockwave signals
While atmospheric shockwaves are a potential source of spurious signals in gravitational-wave detectors, they are easy to veto using environmental sensors. One need simply place infrasound microphones outside the buildings and test-mass vacuum enclosures. If these sensors detect a pressure change of more than a millibar over timescales of 50–100 milliseconds, then one might expect spurious signals with dimensionless amplitude of $`10^{22}`$ in the 10–20 Hz frequency range. The stretch of data containing the potential spurion can then be discarded.
Alternatively, if the same shock profile is detected in an array of at least three sound sensors, then one can determine the direction of propagation of the shock and predict the actual induced test-mass motions. The spurious signal could then be subtracted out of the data stream. This is a much trickier procedure, and would only be necessary in the unlikely event that significant amounts of data were being corrupted.
In either case, it is clear that infrasound sensors will be important environmental monitors for advanced interferometric detectors.
## V High-speed objects
Another potential source of spurious signals in the interferometer is the gravity gradient caused by the motion of an individual massive object past the interferometer end station, or the collision of such an object with the end station. The latter is particularly serious, since the sudden deceleration of the object can produce a signal at high frequencies. The issue of human-generated gravity gradient noise has been addressed in Thorne\_K:1999 , but there are other sources outside the facility that must be considered, such as stray bullets and wind-borne debris. In particular, the Hanford LIGO facility is plagued by tumbleweeds, which can produce non-negligible gravity gradient signals.
The general formula for the spurious gravitational-wave signal produced by a moving object is:
$$\stackrel{~}{h}(f)=\frac{GM}{4\pi ^2Lf^2}_{\mathrm{}}^{\mathrm{}}\frac{x(t)}{r^3(t)}e^{2\pi ift}𝑑t,$$
(30)
where $`M`$ is the mass of the object, $`r(t)`$ is its distance from the test mass as a function of time, and $`x(t)`$ is its distance from the end mass in the direction parallel to the interferometer arm. For an object traveling parallel to the ground in a straight line at speed $`v`$, Eq. (30) becomes \[as in Eq. (14)\]:
$`\stackrel{~}{h}(f)`$ $`=`$ $`{\displaystyle \frac{GM}{Lv^2\pi f}}[K_0\left({\displaystyle \frac{2\pi fr_{\mathrm{min}}}{v}}\right)\mathrm{cos}\theta `$ (31)
$`iK_1\left({\displaystyle \frac{2\pi fr_{\mathrm{min}}}{v}}\right)\mathrm{sin}\theta \mathrm{cos}\psi ]e^{2\pi ift_0},`$
where $`\theta `$ is the angle between the line of motion and the interferometer arm, $`r_{\mathrm{min}}`$ is the distance of closest approach between the object and the test mass, $`\psi `$ is the angle projecting this distance onto the ground, and $`t_0`$ is the time of closest approach. $`K_0`$ and $`K_1`$ are modified Bessel functions of the second kind of order 0 and 1, respectively.
Under moderately windy conditions (wind speeds up to 15m/s or so), tumbleweeds at the Hanford LIGO facility will bounce along the ground at 5–10m/s. In stronger winds, the tumbleweeds become airborne, with speeds approaching the wind speed; they can fly right over the LIGO buildings, or impact with considerable force. The same may be true of wind-borne debris at other interferometer facilities. However, the value of $`r_{\mathrm{min}}`$ is usually at least 5m, so for frequencies above 10 Hz one has $`2\pi fr_{\mathrm{min}}/v10`$ even for very strong winds ($`v30`$m/s). In this regime the Bessel functions are exponentially damped, so these objects will not produce significant gravity gradient signals simply by blowing past the instrument. A rifle bullet, on the other hand, might be moving around 1000m/s, putting us in the small-argument regime of the Bessel functions, where $`K_0(x)\mathrm{ln}(x)`$ and $`K_1(x)x^1`$. Taking the most dangerous geometry $`\theta =\pi /2`$, $`\psi =0`$, and assuming a bullet mass of around 5 grammes, this gives a dimensionless signal amplitude near 10 Hz of:
$`2|f\stackrel{~}{h}(f)|`$ $``$ $`1.6\times 10^{22}\left({\displaystyle \frac{M}{5\mathrm{g}}}\right)\left({\displaystyle \frac{1000\mathrm{m}/\mathrm{s}}{v}}\right)`$ (32)
$`\times \left({\displaystyle \frac{5\mathrm{m}}{r_{\mathrm{min}}}}\right)\left({\displaystyle \frac{10\mathrm{H}\mathrm{z}}{f}}\right).`$
This gives a signal-to-noise ratio of about 1 at 10 Hz, a bit below the detectable threshold. In fact, even if one fine-tunes the bullet speed $`v`$, the largest signal-to-noise ratio that one can get at 10 Hz is about 2, for a speed of around 250m/s. Since events with signals less than about 5 times the noise will probably be ignored in any case, one can conclude that objects flying past a test mass are not likely to be serious sources spurious events.
If an object does not pass smoothly by the interferometer but instead collides with an end station, the signal at $`10`$ Hz can be large even for slow-moving objects: it is the deceleration time, not the end-station-crossing time, that sets the frequency scale of the signal. Suppose an object collides end-on with the end station at a speed $`v`$, coming to a stop within a distance $`d`$ at constant acceleration. Let $`t=0`$ denote the time that the object comes to rest. The motion of the object is then given by:
$$r(t)=x(t)=\{\begin{array}{cc}r_{\mathrm{min}}dvt\hfill & t2d/v\\ r_{\mathrm{min}}+(vt)^2/4d\hfill & 2d/vt0\\ r_{\mathrm{min}}\hfill & t0\end{array}.$$
(33)
The Fourier transform of $`x(t)/r^3(t)=1/r^2(t)`$ is tricky to do analytically, so I have relied on numerical fast Fourier transforms, and then made approximate analytic fits to the result. However, one can qualitatively predict the shape of the signal in frequency space. The function $`1/r^2(t)`$ starts out near zero and then slowly rises over a timescale $`r_{\mathrm{min}}/v`$ to a value $`r_{\mathrm{min}}^2`$, then quickly levels off at that value over a timescale $`d/v`$. So on frequency scales $`v/r_{\mathrm{min}}`$, $`1/r^2(t)`$ looks like a step function, whose Fourier transform goes as $`1/f`$. On frequency scales $`v/r_{\mathrm{min}}`$ but $`v/d`$, one sees the deceleration as a cusp (discontinuous first derivative), giving a Fourier transform that goes as $`1/f^2`$. On frequency scales $`v/d`$, the deceleration appears smooth, but the onset of deceleration in sudden, giving a $`1/f^3`$ behaviour. The following gives a good fit to the numerical Fourier transform:
$`\left|{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{e^{2\pi ift}}{r^2(t)}}𝑑t\right|`$ $``$ $`{\displaystyle \frac{1}{vr_{\mathrm{min}}}}[5.9\left({\displaystyle \frac{fr_{\mathrm{min}}}{v}}\right)+14\left({\displaystyle \frac{fr_{\mathrm{min}}}{v}}\right)^2`$ (34)
$`+59\left({\displaystyle \frac{f^3r_{\mathrm{min}}^2d}{v^3}}\right)]^1.`$
The two breakpoints separating the three branches are at frequencies $`f_1=0.4v/r_{\mathrm{min}}`$ and $`f_2=0.24v/d`$. More precisely, since the deceleration occurs over a well-defined time $`2d/v`$, the third branch of the Fourier transform is oscillatory with nodes every $`v/2d`$ in frequency space; the functional fit in Eq. (34) is an envelope containing these oscillations.
For a 5g bullet striking an end station at 1000m/s, the signal at 10 Hz is dominated by the low-frequency tail regardless of stopping distance. The dimensionless signal amplitude is then:
$$2|f\stackrel{~}{h}(f)|3\times 10^{22}\left(\frac{M}{5\mathrm{g}}\right)\left(\frac{10\mathrm{H}\mathrm{z}}{f}\right)^2\left(\frac{5\mathrm{m}}{r_{\mathrm{min}}}\right)^2.$$
(35)
The signal amplitude is about the same as for a passing bullet, although the dependence on $`f`$ (and $`r_{\mathrm{min}}`$ and $`v`$) is different. For a tumbleweed or other wind-borne object, by contrast, the sudden deceleration can create significant high-frequency noise. A typical tumbleweed at Hanford has a mass of $`M=0.1`$kg and a diameter of 0.4m, and can be compressed by about half that amount ($`d=0.2`$m). Larger weeds can be twice as large in diameter, putting their masses in the 0.5–1kg range Raab\_F:1999 . For moderate to high wind speeds ($`v=10`$–30m/s), the signal at frequencies above 10 Hz is dominated by the second branch of Eq. (34), giving a dimensionless signal amplitude of:
$`2|f\stackrel{~}{h}(f)|`$ $``$ $`5\times 10^{21}\left({\displaystyle \frac{M}{1\mathrm{k}\mathrm{g}}}\right)\left({\displaystyle \frac{v}{10\mathrm{m}/\mathrm{s}}}\right)`$ (36)
$`\times \left({\displaystyle \frac{10\mathrm{H}\mathrm{z}}{f}}\right)^3\left({\displaystyle \frac{5\mathrm{m}}{r_{\mathrm{min}}}}\right)^3.`$
Thus even a “typical” 0.1kg weed at these speeds will produce a signal-to-noise ratio of around 4 at 10 Hz in advanced LIGO interferometers, which is in danger of being interpreted as a real gravitational-wave event. In a more extreme case, a 1kg tumbleweed blown airborne by a strong 30m/s wind will produce a signal 100 times higher than the noise at 10 Hz, which is easily detectable.
Fig. 9 shows the signal spectra for the objects discussed above, plotted against the dimensionless noise amplitude expected in advanced LIGO detectors.
Since tumbleweeds are a potential source of spurious detectable events, one should consider ways to reduce the tumbleweed gravity gradient noise. Fortunately, the signal goes as $`r_{\mathrm{min}}^3`$, so a simple fence preventing the weeds from approaching the end station should be sufficient. A fence 30m out from the end station will reduce the signal-to-noise ratio to 1 for tumbleweed masses up to 1kg and speeds up to 30m/s, reducing the risk of spurious events.
## VI Conclusions
This paper has studied two sources of background gravity gradient noise, from infrasonic atmospheric pressure waves and from wind-advected temperature perturbations, in order to determine whether they constitute a limiting noise floor for interferometric gravitational-wave detectors—in particular, for the “advanced” LIGO detectors projected in Abramovici\_A:1992 . The paper also analyzed two sources of gravity gradient signals, from transient atmospheric shockwaves and from massive airborne bodies, to determine whether they would constitute detectable spurious events in these interferometers. The following summarizes the results and suggests possible further work that may need to be done.
Current estimates suggest that infrasonic pressure waves will not be a significant source of gravity gradient noise, being over two orders of magnitude below the advanced LIGO noise floor at 10 Hz. Nonetheless, these estimates are not based on actual noise measurements at an interferometer site, so infrasound measurements at these sites are recommended to confirm them. Further empirical studies might also analyze the specific effects of building shapes and of infrasound coherence lengths on the noise spectrum, although these refinements would likely only serve to reduce noise estimates above 15 Hz or so.
Wind-advected temperature perturbations, although the dominant source of atmospheric density fluctuations, do not produce significant high-frequency gravity gradient noise, due to the long times that any particular pocket of warm or cool air spends in the vicinity of an interferometer test mass. A possible exception is when the airflow forms vortices around the interferometer buildings, since this will produce a noise spectrum peaked around the typical vortex circulation frequencies near the test mass. The current crude analysis of these effects suggests that the noise is still an order of magnitude below the advanced LIGO noise floor at 10 Hz even in the worst-case scenarios, but the model could be improved significantly. Numerical models of the airflow and of temperature perturbations near an interferometer building may be required to settle this issue definitively.
Gravity gradients from atmospheric shockwaves are potentially serious sources of spurious signals in interferometric gravitational-wave detectors. For instance, the sonic boom from a supersonic aircraft overflying an advanced LIGO detector could produce signal-to-noise ratios of several hundred. Although such overflights are expected to be rare or nonexistent, they point out the potential seriousness of shocks from weaker or more distant sources, even if the signals are several orders of magnitude smaller. It is therefore *strongly* recommended that advanced interferometric detectors include infrasonic detectors as environmental monitors. Such sensors could easily be employed to veto spurious atmospheric gravity-gradient events.
Gravity gradients from wind-borne objects such as tumbleweeds are another possible source of spurious events in gravitational-wave detectors, if these objects are allowed to collide with the buildings housing the interferometers. Fences or other structures should be used to keep these objects at least 30 metres from the test masses, in order to eliminate the risk of spurious signals.
Obviously the number of things that can affect interferometric detectors through gravity gradient forces is immense; I have considered here only the few sources that I considered the most worrisome. I encourage other researchers to consider the implications of this often-neglected effect.
###### Acknowledgements.
This work was supported by NSF grant PHY-9424337. I would like to thank Kip Thorne for his help and support throughout this project, as well as Brad Sturtevant, Brian Kern, and Scott Hughes for their insightful discussions about atmospheric phenomena. I am especially grateful to Fred Raab for collecting the vital statistics of Hanford tumbleweeds.
## Appendix A The temperature noise spectrum
This appendix presents a more rigorous mathematical derivation of Eq. (10) used in Sec. III.
Consider a time-varying field of temperature perturbations $`\delta T(𝐫,t)`$ about some average temperature $`T`$. This produces a gravitational perturbation $`g_x(t)=𝑑VG\rho xr^3(\delta T/T)`$, where $`\rho `$ is the average air density, and the $`x`$-axis is along the interferometer arm. The spectral density of gravity gradient noise $`S_g(|f|)`$ is given by twice the Fourier transform of the gravity autocorrelation $`C_g(\tau )=g(t)g(t+\tau )`$. Thus:
$`S_g(|f|)`$ $`=`$ $`2\left({\displaystyle \frac{G\rho }{T}}\right)^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau {\displaystyle 𝑑V𝑑V^{}\frac{xx^{}}{r^3(r^{})^3}}`$ (37)
$`\delta T(𝐫,t)\delta T(𝐫^{},t+\tau )e^{2\pi if\tau }.`$
The temperature noise measured at a point $`𝐫_0`$, on the other hand, is given simply by:
$$S_T(|f|)=_{\mathrm{}}^{\mathrm{}}𝑑\tau \delta T(𝐫_0,t)\delta T(𝐫_0,t+\tau )e^{2\pi if\tau }.$$
(38)
On sufficiently small scales, the temperature perturbations in the Earth’s turbulent boundary layer can be treated as homogeneous and isotropic. The expected squared temperature difference between two points is then a function only of their separation: $`[\delta T(𝐫)\delta T(𝐫+\mathrm{\Delta }𝐫)]^2=D_T(\mathrm{\Delta }𝐫)`$. The function $`D_T(\mathrm{\Delta }r)`$ is called the *temperature structure function* of the atmosphere, and for small $`\mathrm{\Delta }r`$ reduces to a power law $`D_T(r)=c_T^2\mathrm{\Delta }r^p`$. If a wind with speed $`v`$ blows these perturbations past a measuring station, the temperature autocorrelation is $`\delta T(𝐫_0,t)\delta T(𝐫_0,t+\tau )=\sigma _T^2(1/2)c_T^2(v\tau )^p`$ for small $`\tau `$, where $`\sigma _T^2`$ is the mean squared temperature fluctuation. This results in a high-frequency power law tail:
$$S_T(|f|)=c_T^2v^p(2\pi f)^{(p+1)}\mathrm{\Gamma }(p+1)\mathrm{sin}(p\pi /2).$$
(39)
Turbulent mixing theory, as well as micrometeorological measurements of $`S_T(|f|)`$, show that the value of $`p`$ is normally 2/3, characteristic of a type of turbulence known as Kolmogorov turbulence. See, for example, Busch\_N:1972 for discussion of this type of turbulent mixing, also Coulman\_C:1985 and references therein.
We are interested in $`S_g(|f|)`$, which is somewhat trickier to calculate than $`S_T(|f|)`$, since it involves a correlation between points separated in space as well as time. However, chaotic turbulence will almost certainly destroy high-frequency correlations between widely separated points, so the high-frequency behaviour of $`S_g(|f|)`$ will come from correlations between nearby points. That is, the high-frequency support of $`\delta T(𝐫,t)\delta T(𝐫^{},t+\tau )`$ will come from those points $`𝐫^{}`$ at time $`t+\tau `$ whose fluid elements were near $`𝐫`$ at time $`t`$.
Consider two fluid elements moving along paths $`S`$ and $`S^{}`$ passing through $`𝐫`$ and $`𝐫^{}`$, respectively. This is shown schematically in Fig. 10. Let $`r_0`$ be the distance from $`𝐫`$ to the nearest point on $`S^{}`$, and let $`\tau _0`$ be the time it would take a pocket of air at $`𝐫^{}`$ to be carried to this point on $`S^{}`$. In order for these points $`𝐫`$ and $`𝐫^{}`$ to contribute to the high-frequency component of the spectrum, the distance $`r_0`$ must be fairly small, of the order $`v/f`$ where $`v`$ is the wind speed past $`𝐫`$. I treat the streamlines as relatively straight over these scales, in which case the temperature perturbations at $`(𝐫,t)`$ and $`(𝐫^{},t+\tau )`$ correspond to physical pockets of air separated by a distance $`\sqrt{r_0^2+v(𝐫)^2(\tau \tau _0)^2}`$, for $`\tau `$ near $`\tau _0`$. Assuming that the $`\tau `$-dependence of the correlation function is due entirely to this advection, the correlation function can be written explicitly as:
$$\delta T(𝐫,t)\delta T(𝐫^{},t+\tau )=\sigma _T^2\frac{c_T^2}{2}\left[r_0^2+v^2(\tau \tau _0)^2\right]^{p/2},$$
(40)
where typically $`p2/3`$.
This term contains the entire $`\tau `$-dependence of the gravity perturbation in Eq. (37), so the first integral I do is the Fourier integral over $`\tau `$. This integral has the form $`_{\mathrm{}}^{\mathrm{}}(\beta ^2+x^2)^{\nu 1/2}e^{iax}𝑑x=2\pi ^{1/2}(2\beta /a)^\nu \mathrm{cos}(\pi \nu )\mathrm{\Gamma }(\nu +1/2)K_\nu (a\beta )`$. Formally the integral diverges for $`\nu 1/2`$, but the closed-form expression remains approximately correct for large $`a`$ provided the integrand is cut off smoothly for large $`x1/a`$. Physically this corresponds to the fact that a smooth, large-scale cutoff of the temperature correlations in Eq. (40) will not affect the high-frequency component of the temperature noise. For horizontal winds near the ground the spatial correlation function is cut off on horizontal distance scales of $`50`$ times the fluid elements’ altitude $`z`$ above ground, giving a low-frequency cutoff around $`0.02v/z`$ (Fig. 1.A4 of Busch\_N:1972 ). Typically this will be below the relevant frequency range for interferometric detectors; I ignore it to obtain pessimistic (upper-limit) noise estimates. The high-frequency noise tail is then:
$`S_g(|f|)`$ $`=`$ $`\left({\displaystyle \frac{G\rho }{T}}\right)^2{\displaystyle }dV{\displaystyle }dV^{}[{\displaystyle \frac{xx^{}}{r^3(r^{})^3}}`$ (41)
$`\times `$ $`c_T^2v^p(\sqrt{2}\pi f)^{(p+1)}a(p)\left({\displaystyle \frac{2\pi fr_0}{v}}\right)^{(p+1)/2}`$
$`\times K_{(p+1)/2}(2\pi fr_0/v)e^{2\pi if\tau _0}],`$
where $`a(p)=2\pi ^{1/2}\mathrm{cos}(\pi [p+1]/2)\mathrm{\Gamma }([p+2]/2)0.873`$ for $`p=2/3`$.
Next is the integral over $`dV^{}`$. The exponential decay of the Bessel function $`K_{(p+1)/2}`$ restricts the support of this integral to values $`r_0v/2\pi f`$, representing the fact that high-frequency fluctuations can only arise from the rapid change of the spatial correlation function over small lengthscales. This range in $`r_0`$ defines a narrow bundle of streamlines $`S^{}`$ about the streamline $`S`$ passing through $`𝐫`$. I assume that the size of this bundle is less than the distance from the test mass to the bundle, so that the values of $`x^{}`$ and $`r^{}`$ on a given $`S^{}`$ can be replaced with the nearby values on $`S`$. Now for points near $`𝐫`$ the volume element $`dV^{}`$ can be written in terms of the new parameters $`r_0`$ and $`\tau _0`$ as $`dV^{}=2\pi r_0dr_0v(𝐫)d\tau _0`$. Since the airflow, being very subsonic, is nearly incompressible, the volume element retains this form for all points along the bundle: if $`v(𝐫^{})`$ decreases below $`v(𝐫)`$, for instance, the length element $`v(𝐫^{})d\tau _0`$ will decrease, but the cross-section of the bundle (i.e., the relevant range of $`2\pi r_0dr_0`$) will increase to compensate. Plugging in this volume element and integrating over $`r_0`$, one obtains:
$`S_g(|f|)`$ $`=`$ $`\left({\displaystyle \frac{G\rho }{T}}\right)^2{\displaystyle }dV{\displaystyle }vd\tau _0[{\displaystyle \frac{xx^{}}{r^3(r^{})^3}}c_T^2v^p(\pi f)^{(p+1)}`$ (42)
$`\times `$ $`a(p)2\pi \mathrm{\Gamma }([p+3]/2)\left({\displaystyle \frac{v}{2\pi f}}\right)^2e^{2\pi if\tau _0}].`$
Let $`t^{}`$ be a new time coordinate denoting the time it takes for an air pocket to reach $`𝐫^{}`$ from some fixed reference plane that crosses all streamlines orthogonally, and $`t`$ be the corresponding coordinate for the point $`𝐫`$. Then $`\tau _0`$ is just $`t^{}t`$, and the volume element $`dV`$ can be written as $`wdtdA`$, where $`dA`$ is a cross-sectional area element on the reference plane, and $`w`$ is the wind speed across that area element. Plugging this in, and ignoring the spatial separation between $`S`$ and $`S^{}`$, one obtains:
$`S_g(|f|)`$ $`=`$ $`\left({\displaystyle \frac{G\rho }{T}}\right)^2c_T^2(\pi f)^{(p+3)}a(p)(\pi /2)\mathrm{\Gamma }([p+3]/2)`$ (43)
$`\times `$ $`{\displaystyle _{\{S\}}}w𝑑A\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{x}{r^3}}v^{p+3}e^{2\pi ift}𝑑t\right)`$
$`\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{x^{}}{(r^{})^3}}e^{2\pi ift}𝑑t^{}\right),`$
where I have used $`_{\{S\}}`$ to denote an integral over the entire reference plane; i.e., over all streamlines $`S`$.
The noise in the gravitational-wave signal $`h(t)`$ due to the gradients at a given test mass is $`S_h(|f|)=(2\pi f)^4S_g(|f|)/L^2`$, where $`L`$ is the interferometer arm length. The noise at each test mass adds incoherently to the total signal. Combining these with Eq. (43) yields the result given in Eq. (10).
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# 1 Introduction
## 1 Introduction
The idea of introducing the statistical description into the soliton theory has its origin in the work of Zakharov on kinetic equation for the rarefied soliton gas. In this work, a new object, an infinite sequence of KdV solitons on the $`x`$-axis was considered. An accurate limit for the $`N`$\- soliton solution as $`N\mathrm{}`$ is , however, very nontrivial and has been investigated so far only for the case of a very special phase distribution , which yields the reflectionless potential decaying at one of infinities $`x\pm \mathrm{}`$. This kind of infinite - soliton solutions does not provide directly the spatial uniformity which is necessary for description of the soliton gas in thermodynamically equilibrium state.
Another view on the problem was proposed by Lax who noticed that ‘the weak limits of oscillatory sequences of dispersive compressible flows show a remarkable resemblance to ensemble averages of classical turbulence theory’. To constitute such a deterministic analogue of turbulence one has to construct the correspondent ensemble of flows and provide it with the physically plausible measure which form a stochastic process. It is clear from the very beginning that for description of the ‘soliton turbulence’ we are interested in the stationary ergodic processes.
The consistent way of application of stochastic description to the soliton theory was proposed in , where the notion of stochastic soliton lattice (SSL) was introduced. The basic idea of such a description lies in the fact that the finite-gap solutions of completely integrable equations are almost periodic functions and posess, therefore, their natural stochastic structure determined by the compact shift group with the uniform (Haar) measure.
In this work, we investigate $`N`$-gap stochastic soliton lattices $`q(x)=\nu _N(x)`$ of the KdV equation on a special band-gap scaling of the spectral Riemann surface of the complex parameter $`E`$ when the width of gaps is $`O(1/N)`$ while the bands are exponentially narrow $`\mathrm{exp}(N)`$, $`N1`$. Along with this, the small parameter $`1/N`$ in our consideration does not appear explicitely in the Schrödinger equation $`(d_{xx}^2+q(x))\varphi =E\varphi `$.
The basic results of the present work are:
* the process $`\zeta (x)=\underset{N\mathrm{}}{lim}\nu _N(x)`$ exists. It is an ergodic stationary random process each realization of which satisfies the KdV equation. The phase space of $`\zeta (x)`$ is a one-dimensional space with random Poisson measure.
* the considered limit is of a thermodynamic type (the integrated density of states remains finite as $`N\mathrm{}`$ for the chosen scaling of the band-gap stucture).
* as $`\underset{N\mathrm{}}{lim}\frac{band}{gap}0`$ the thermodynamic limit of stochastic soliton lattices $`\zeta (x)`$ can be interpreted as a soliton turbulence.
* we calculate the rotation number and some ensemble averages for the soliton turbulence.
* we show that the zero-density limit of the soliton turbulence yields the Frish – Lloyd potential of the quantum theory of disordered systems , .
## 2 Basic Preliminaries
The soliton lattice (SL) ($`N`$-gap potential of the Schrödinger equation) is given by the known Its – Matveev formula (see for instance )
$$u_N(x;𝐫)=C(𝐫)2_{xx}^2\mathrm{log}\mathrm{\Theta }[𝐳(𝐱)|B(𝐫)],$$
(1)
$$x𝐑,𝐫=(r_1,\mathrm{},r_{2N+1}),r_1<r_2<\mathrm{}<r_{2N}<r_{2N+1}.$$
Here
$$\mathrm{\Theta }[𝐳|B(𝐫)]=\underset{𝐦}{}\mathrm{exp}\{\pi i[2(𝐦,𝐳)+(𝐦,B𝐦)]\},$$
(2)
$$𝐦=(m_1,\mathrm{},m_N)𝐙^N,𝐳𝐂^N$$
is the Jacobi theta-function of the hyperelliptic Riemann surface of genus $`N`$
$$\mu ^2=\underset{j=1}{\overset{2N+1}{}}(Er_j)R_{2N+1}(E,𝐫)$$
(3)
with cuts along the bands.
The Riemann matrix $`B(𝐫)`$ and the constant $`C(𝐫)`$ are expressed in terms of the basis holomorphic differentials $`\psi _j`$ :
$$B_{ij}=\underset{\beta _j}{}\psi _i,C(𝐫)=\underset{j=1}{\overset{N}{}}r_j2\underset{j=1}{\overset{N}{}}\underset{\alpha _j}{}E\psi _j,$$
(4)
Here
$$\psi _j=\underset{k=0}{\overset{N1}{}}a_{jk}\frac{E^k}{\sqrt{R_{2N+1}(𝐫,E)}}dE,$$
(5)
and dependence $`a_{jk}(𝐫)`$ is given by the normalization
$$\underset{\alpha _k}{}\psi _j=\delta _{jk},$$
(6)
where the $`\alpha `$-cycles surround the bands clockwise, and the $`\beta `$\- cycles are canonically conjugated to $`\alpha `$’s such that the contour $`\beta _j`$ starts from the $`j`$-th cut, then goes to $`+\mathrm{}`$ and returns on the lower cut.
The imaginary phases $`𝐳(x)`$ are given by the formula
$$𝐳=2i𝐚_{N1}x+𝐝,$$
(7)
where d is the initial imaginary phase vector.
Using the substitution
$$𝐳=\frac{1}{2\pi }B𝐲$$
(8)
we rewrite ( 1) in the form with real phases $`𝐲`$
$$u_N(x|𝐫)=u_N(y_1(x),\mathrm{},y_N(x)|𝐫),$$
(9)
$$y_j(x)=k_jx+f_j,f_j(mod2\pi ),$$
(10)
$$𝐤=4\pi iB^1𝐚_{N1}.$$
(11)
Here $`𝐤=𝐤(𝐫)=(k_1,\mathrm{},k_N)`$ is the wave number vector, and $`𝐟=(f_1,\mathrm{},f_N)`$ is the initial (angle) phase vector, $`\pi <f_j\pi `$. Note that in (9)
$$u_N(y_1,\mathrm{},y_j+2\pi ,\mathrm{},y_N|𝐫)=u_N(y_1,\mathrm{},y_j,\mathrm{},y_N|𝐫),$$
(12)
that is $`u_N(x|𝐫)`$ is $`N`$-quasiperiodic in $`x`$ , (hereafter we consider only incommensurate $`k_j,j=1,\mathrm{},N)`$. As for any almost periodic function, we have for $`u_N(x|𝐫)`$ the Fourier representation
$$u_N(x|𝐫)=\underset{j}{}c_je^{i(l_jx+h_j)},$$
(13)
where $`c_j,l_j,h_j`$ are real, $`\pi <h_j\pi `$. ($`l_j𝐌`$, where $`𝐌`$ is a frequency-module and $`h_j`$ are the integer linear combinations of $`f_j(mod2\pi )`$).
The KdV evolution of (9) is isospectral and is described by the linear motion of the phases on the Jacobian:
$$u_N(x,t|𝐫)=(y_1(x,t),\mathrm{},y_N(x,t)|𝐫),$$
(14)
$$y_j(x,t)=k_jx+\omega _jt+f_j,$$
(15)
where the frequency vector
$$\omega =\omega (𝐫)=(\omega _1,\mathrm{},\omega _N)=8\pi iB^1(𝐚_{N1}\underset{j=1}{\overset{2N+1}{}}r_j+2𝐚_{N2}).$$
(16)
## 3 Stochastic Soliton Lattices
It is well known that any almost periodic function generates the stochastic stationary process (see , ).
Definition 1., The stochastic process generated by SL $`u_N(x|𝐫)`$ we call Stochastic Soliton Lattice (SSL) and denote as $`\nu _N(x|𝐫)`$.
The general construction of $`\nu _N(x|𝐫)`$ admits a very simple and clear description. The realization set $`\mathrm{\Omega }`$ of it consists of functions (9), (10) where $`𝐟Tor^N`$ ; $`Tor^N`$ is $`N`$-dimensional torus $`(\pi ,\pi ]^N`$. The probability measure $`d\mu `$ is the uniform (Lebesque) measure on the torus. It corresponds to the description of $`\nu _N(x|𝐫)`$ following from(9),(10):
$$\nu _N(x|𝐫)=u_N(\mathrm{}k_jx+\varphi _j\mathrm{}|𝐫)$$
(17)
where $`\varphi _1,\mathrm{},\varphi _N`$ are independent random values uniformly distributed on $`(\pi ,\pi ]`$, that is $`\varphi =(\varphi _1,\mathrm{},\varphi _N)`$ is uniformly distributed on $`Tor^N`$. As $`k_j`$ are incommensurate then $`\nu _N(x|𝐫)`$ is an ergodic process .
As $`\nu _N(x|𝐫)`$ is the stationary process then it has the Stone - Kolmogorov spectral decomposition ; due to the ergodicity this decomposition has the form (cf. (13))
$$\nu _N(x|𝐫)=\underset{j}{}c_je^{i(l_jx+\theta _j)},$$
(18)
where $`\theta _j`$ are uniformly distributed on $`(\pi ,\pi ]`$ noncorrelated random values .
The well known formula (Bochner - Khintchin) gives us the covariance function $`K(h)`$ of the stationary process $`\nu _N(x|𝐫)`$:
$$K(h)\widehat{\nu }_N(x|𝐫)\widehat{\nu }_N(x+h|𝐫)=\underset{j}{}|c_j|^2e^{ił_jh}.$$
(19)
Here $`\widehat{\xi }(x)\xi (x)\xi `$ is the centered process.
Theorem. Consider the KdV equation $`u_t6uu_x+u_{xxx}=0`$ as the equation describing the evolution in the phase space of stationary processes. Let the initial data has the form of the SSL:
$$u(x,0)=\nu _N(x|𝐫).$$
(20)
Then the solution of the KdV is
$$u(x,t)=\nu _N(x|𝐫)=u(x,0),$$
(21)
that is $`u(x,0)=\nu _N(x|𝐫)`$ is a stationary point.
Proof . The evolution of realizations (9) ,(10) is described by (14) ,(15). For any moment $`t`$ one can introduce the new ‘initial phase’ $`f_j^{}=\omega _jt+f_j`$ which is also uniformly distributed on $`Tor^N`$. Therefore, the KdV evolution changes neither realization set $`\mathrm{\Omega }`$ nor the probability measure. Q.E.D.
Now we present some expressions for the ensemble averages which will be needed in the future. If $`Q(f)`$ is an arbitrary smooth function then the mean value $`Q(\xi )`$ of the stochastic process $`\xi `$ is
$$Q(\xi )=\underset{\mathrm{\Omega }}{}𝑑\mu Q(f),f\mathrm{\Omega }.$$
(22)
(We recall that $`\xi (x)\{\mathrm{\Omega }=\{f(x)\},B,\mu \}`$, and $`B`$ is some $`\sigma `$-algebra of measurable Borel sets of $`\mathrm{\Omega }`$).
For $`\nu _N(x|𝐫)`$ we have
$$Q(\nu _N(x|𝐫)=\frac{1}{(2\pi )^N}\underset{\pi }{\overset{\pi }{}}\mathrm{}\underset{\pi }{\overset{\pi }{}}d\varphi _1\mathrm{}d\varphi _NQ(\nu _N(x|𝐫))=\underset{L\mathrm{}}{lim}\frac{1}{L}\underset{0}{\overset{L}{}}Q(\nu _N(x|𝐫)),$$
(23)
which constitutes ergodicity of SSL. Direct calculation using (1),(2) gives surprisingly simple formulas for two first moments of the value $`\nu _N(x|𝐫)`$ Namely,
$$\nu _N(x|𝐫)=C(𝐫)\frac{1}{\pi i}(𝐤,B𝐤),\nu _N^2(x|𝐫)=\frac{1}{3\pi i}(\omega ,B𝐤),$$
(24)
and the vectors $`𝐤(𝐫),\omega (𝐫)`$ are given by (11) , (16). It should be noted that expressions (24) are obtained for the particular choice of the canonical basis of cycles (see (6)). However, namely this normalizatiton is preferrable for our consideration.
## 4 Rotation Number in Stochastic Soliton Lattices
Consider the Schrödinger equation with an almost periodic potential $`q(x)`$
$$(_{xx}^2+q(x))\varphi =E\varphi ,x𝐑.$$
(25)
The potential $`q(x)`$ has an important characteristics, the rotation number, which is defined for real $`E`$ as \[JM\]
$$\alpha (E)=\underset{x\mathrm{}}{lim}\frac{1}{x}\mathrm{arg}(\varphi ^{}(x,E)+i\varphi (x,E)).$$
(26)
We will use also the integrated density of states $`𝒩(E)`$ which is connected with the rotation number by a simple relation (here we improve some inaccuracies in , ):
$$𝒩(E)=\frac{1}{\pi }\alpha (E).$$
(27)
If $`q(x)`$ is N-gap potential , $`q(x)=u_N(x|𝐫)`$, then
$$\alpha _N(E)=Re\underset{\mathrm{}}{\overset{E}{}}𝑑p(E^{}),$$
(28)
where $`dp(E)`$ is the quasimomentum differential , :
$$dp(E)=\frac{E^N+b_{N1}E^{N1}+\mathrm{}+b_0}{\sqrt{R_{2N+1}(E;𝐫)}},\underset{\beta _j}{}𝑑p(E)=0,j=1,\mathrm{},N.$$
(29)
The integrals of $`dp`$ over the $`\alpha `$ \- cycles are known , to give the components of the wave number vector (11)
$$\underset{\alpha _j}{}𝑑p(E)=k_j.$$
(30)
Then,
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{M(E)}{}}}k_j\text{if}E\text{gap}_M`$
$`\alpha _N(E)`$ $`=`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{M(E)}{}}}k_j+{\displaystyle \underset{r_{2M1}}{\overset{E}{}}}𝑑p(E^{})\text{if}E\text{band}_M,`$
where $`M(E)`$ is the number of the band nearest to $`E`$ from the left $`M(E)N,M(0)=N,M(1)=0`$.
We note that the formula (4) gives an effectivization for the case of finite-gap potentials of the statement in that the values of the function $`2\alpha (E)`$ if $`E\{\text{gap}\}`$ belong to the frequency module of the almost periodic potential.
As the SL (9), (10) is the quasiperiodic function both in $`x`$ and $`t`$ one can also formally introduce the temporal rotation number for the SSL by
$$\alpha _N^t(E)=Re\underset{\mathrm{}}{\overset{E}{}}𝑑q(E^{})$$
(32)
where $`dq(E)`$ is the quasienergy differential , ($`dq(E)=(E^{N+1}+c_NE^N+\mathrm{}+c_0)dE/\sqrt{R_{2N+1}(E;𝐫)}`$, $`\underset{\beta _j}{}𝑑q(E)=0,\underset{\alpha _j}{}𝑑q(E)=\omega _j`$, where $`\omega _j`$ are the frequences (16)).
Since the value of $`\alpha _N(E)`$ is the same for any realization from the set $`\mathrm{\Omega }`$ from Def.1 then it is the spectral characteristics of the whole SSL $`\nu _N(x|𝐫)`$ (deterministic characteristics of the stochastic process). Due to the Theorem 1 the rotation number does not change under the KdV evolution.
We introduce also the full integrated density of states in the SSL $`\lambda `$ which is a number:
$$\lambda 𝒩_N(\mathrm{})=\frac{1}{2\pi }\underset{j=1}{\overset{N}{}}k_j.$$
(33)
One can see that the full integrated density of states has in the SSL the natural meaning of the mean number of waves per unit length.
## 5 Thermodynamic Limit of Stochastic Soliton Lattices
Let the nontrivial (band-gap) part of the spectrum lies in the interval $`(1,0)`$ of the real $`E`$-axis: $`\alpha _N(1)=0,\alpha _N(0)=\pi \lambda `$. For $`N1`$ we consider the following scaling of the band-gap structure
$$\text{gaps}(E)\frac{1}{N}\text{bands}(E)\mathrm{exp}(N)N1,E(1,0).$$
(34)
Using (11) one can show that the scaling (34) implies the following behavior for the wave numbers: $`k_j1/N`$, which provides finiteness of the rotation number $`\alpha _N(E)`$ (and of the integrated density of states $`𝒩_N(E)`$) as $`N\mathrm{}`$:
$$\alpha (E)\underset{N\mathrm{}}{lim}\alpha _N(E)=\underset{N\mathrm{}}{lim}\frac{1}{2}\underset{j=1}{\overset{M(E)}{}}k_j<\mathrm{},MN.$$
(35)
Due to this property we call the scaling (34) and the correspondent limit as $`N\mathrm{}`$ the thermodynamic ones.
It should be noted that the thermodynamic scaling (34) appears when one considers the quasiclassical asymptotics of the spectrum for periodic potentials , , . In contrast to indicated works, however, the small parameter $`1/N`$ in our consideration is contained only in the band-gap structure of the spectrum and does not appear explicitely in the Schrödinger equation (25). As a result, we will show that the limit for the finite-gap potentials on the thermodynamic scaling exists as $`N\mathrm{}`$ in a strong sense. We emphasize the difference of this point from the known Lax – Levermore – Venakides (LLV) approach (see , , ) to the limiting passage $`lim_N\mathrm{}u_N`$ where due to rescaling $`xNx`$ the above limit exists only in a weak sense. The LLV consideration provides information about the slow modulations (macrostructure) appearing in the $`N`$-soliton (or $`N`$-gap) KdV solution under the time evolution. We, on the contrary, are interested in the limiting microstructure of nonmodulated $`N`$-phase wave as $`N\mathrm{}`$. Futhermore, we are going to show that the thermodynamic limit exists for the whole stochastic process (SSL) which implies not only compactness of the realization set but also convergence of the probability measure.
Now we describe the thermodynamic scaling more in detail. Following we introduce the lattice of points
$$1\eta _1>\eta _2>\mathrm{}>\eta _N0,$$
(36)
where
$$\eta _j^2=\frac{1}{2}\left(r_{2j1}+r_{2j}\right)$$
are centers of bands.
We define two continuous functions on the lattice:
1. The normalized density of bands: $`\phi (\eta )d\eta \frac{\text{number of bands in}(E,E+dE)}{N}`$
$$\phi (\eta _j)=\frac{1}{N(\eta _j\eta _{j+1})}+O(\frac{1}{N}),\underset{0}{\overset{1}{}}\phi (\eta )𝑑\eta =1,\eta ^2=E(0,1).$$
(37)
2. The normalized logarithmic band width
$$\gamma (\eta _j)=\frac{1}{N}\mathrm{log}(r_{2j}r_{2j1})+O(\frac{1}{N}).$$
(38)
Now one can easily represent the relationship (35) for the thermodynamic limit of the rotation number in a continuum form:
$$2d\alpha (\eta ^2)=\phi (\eta )k(\eta )d\eta .$$
(39)
where $`k(\eta )`$ is the continu of the wave number vector (11).
For the temporal rotation number $`\alpha _N^t(E)`$ (32) in the thermodynamic limit we have the equation which is analogous to (39):
$$2d\alpha ^t(\eta ^2)=\phi (\eta )\omega (\eta )d\eta .$$
(40)
where $`\omega (\eta )`$ is the continuum limit for the frequency vector (16).
To find a continuum limit for the wave number and the frequency vectors given by (11) and (16) one has to know the continuum limits of the period matrix $`B`$ (4) and of the vectors $`𝐚_{N1}`$ and $`𝐚_{N2}`$ from (5). It is clear that all quantities given on the Riemann surface can be computed now in terms of two functions (densities) $`\phi (\eta )`$ and $`\gamma (\eta )`$ completely describing the band - gap structure.
In particular, the period matrix (4) and the vectors $`𝐚_{N1}`$ and $`(𝐚_{N1}\underset{j=1}{\overset{2N+1}{}}r_j+2𝐚_{N2})`$ take the form :
$$BB(\eta ,\mu )=\frac{i}{\pi }\left(\mathrm{log}\left|\frac{\eta \mu }{\eta +\mu }\right|\phi (\mu )+\gamma (\mu )\delta (\eta \mu )\right),$$
(41)
where $`\delta (x)`$ is the Dirac delta-function;
$$𝐚_{N1}\frac{\eta }{2\pi },𝐚_{N1}\underset{j=1}{\overset{2N+1}{}}r_j+2𝐚_{N2}\frac{2\eta ^3}{\pi }.$$
(42)
Now we represent the general relationship (11) between the wave number vector $`𝐤`$ and the vector $`𝐚_{N1}`$ in the form suitable for the limiting transition
$$𝐤B=4\pi i𝐚_{N1}$$
(43)
Then applying (39), (41), (42) to (43) we arrive at the integral equation for the rotation number in the thermodynamic limit :
$$\frac{1}{\pi }\phi (\eta )d\eta \underset{0}{\overset{1}{}}\mathrm{log}\left|\frac{\eta \mu }{\eta +\mu }\right|d\alpha (\mu ^2)+\gamma (\eta )d\alpha (\eta ^2)=\eta \phi (\eta )d\eta $$
(44)
Similarly, for the thermodynamic limit of the temporal rotation number one gets with the aid of (40), (41), (42) and (16):
$$\frac{1}{\pi }\phi (\eta )d\eta \underset{0}{\overset{1}{}}\mathrm{log}\left|\frac{\eta \mu }{\eta +\mu }\right|d\alpha ^t(\mu ^2)+\gamma (\eta )d\alpha ^t(\eta ^2)=4\eta ^3\phi (\eta )d\eta .$$
(45)
Thus, we have established the existence of the thermodynamic limit for the rotation number in SSL (as a matter of fact, the rotation number, as well as in the finite-gap case, is a deterministic function).
The rotation number is known to be an imaginary part of the Floquet exponent $`w(E)`$ (for the finite-gap potentials $`w(E)=i𝑑p(E)`$) . Then we have with the aid of Herglotz formula:
$$w(E)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{\alpha (z)dz}{Ez},ImE>0.$$
(46)
But following Kotani , if the Floquet exponent satisfies the conditions
$`Rew(E+i0)=0\text{a.e.}\text{on}[0,+\mathrm{})`$ (47)
$`{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}Rew(E+i0)d\alpha (E)=0,`$
then there exists an ergodic stationary random process whose Floquet exponent is $`w(E)`$. One can see that for SSL conditions (47) are satisfied due to the properties of the quasimomentum (29). Therefore the thermodynamic limit of SSL exists.
Remark. It can be shown that the existence of the limit for the rotation number guarantees the existence of the limit for the Weyl function. Then the compactness of the realization set can be deduced from the results of Marchenko .
As $`band/gap0`$ in the thermodynamic limit (see (34)), then it is natural to call the thermodynamic limit $`\zeta (x)`$ of the SSL $`\nu _N(x)`$ a soliton turbulence
$$\zeta (x)=\underset{N\mathrm{}}{lim}\nu _N(x).$$
(48)
We believe that the spectrum of a typical realization of the soliton turbulence is pure point on $`(1,0)`$ but exact proof of this conjecture is missed yet.
The ensemble averages (moments) (22) are the important characteristics of the soliton turbulence. To compute three first moments we make use of the expressions (24) obtained for the finite-gap SSL. These expressions admit direct limiting passage as $`N\mathrm{}`$ in the thermodynamic scaling (34).
Using results of it is not hard to show that the constant $`C(𝐫)`$ (4) in the Its-Matveev formula (1) vanishes in the thermodynamic limit. Then, substituting (39), (40) into (24) we have for the two first moments in soliton turbulence:
$$\zeta =2\lambda \underset{1}{\overset{0}{}}𝑑\stackrel{~}{\alpha }(E)=2\lambda ,$$
(49)
$$\zeta ^2=\frac{4\lambda }{3}\underset{1}{\overset{0}{}}E𝑑\stackrel{~}{\alpha }(E)=\frac{4\lambda }{3}\overline{E}.$$
(50)
Here
$$d\stackrel{~}{\alpha }(E)=\frac{d\alpha }{\pi \lambda },\underset{1}{\overset{0}{}}𝑑\stackrel{~}{\alpha }=1,\overline{f(E)}=\underset{1}{\overset{0}{}}f(E)𝑑\stackrel{~}{\alpha }(E).$$
Thus, we have expressed the ensemble averages in soliton turbulence through the averages over the spectrum with the measure $`d\stackrel{~}{\alpha }(E)`$. One should also note that the Floquet exponent $`w(E)`$ (46) provides the averaged Kruskal integrals as the coefficients in the decomposition as $`E\mathrm{}`$.
## 6 Poissonic Properties of Soliton Turbulence
We study the phase space of the soliton turbulence. We recall that in the finite-gap SSL the phase space is the $`N`$ \- dimensional torus with the uniform (Lebesque) measure. For performing the thermodynamic limit it is convenient to introduce the linear phases instead of the angle phases $`\varphi _j`$ in the SSL (17)
$$l_j\frac{\varphi _j}{k_j},$$
(51)
where $`l_j(j=1,\mathrm{},N)`$ are independent random values uniformly distributed on $`(\frac{\pi }{k_j},\frac{\pi }{k_j}]`$ respectively.
The further consideration is close to standard procedure of the thermodynamic limiting passage in the ergodic theory of ideal gas (see for instance , ). We now turn to a new space with the aid of factorization of the torus by interchange group $`S_N`$. The sets of $`N`$ points $`l_j`$ serve now as the points of the obtained space $`Q_N=Tor^N/S_N`$. The measurable sets in the factorized space consist of realizations having $`0,1,\mathrm{},N`$ linear phases in the interval $`\mathrm{\Delta }`$ of the $`x`$ axis.
We introduce the random value $`\xi _j=\chi _{(0,1)}(l_j)`$ which is the number of hitting of the particular phase $`l_j`$ into the fixed interval $`(0,1)𝐑`$ (we suppose that $`\pi /k_j1`$). The variable $`\xi _j`$ takes two values: $`1`$ and $`0`$ with the probabilities $`p_j(1)=k_j/2\pi `$ and $`p_j(0)=q_j=1p_j=1k_j/2\pi `$ . The generating function $`\phi _j(z)`$ for $`\xi _j`$ is
$$\phi _j(z)=(1p_j)+zp_j$$
(52)
The sum $`\xi ^{(N)}\underset{j=1}{\overset{N}{}}\xi _j`$ is the number of hitting of all linear phases into $`(0,1)`$. As $`\xi _j`$ are independent random values, then the generating function for $`\xi ^{(N)}`$ has the form
$$\phi ^{(N)}(z)=\underset{j=1}{\overset{N}{}}\phi _j(z)=\underset{j=1}{\overset{N}{}}(1+(z1)p_j)=\underset{j=1}{\overset{N}{}}(1+\frac{(z1)k_j}{2\pi }).$$
(53)
On the thermodynamic scaling (34) $`k_j=O(N^1)`$. Then taking the logarithm of (53) we obtain
$$\mathrm{ln}\phi ^{(N)}(z)=(z1)\frac{1}{2\pi }\underset{j=1}{\overset{N}{}}k_j+O(N^1).$$
Therefore taking into account (33) one has
$$\phi ^{(\mathrm{})}(z)=\mathrm{exp}\{(z1)\lambda \}=\underset{n=0}{\overset{\mathrm{}}{}}z^n\left(\frac{e^\lambda \lambda ^n}{n!}\right).$$
(54)
Thus, the right-hand part of (54) is the generating function for the Poisson distribution with the parameter $`\lambda `$ ( which is the full integrated density of states in our case). Therefore, the limiting measure in the configurational space is the random Poisson measure (Poissonic white noise).
We note that the Poissonic white noise can be described as the random collection of the points $`\{l_j\}`$ (the linear phases in our case) on the $`x`$ \- axis such that the distances $`s_k`$ between them are independent random values distributed exponentially with the density
$$f(s)=\lambda \mathrm{exp}(\lambda s).$$
(55)
We also note that the Poisson parameter $`\lambda `$ according to (49) can be expressed through the first moment in the soliton turbulence:
$$\lambda =\frac{\zeta }{2}.$$
The value $`\lambda `$ can be then interpreted as a density of the soliton turbulence.
## 7 The Frish – Lloyd Potential as a Zero-Density Limit of the Soliton Turbulence
We consider the soliton turbulence with the small integrated density of states. We introduce a small parameter $`\epsilon 1`$ and make a renormalization:
$$\alpha (E)=\epsilon \stackrel{~}{\alpha }(E),\lambda =\epsilon \stackrel{~}{\lambda }.$$
(56)
This implies the following transformations for the functions $`\phi (\eta )`$ and $`\gamma (\eta )`$ (37), (38) characterizing the Riemann surface in the thermodynamic scaling
$$\phi (\eta )=\phi (\eta ),\gamma (\eta )=\frac{\stackrel{~}{\gamma }(\eta )}{\epsilon },\stackrel{~}{\gamma }(\eta )=O(1).$$
(57)
We choose $`\stackrel{~}{\gamma }(\eta )=c\eta `$ which corresponds to imposing the periodicity condition : $`k_j=2\pi c^1(\epsilon N^1)`$. Then, it follows from (33), (56) that $`\stackrel{~}{\lambda }=1/c`$.
We also make rescaling of the spatial variables
$$x=\frac{\stackrel{~}{x}}{\epsilon },l_j=\frac{\stackrel{~}{l}_j}{\epsilon }.$$
(58)
Consider now the SSL with a small integrated density of states on the thermodynamic scaling (34). We make use of the Its – Matveev formula (1):
$$\nu _N(\stackrel{~}{x})2\epsilon ^2_{\stackrel{~}{x}\stackrel{~}{x}}^2\mathrm{log}\mathrm{\Theta }_N(\mathrm{},\frac{2ia_{N1,j}(\stackrel{~}{x}\stackrel{~}{l}_j)}{\epsilon },\mathrm{}|B),$$
(59)
where according to (41), (42)
$$B_{jj}\frac{\stackrel{~}{\gamma }(\eta _j)}{\epsilon },B_{ij}\mathrm{log}\left|\frac{\eta _i\eta _j}{\eta _i+\eta _j}\right|,$$
(60)
$$a_{N1,j}\frac{\eta _j}{2\pi }.$$
(61)
An analysis of the explicit expression for the theta-function (2) provided (60), (61) shows that on the set of realizations excluding the set of a small measure one can neglect the contribution of the off-diagonal part of the period matrix $`B`$ into the solution (59) (one suppose the inequality
$$\underset{j=1}{\overset{N}{}}\frac{\stackrel{~}{\gamma }(\eta _j)}{\epsilon }n_j^2\underset{ij=1}{\overset{N}{}}\mathrm{log}\left|\frac{\eta _i\eta _j}{\eta _i+\eta _j}\right|n_in_j,N1.$$
(62)
to be satisfied outside of the indicated set of a small measure).
Then asymptotically in $`N`$, $`\epsilon `$, ($`N1\epsilon 1,N\epsilon 1`$ ) the following factorization is valid :
$$\mathrm{\Theta }_N\underset{j=1}{\overset{N}{}}\mathrm{\Theta }(\eta _j\frac{(\stackrel{~}{x}\stackrel{~}{l}_j)}{\pi \epsilon }|\tau _j),$$
(63)
where $`\mathrm{\Theta }(y)`$ is the one-dimensional theta-function ($`\mathrm{\Theta }_3`$ in standard literature (see for ex. )),
$$\tau _j=\frac{K(1m_j)}{K(m_j)}0,m_j1\frac{1}{\eta _j}\mathrm{exp}(\frac{c}{\epsilon }\eta _j)1.$$
(64)
Here $`K(m)`$ is the complete elliptic integral of the first kind, $`m_j`$ are the ellipticity parameters (do not confuse with the multiintegers in (2), (62) ); $`m_j=1`$ corresponds to the soliton limit.
Consider the rescaled SSL of a small density $`\nu _N(\stackrel{~}{x})/\epsilon ^2`$. Then it follows from (59), (63), (64) that
$$\frac{\nu _N(\stackrel{~}{x})}{\epsilon ^2}2\underset{j=1}{\overset{N}{}}\eta _j^2\frac{sech^2\left(\eta _j\frac{\stackrel{~}{x}\stackrel{~}{l}_j}{\pi \epsilon }\right)}{\epsilon ^2}.$$
(65)
Then passing to a limit
$$\epsilon 0,N\mathrm{},\epsilon N\mathrm{}$$
(66)
one obtains
$$\stackrel{~}{\zeta }(\stackrel{~}{x})=lim\frac{\nu _N(\stackrel{~}{x})}{\epsilon ^2}=4\underset{j}{}a_j^2\delta (\stackrel{~}{x}\stackrel{~}{l}_j),$$
(67)
where, according to the results of Sec. 5, $`\stackrel{~}{l}_j`$ are random points on the $`\stackrel{~}{x}`$ \- axis. The distances $`\stackrel{~}{s}_k`$ between these points are independent random values distributed with the density
$$f(\stackrel{~}{s})=\frac{1}{c}\mathrm{exp}(\frac{\stackrel{~}{s}}{c}).$$
(68)
The amplitudes $`a_j`$ are independent among themselves and independent of $`\stackrel{~}{l}_j`$ random values distributed with the density $`\phi (a)`$. The distribution (67) is so-called complex Poisson white noise; it is called the Frish – Lloyd potential in the quantum theory of disordered systems , .
Aknowledgements
We are grateful to M. Freidlin, P.Miller and S.Novikov for useful discussions.
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# 1 Introduction and summary
## 1 Introduction and summary
Shortly after the first observation of the fractional quantum Hall (fqH) effect , R.B. Laughlin proposed that this phenomenon has its origin in a new state of matter, which is formed by electrons but nevertheless admits excitations of fractional charge . The experimental evidence for the existence of fractionally charged excitations includes the results of shot-noise measurements for charge transport through fqH samples .
The unusual properties of the fqH quasi-particles do not stop at their fractional charge; in addition, the (bulk and edge) quasi-particles exhibit various forms of fractional quantum statistics .
In a previous paper we initiated a program that aims at computing various transport properties of fqH systems in a formalism that makes direct reference to the (fractionally) charged quasi-particles. In particular, we presented a quasi-particle basis of the edge Conformal Field Theory (CFT) for the $`\nu =\frac{1}{p}`$ principal Laughlin states. We demonstrated that the CFT quasi-particles satisfy a form a fractional statistics that is closely related to Haldane’s notion of ‘fractional exclusion statistics’ .
The CFT for the $`\nu =1/p`$ fqH edges can be identified with a continuum ($`N\mathrm{}`$) limit of the Calogero-Sutherland (CS) model for quantum mechanics with inverse square exchange. The natural quasi-particles for the fqH system correspond to eigenstates of the CS hamiltonian. In the work of many groups, the eigenstates of hamiltonians of the CS-type have been understood in terms of Jack symmetric polynomials. In particular, S. Iso presented two alternative Jack polynomial bases for the continuum CS theory. In section 3.4 below, we discuss the precise correspondence between these CS bases and the fqH basis we presented in .
It has been recognized by many groups that quantum many body systems with inverse square exchange come close to being ‘ideal gases of fractional statistics particles’. Supporting this claim are the observations that (i) equilibrium thermodynamic quantities can be evaluated with the help of 1-particle distribution functions for fractional statistics and, (ii) the zero-temperature correlation functions of simple operators (such as electron and hole operators) are mediated by intermediate states with a minimal number of propagating quasi-particles \[see section 4\].
When turning to finite temperature, one quickly discovers that the second ‘free gas’ property (ii) no longer holds. There are important many body effects which can not be ignored when computing correlation functions of physical operators at finite temperature. This implies that great care needs to be taken when setting up arguments that link the $`T`$-dependence of physical observables ($`I`$-$`V`$ and shot-noise characteristics in particular) to the fractional statistics of the fundamental charge carriers .
In this paper we turn to the problem of extending the formalism based on quasi-particles with fractional statistics to finite temperatures. We argue that specific finite-$`T`$ correlation functions can be written in a so-called form factor expansion. Such expansions start with a term that refers to a minimal number of quasi-particles; to this leading term successive corrections are added that refer to more and more quasi-particles participating in physical process described by the correlation function. To test the validity of the proposed expansion, we explicitly evaluate the form factor expansion for a specific finite-$`T`$ Green’s function, collecting all terms that refer to one or two quasi-particles. The numerical results show a rather rapid convergence to the (known) exact expression.
The form factor expansion that we propose is similar in spirit to expressions that were proposed in (see also ) in the context of integrable quantum field theories whose structure is set by a factorizable scattering matrix. Despite this similarity, the two approaches are rather different: in our approach we nowhere rely on scattering data or on the associated form factor axioms, but instead perform explicit computations in a CFT that is regularized by the finite size $`L`$ of the spatial direction.
This paper is organized as follows.
In section 2 we review results obtained in our earlier paper . We present a basis for the edge theories for the $`\nu =\frac{1}{p}`$ fqH state, employing edge electrons and edge quasi-holes as the fundamental charged quasi-particles. We describe the associated 1-particle thermodynamic distribution functions and discuss the fundamental duality between the two types of quasi-particles.
In section 3 we discuss the continuum CS models and the associated Jack polynomial bases as first given by Iso. We give an explicit 1-1 connection between the states in the CS basis and the states in the fqH basis.
In section 4 we turn to form factors. In 4.1 we discuss the ones relevant for $`T=0`$, while in 4.2 we present simple examples of form factors that contribute at non-zero $`T`$. The expressions that we obtain make clear that the simple picture of ‘an ideal gas’ breaks down at finite temperature. A rather general set of selection rules is presented in 4.3, while in section 4.4 we briefly discuss the extent to which our explicit results can be understood from an axiomatic approach based on a factorized S-matrix.
In section 5 we compute 1 and 2-particle diagonal form factors for the ‘edge electron counting operator’. These form factors are then used to evaluate the leading terms in a form factor expansion for the finite-$`T`$ Green’s function that determines the $`I`$-$`V`$ characteristics of edge tunneling processes.
In section 6 we present our conclusions. In the appendices we specify algebraic properties of the charged edge operators for $`\nu =\frac{1}{2}`$, we summarize some relevant results on Jack polynomials, and we provide a proof of the form factor selection rules presented in section 4.3.
## 2 Quasi-particle basis of fqH edge theories
### 2.1 fqH edge theories as CFT
It is well known that many (though not all!) aspects of the low energy dynamics of fqH systems can be captured by an effective edge theory. These edge theories are so-called chiral Luttinger Liquids or chiral CFTs. For the specific case of a principal Laughlin state at filling $`\nu =1/p`$, the CFT describing a single edge (in isolation) is a specific $`c=1`$ CFT.
The standard interpretation of this effective theory is in terms of a chiral boson, which is identified with a quantized density wave (magnetoplasmon) along the edge of the fqH sample. By exploiting the relation between the chiral anomaly and the quantized Hall conductance, one finds that the chiral boson field is compactified on a circle of radius $`R^2=p`$ . This construction directly leads to a space of states (the so-called chiral Hilbert space) with partition function
$$Z^{1/p}(q)=\underset{Q=\mathrm{}}{\overset{\mathrm{}}{}}\frac{q^{\frac{Q^2}{2p}}}{(q)_{\mathrm{}}},$$
(2.1)
with $`(q)_{\mathrm{}}=_{l=1}^{\mathrm{}}(1q^l)`$ and $`q=e^{\beta \frac{2\pi }{L}\frac{1}{\rho _0}}`$. \[The 1-particle energies are spaced by $`l\frac{2\pi }{L}\frac{1}{\rho _0}`$ with $`l`$ integer and $`\rho _0`$ the density of states per unit length, $`\rho _0=(\mathrm{}v_F)^1`$.\]
In eq. (2.1), the parameter $`Q`$ is the electric $`U(1)`$ charge in units of $`\frac{e}{p}`$. The Hilbert space is obtained as a collection of charge sectors $`Q`$. Within each sector, there is a leading state of minimal energy $`\frac{Q^2}{2p}`$; all other states in that sector are reached via the application of (neutral) Fourier modes of the density operator. Together, these modes form a $`U(1)`$ Kac-Moody algebra, and the factor $`\frac{1}{(q)_{\mathrm{}}}`$ is the well-known character of a highest weight module of this affine algebra.
In our earlier paper , we proposed that the CFT for the $`\nu =\frac{1}{p}`$ fqH edge can be interpreted in terms of a set of fundamental charged quasi-particles. We have worked out a formulation in terms of edge electrons (of charge $`e`$) and edge quasi-holes (of charge $`+\frac{e}{p}`$). Our main motivation has been that, using this novel formulation, one learns how to understand some of the unusual and spectacular phenomenology of the fqH systems as manifestations of unusual properties (fractional charge and fractional statistics in particular) of their fundamental quasi-particles.
### 2.2 The fqH quasi-particle basis
In , we demonstrated how the collection of states (2.1) can be understood as a collection of multi-particle states, the fundamental (quasi-)particles being the edge electron and the edge quasi-hole, of charge $`Q=p`$ and $`Q=+1`$, respectively.
The edge electron and quasi-holes are described by the conformal, primary fields
$$J^{(p)}(z)=\underset{t}{}J_t^{(p)}z^{t\frac{p}{2}},\varphi ^+(z)=\underset{s}{}\varphi _s^+z^{s\frac{1}{2p}}.$$
(2.2)
Clearly, one can employ the Fourier modes $`J_t^{(p)}`$ and $`\varphi _s^+`$ as ‘creation operators’ for the corresponding quasi-particles. \[The reason why we put quotation marks here will soon become clear.\] In we identified a collection of multi-$`J`$, multi-$`\varphi `$ states which together form a basis for the chiral Hilbert space. It is given by the states (from now on we omit the charge superindex on the operators $`J`$, $`\varphi `$)
$`J_{(2M1)\frac{p}{2}+Qm_M}\mathrm{}J_{\frac{p}{2}+Qm_1}\varphi _{(2N1)\frac{1}{2p}\frac{Q}{p}n_N}\mathrm{}\varphi _{\frac{1}{2p}\frac{Q}{p}n_1}|Q`$
$`\mathrm{with}m_Mm_{M1}\mathrm{}m_10,n_Nn_{N1}\mathrm{}n_10,`$
$`n_1>0\mathrm{if}Q<0,`$ (2.3)
where $`|Q`$ denotes the lowest energy state of charge $`Q\frac{e}{p}`$, with $`Q`$ taking the values $`(p1)`$, $`(p2)`$, $`\mathrm{}`$, $`1`$, $`0`$.
The associated character identity is
$$Z^{1/p}(q)=\underset{Q=(p1)}{\overset{0}{}}q^{\frac{Q^2}{2p}}Z_Q^{\mathrm{qh}}(q)Z_Q^\mathrm{e}(q),$$
(2.4)
where the factor $`q^{\frac{Q^2}{2p}}`$ takes into account the energy of the initial states and we denoted by $`Z_Q^{\mathrm{qh}}`$, $`Z_Q^\mathrm{e}`$ the partition functions for quasi-holes and electrons in the sector with vacuum charge $`Q`$. They are naturally written as
$$Z_Q^{\mathrm{qh}}=\underset{N=0}{\overset{\mathrm{}}{}}\frac{q^{\frac{1}{2p}(N^2+2QN)+(1\delta _{Q,0})N}}{(q)_N},Z_Q^\mathrm{e}=\underset{M=0}{\overset{\mathrm{}}{}}\frac{q^{\frac{p}{2}M^2QM}}{(q)_M},$$
(2.5)
with $`(q)_L=_{l=1}^L(1q^l)`$. The identity (2.4) can be rigorously established by employing partition counting theorems that are available in the mathematical literature (see for a discussion).
An important special case is $`p=1`$, where the quasi-particle basis (2.3) is the standard multi-particle basis in a theory of free, charged but spin-less fermions.
### 2.3 Fractional statistics and duality
#### 2.3.1 Fractional exclusion statistics
In a 1991 paper , F.D.M. Haldane proposed the notion ‘fractional exclusion statistics’, as a tool for the analysis of strongly correlated many-body systems. The central assumption that is made concerns the way a many-body spectrum is built by filling available one-particle states. In words, it is assumed that the act of filling a one-particle state effectively reduces the dimension of the space of remaining one-particle states by an amount $`g`$. The choices $`g=1`$, $`g=0`$ correspond to fermions and bosons, respectively. The thermodynamics for general ‘$`g`$-ons’, and in particular the appropriate generalization of the Fermi-Dirac distribution function, have been obtained in . The so-called IOW equations
$$\overline{n}_g(ϵ)=\frac{1}{[w(ϵ)+g]},\mathrm{with}[w(ϵ)]^g[1+w(ϵ)]^{1g}=e^{\beta (ϵ\mu )}$$
(2.6)
provide an implicit expression for the 1-particle distribution function $`\overline{n}_g(ϵ)`$ for $`g`$-ons at temperature $`T`$ and chemical potential $`\mu `$. The solutions $`\overline{n}_g(ϵ)`$ have limiting value $`\overline{n}_g^{\mathrm{max}}=\frac{1}{g}`$ for $`ϵ\mathrm{}`$.
#### 2.3.2 Spectral shift statistics and the fqH effect
In we analyzed the exclusion statistics behind the states (2.3) that form a basis for the chiral Hilbert space for a $`\nu =\frac{1}{p}`$ fqH edge. This analysis employed a technique, first proposed in , based on recursion relations satisfied by truncated chiral partition functions. The remarkable conclusions are that
* the ‘microscopic’ state-filling rules differ from those formulated by Haldane, but
* the state counting, and thereby the 1-particle thermodynamic distribution functions, agree with those associated to fractional exclusion statistics, i.e. the distribution functions are solutions of the IOW equations.
The precise statement is that the edge electrons are described by the IOW distribution with $`g=p`$, while the edge quasi-hole states are thermally occupied according to the distribution with $`\stackrel{~}{g}=1/p`$. It is important to remark that there is no mutual statistics between the two types of excitations.
For later reference, we list the explicit expressions for the distribution functions for the case $`p=2`$
$$\overline{n}_2(ϵ)=\frac{1}{2}\left(1\frac{1}{\sqrt{1+4e^{\beta (ϵ\mu )}}}\right),\overline{n}_{{\scriptscriptstyle \frac{1}{2}}}(ϵ)=\frac{2}{\sqrt{1+4e^{2\beta (ϵ\mu )}}}.$$
(2.7)
#### 2.3.3 Duality
The distribution functions for fractional exclusion statistics with parameters $`g`$ and $`\stackrel{~}{g}=1/g`$ satisfy the following duality relation
$$g\overline{n}_g(ϵ)=1\stackrel{~}{g}\overline{n}_{\stackrel{~}{g}}(\stackrel{~}{g}ϵ).$$
(2.8)
The interpretation of this result is that the $`\stackrel{~}{g}`$ quanta with positive energy act as holes in the ground state distribution of negative energy $`g`$-quanta.
Translating back to the $`\nu =\frac{1}{p}`$ fqH edge, we observe a fundamental duality between edge electrons and edge quasi-holes, in agreement with the physical interpretation of strong and weak backscattering limits of edge-to-edge tunneling processes . Under the duality, the removal of a single edge electron corresponds to the creation of a total of $`p`$ quasi-holes.
This duality further implies that, when setting up a quasi-particle description for fqH edges, we can opt for (i) either quasi-holes òr edge electrons, with energies over the full range $`\mathrm{}<ϵ<\mathrm{}`$ (‘particle picture’), or (ii) a combination of both types of quasi-particles, each having positive energies only (‘excitation picture’). The option (ii) is the one realized in the fqH basis of section 2.2. In section 3 we shall discuss the CS bases proposed by Iso, which in a sense uses the option (i).
### 2.4 Equilibrium quantities
#### 2.4.1 Specific heat
The specific heat of a conformal field theory is well-known to be proportional to the central charge $`c_{CFT}`$
$$\frac{C(T)}{L}=\gamma \rho _0k_B^2T,\gamma =\frac{\pi }{6}c_{CFT},$$
(2.9)
where $`\rho _0=(\mathrm{}v_F)^1`$ is the density of states per unit length. In we demonstrated that the fqH quasi-particle basis specified in (2.3) leads to (with $`g=p`$, $`\stackrel{~}{g}=\frac{1}{p}`$)
$$\gamma =\gamma _{g,+}+\gamma _{\stackrel{~}{g},+}$$
(2.10)
with
$$\gamma _{g,+}=_\beta _0^{\mathrm{}}𝑑ϵϵ\overline{n}_g(ϵ),\gamma _{\stackrel{~}{g},+}=_\beta _0^{\mathrm{}}𝑑ϵϵ\overline{n}_{\stackrel{~}{g}}(ϵ).$$
(2.11)
One finds that, while the individual contributions $`\gamma _{g,+}`$, $`\gamma _{\stackrel{~}{g},+}`$ depend on $`g`$, their sum is equal to $`\frac{\pi }{6}`$ for all $`g`$, in agreement with $`c_{\mathrm{CFT}}=1`$. For $`g=2`$, $`\stackrel{~}{g}=\frac{1}{2}`$ one has
$$\gamma _{2,+}=\frac{\pi }{6}\frac{2}{5},\gamma _{{\scriptscriptstyle \frac{1}{2}},+}=\frac{\pi }{6}\frac{3}{5}.$$
(2.12)
#### 2.4.2 Hall conductance
By a simple back-of-the-envelope argument, the Hall conductance is related to the edge capacitance, i.e. to the charge $`\mathrm{\Delta }Q`$ that is accumulated on a given edge in response to an applied voltage $`V`$. One quickly derives the following expression for the Hall conductance
$$G/[\frac{e^2}{h}]=\frac{1}{eV}\left[_0^{\mathrm{}}𝑑ϵ\overline{n}_p(ϵ+eV)+\frac{1}{p}_0^{\mathrm{}}𝑑ϵ\overline{n}_{\frac{1}{p}}(ϵ\frac{e}{p}V)\right].$$
(2.13)
Using the properties of the distribution functions, one shows that this expression is independent of the temperature, and gives $`G=\frac{1}{p}\frac{e^2}{h}`$ for the $`\nu =\frac{1}{p}`$ edge. For $`T=0`$ eq. (2.13) reduces to
$$G=\overline{n}_g^{\mathrm{max}}\frac{q^2}{h},$$
(2.14)
with $`q`$ the charge and $`g`$ the statistics parameter of the quasi-particles that are pulled into the edge by the applied voltage. Depending on the sign of $`V`$ these are the edge electrons ($`q=e`$, $`\overline{n}^{\mathrm{max}}=\frac{1}{p}`$) or the quasi-holes ($`q=\frac{e}{p}`$, $`\overline{n}^{\mathrm{max}}=p`$); both give the canonical value of the Hall conductance.
## 3 CS models and fractional statistics
### 3.1 Inverse square exchange in the CFT setting
While the states specified in (2.3) form a complete basis of the chiral Hilbert space, they are not mutually orthogonal, and as such they do not form a proper starting point for further analysis. In principle one may go through a orthogonalization procedure to arrive at a proper canonical quasi-particle basis. We shall here reach this goal in a more efficient way, by exploiting the close connection with so-called Calogero-Sutherland (CS) models of many-body quantum mechanics.
The CS model describes the (non-relativistic) quantum mechanics of $`N`$ particles on a circle, with 2-body interaction that is proportional to the inverse square of the chord distance between the particles. In , S. Iso demonstrated that the limit $`N\mathrm{}`$ of a CS model with interaction strength $`p(p1)`$ can be identified with the $`c=1`$ CFT of a chiral boson compactified on a circle with radius $`R^2=p`$. Iso also specified a collective hamiltonian $`H_{CS}`$, acting in the CFT Hilbert space, whose eigenstates precisely correspond to the multi-particle states of the underlying CS model.
It turns out that the eigenstates of $`H_{CS}`$ are in 1-1 correspondence with the states specified in eq. (2.3): by adding subleading terms to the expressions in (2.3), one arrives at a (orthogonal and complete) set of eigenstates of $`H_{CS}`$. Comparing with Iso’s formulation, one finds that the ‘superfermions’ of correspond to what we call edge electrons and the ‘anyons’ of are the edge quasi-holes of the fqH system. Nevertheless, the ‘CS basis’ specified by Iso differs from the ‘fqH basis’ of this paper through the way in which the quasi-particle content of a given state is specified. In subsection 3.4 below we shall spell out the precise correspondence between the two formulations.
The spectrum of the CS models has been analyzed with the help of so-called Jack symmetric polynomials, which provide explicit wave functions and eigenstates for the CS hamiltonian. With the help of ‘Jack technology’, important conjectures about zero-temperature correlation functions of models with inverse square exchange have been proven . In the present paper, where we are interested in finite-temperature correlation functions, we shall apply the ‘Jack technology’ to obtain a set of selection rules on form factors that are relevant for computations at finite temperature. We shall complement these considerations with explicit computations of form factors involving states with up to two quasi-particles.
### 3.2 The hamiltonian $`H_{CS}`$ and the fqH basis
To specify the operator $`H_{CS}`$, we employ a free chiral scalar field $`\phi (z)`$. In terms of this scalar field, the charged fields $`J`$ and $`\varphi `$ take the form of so-called vertex operators,
$$J(z)=e^{i\sqrt{p}\phi }(z),\varphi (z)=e^{i\frac{1}{\sqrt{p}}\phi }(z).$$
(3.1)
The operator $`Q=(i\sqrt{p}\phi )`$ measures the electric charge in units $`\frac{e}{p}`$. Following , we define
$$H_{CS}=\frac{p1}{p}\underset{l=0}{\overset{\mathrm{}}{}}(l+1)(i\sqrt{p}\phi )_{l1}(i\sqrt{p}\phi )_{l+1}+\frac{1}{3p}\left[(i\sqrt{p}\phi )^3\right]_0,$$
(3.2)
where $`\phi (z)=_l(\phi )_lz^{l1}`$ and where the second term on the r.h.s. denotes the zero-mode of the normal ordered product of three factors $`(i\sqrt{p}\phi )(z)`$. As a first result, one finds the following action of $`H_{CS}`$ on states containing a single quasi-particle of charge $`Q=+1`$ or $`Q=p`$
$`H_{CS}\varphi _{\frac{1}{2p}n}|0=h_\varphi (n)\varphi _{\frac{1}{2p}n}|0,h_\varphi (n)=\left[{\displaystyle \frac{1}{3p}}+pn(n+{\displaystyle \frac{1}{p}})\right]`$
$`H_{CS}J_{\frac{p}{2}m}|0=h_J(m)J_{\frac{p}{2}m}|0,h_J(m)=\left[{\displaystyle \frac{p^2}{3}}m(p+m)\right].`$
We would like to stress that the fact that both $`J_s`$ and $`\varphi _t`$ diagonalize $`H_{CS}`$ is quite non-trivial. If one evaluates $`H_{CS}`$ on any vertex operator $`\varphi ^{(Q)}`$ (of charge $`Q\frac{e}{p}`$), one typically runs into the field product $`(T\varphi ^{(Q)})(z)`$, where $`T(z)=\frac{1}{2}(\phi )^2(z)`$ is the stress-energy of the scalar field $`\phi `$. Only for $`Q=1`$ and $`Q=p`$ do such terms cancel and do we find that the quasi-particle states are eigenstates of $`H_{CS}`$.
We can now continue and construct eigenstates of $`H_{CS}`$ which contain several $`\varphi `$ or $`J`$-quanta. What one then finds is that the simple product states specified in (2.3) are not $`H_{CS}`$ eigenstates, but that they rather act as head states that need to be supplemented by a tail of subleading terms. For the $`H_{CS}`$ eigenstate headed by the multi-particle state (2.3) (with unit coefficient), we shall use the notation
$$|\{m_j;n_i\}^Q$$
(3.4)
so that
$`H_{CS}|\{m_j;n_i\}^Q=`$
$`\left[{\displaystyle \frac{Q^3}{3p}}+{\displaystyle \underset{j=1}{\overset{M}{}}}h_J((j1)pQ+m_j)+{\displaystyle \underset{i=1}{\overset{N}{}}}h_\varphi ({\displaystyle \frac{1}{p}}(Q+i1)+n_i)\right]|\{m_j;n_i\}^Q.`$
The states (3.4), with the $`m_j`$, $`n_i`$ and $`Q`$ as specified in and below (2.3), form a complete and orthogonal basis for the chiral Hilbert space.
For the sake of illustration, we give explicit results for a few simple states of the fqH basis. The 1-particle states over the $`Q=0`$ vacuum are given by
$$|\{m_1\}^0=J_{\frac{p}{2}m_1}|0,|\{n_1\}^0=\varphi _{\frac{1}{2p}n_1}|0.$$
(3.6)
The norms of these states can explicitly be evaluated by exploiting the (generalized) commutation relations satisfied by the modes $`J_s`$ and $`\varphi _t`$. One finds
$${}_{}{}^{0}\{m_1\}|\{m_1\}_{}^{0}=C_{m_1}^{(p)},{}_{}{}^{0}\{n_1\}|\{n_1\}_{}^{0}=C_{n_1}^{(\frac{1}{p})},$$
(3.7)
where the $`C_k^{(\alpha )}`$ are the expansion coefficients of $`(1x)^\alpha =_{k0}C_k^{(\alpha )}x^k`$.
In our discussion below we shall often restrict ourselves to the vacuum sector $`Q=0`$, and omit the explicit sector label $`Q`$ on the fqH states (3.4).
For later use, we present the explicit form of the fqH states $`|\{m_2,m_1\}`$ and $`|\{m_1;n_1\}`$ at $`p=2`$
$`|\{m_2,m_1\}=|m_2,m_1+{\displaystyle \frac{2}{m_2m_1+3}}{\displaystyle \underset{l>0}{}}|m_2+l,m_1l,`$ (3.8)
$`|\{m_1;n_1\}=|m_1;n_1{\displaystyle \frac{1}{(m_1+2n_1+1)}}{\displaystyle \underset{l>0}{}}|m_1+l;n_1l,`$ (3.9)
with
$$|m_2,m_1=J_{3m_2}J_{1m_1}|0,|m_1;n_1=J_{1m_1}\varphi _{\frac{1}{4}n_1}|0.$$
(3.10)
By explicit evaluation, we obtain the following norms for these states (again at $`p=2`$)
$`N_{\{m_2,m_1\}}=\{m_2,m_1\}|\{m_2,m_1\}={\displaystyle \frac{m_2m_1+1}{m_2m_1+3}}(m_2+3)(m_1+1),`$
$`N_{\{m_1;n_1\}}=\{m_1;n_1\}|\{m_1;n_1\}={\displaystyle \frac{m_1+2n_1+2}{m_1+2n_1+1}}(m_1+1)C_{n_1}^{(\frac{1}{2})}.`$ (3.11)
### 3.3 Jack polynomials and the CS bases
In the previous section, we specified a complete set of eigenstates of the hamiltonian $`H_{CS}`$ in terms of charged quasi-particles $`J`$ and $`\varphi `$. It is clear that these same eigenstates can be obtained by applying an appropriate polynomial in the modes $`a_n=(\phi )_n`$ of the auxiliary scalar field to a vacuum state $`|q`$. It turns out that the polynomials that are needed are so-called Jack polynomials. In appendix B we briefly discuss some of their relevant properties.
Following Iso , we may specify a basis of CS eigenstates as follows
$$|\{\mu \}_J,q=J_{\{\mu ^{}\}}^{\frac{1}{p}}(\{\sqrt{p}a_n\})|q=J_{\{\mu ^{}\}}^{\frac{1}{p}}|q,$$
(3.12)
with the $`U(1)`$ charge $`q`$ running over all integers, and $`\{\mu \}`$ running over all Young tableaus. The norms of these states are given by
$$\{\mu \}_J,q|\{\mu \}_J,q=j_\mu ^{}^{\frac{1}{p}}.$$
(3.13)
An alternative ‘dual’ basis, consists of the states
$$|\{\nu \}_\varphi ,q=J_{\{\nu ^{}\}}^p(\{\frac{a_n}{\sqrt{p}}\})|q=J_{\{\nu ^{}\}}^p|q,$$
(3.14)
with $`q`$ integer and $`\{\nu \}`$ running over all Young tableaus, with norms given by
$$\{\nu \}_\varphi ,q|\{\nu \}_\varphi ,q=j_\nu ^{}^p.$$
(3.15)
### 3.4 Correspondence between fqH and CS bases
Knowing that both the fqH quasi-particle basis and the CS Jack polynomial basis (3.12) are complete bases of orthogonal eigenstates of the operator $`H_{CS}`$, one quickly concludes that there is a 1-1 identification between these two bases. In this section we explicitly describe this 1-1 correspondence.
We start by observing that the Jack state
$$|\left\{\mu \right\}_J,q=J_{\left\{\mu ^{}\right\}}^{\frac{1}{p}}|q$$
(3.16)
can be written as the sum of a leading state
$$J_{m_{\stackrel{~}{M}}+q+\frac{p}{2}}J_{m_{\stackrel{~}{M}1}+q+\frac{3p}{2}}\mathrm{}J_{m_1+q+\frac{(2\stackrel{~}{M}1)p}{2}}|q+\stackrel{~}{M}p,$$
(3.17)
and a ‘tail’ of sub-leading corrections, where ‘subleading’ refers to the triangular form of $`H_{CS}`$ on states of the form (3.17). Here we identified $`m_j=\mu _{\stackrel{~}{M}+1j}1`$, with $`\stackrel{~}{M}=l(\left\{\mu \right\})`$.
Similarly we identify
$$|\left\{\nu \right\}_\varphi ,q=J_{\left\{\nu ^{}\right\}}^p|q$$
(3.18)
with the $`H_{CS}`$ eigenstate headed by
$$\varphi _{n_{\stackrel{~}{N}}\frac{q}{p}+\frac{1}{2p}}\varphi _{n_{\stackrel{~}{N}1}\frac{q}{p}+\frac{3}{2p}}\mathrm{}\varphi _{n_1\frac{q}{p}+\frac{2\stackrel{~}{N}1}{2p}}|q\stackrel{~}{N},$$
(3.19)
with $`n_i=\nu _{\stackrel{~}{N}+1i}1`$ and $`\stackrel{~}{N}=l(\left\{\nu \right\})`$.
Note that, as they stand, the expressions (3.17) and (3.19) are, in general, not of the form (2.3) that defines a member of the fqH basis.
Using the above, we find the following identifications for fqH states which contain only one type of mode operator
$`|\{n_i\}^Q=|\left\{\nu \right\}_\varphi ,Q+N`$
$`|\{m_j\}^Q=|\left\{\mu \right\}_J,QpM,`$ (3.20)
with $`\nu _i=n_{Ni+1}1`$ for $`\stackrel{~}{N}i0`$ and similarly $`\mu _j=m_{Mj+1}1`$ for $`\stackrel{~}{M}j0`$. Note that only the $`n_i`$ that are non-zero become a part of the tableau $`\{\nu \}`$; the $`\varphi `$-modes with $`n_i=0`$ change the charge of the vacuum without exciting any of the $`(\phi )_n`$ modes. \[A similar remark applies to the $`J`$-modes with $`m_j=0`$.\]
The duality eq. (B.3) on the Jack polynomials leads to the following duality relation for the Jack operators,
$$J_{\left\{\lambda ^{}\right\}}^p=(1)^{|\lambda |}j_{\left\{\lambda \right\}}^pJ_{\left\{\lambda \right\}}^{\frac{1}{p}}.$$
(3.21)
This relation enables us rewrite the states $`|\{n_i\}^Q`$ to either the $`|\left\{\nu ^{}\right\}_J,Q+N`$ or the $`|\left\{\nu \right\}_\varphi ,Q+N`$ form. We can for example rewrite
$`|\{n_i\}^Q`$ $`=`$ $`|\left\{\nu \right\}_\varphi ,Q+N`$ (3.22)
$`=`$ $`(1)^{|\nu |}j_{\left\{\nu ^{}\right\}}^p|\left\{\nu ^{}\right\}_J,Q+N.`$
The last state is equivalent to the state which has
$$J_{\nu _1^{}+Q+N+\frac{p}{2}}J_{\nu _2^{}+Q+N+\frac{3p}{2}}\mathrm{}J_{\nu _{n_N}^{}+Q+N+\frac{(2n_N1)p}{2}}|Q+N+n_Np$$
(3.23)
as its leading state (assuming $`n_1>0`$).
One may explicitly check that the eigenvalues of both the Virasoro zero-mode $`L_0`$ and and the CS hamiltonian $`H_{CS}`$ are invariant under the duality transformation that identifies the state headed by (3.19) to the state headed by (3.23).
If there are no $`J`$-operators present in the fqH head state we can use the duality to transform the $`\varphi `$-operators in the head state into $`J`$-operators and we achieve our goal of identitying the fqH state with a member of the CS $`J`$-basis. If the fqH state has both $`J`$ and $`\varphi `$-operators present, we can still map the $`\varphi `$-operators to a dual set of $`J`$-operators. Starting from the state $`|\{m_j\},\{n_i\}^Q`$ we see that, upon using the identity (3.22), the $`J`$-modes associated to the $`m_j`$ ‘see’ their vacuum shifted from $`|Q`$ to $`|Q+N`$. Since the vacuum charge leads to a shift in the values of the mode-indices (see (3.17)), this means that in the CS basis, the corresponding $`J`$-modes will be labeled by $`m_i+N`$ instead of $`m_i`$. It is important to remark that this shift does not affect the contribution of the $`J`$ modes to the eigenvalue of $`H_{CS}`$ on the state. This is because the eigenvalues of $`H_{CS}`$, as specified in eq. (LABEL:Hcseigenvalue), depend directly on the full mode-indices in the head state and not just on the labels $`m_j`$, as can be seen by comparing eq.(LABEL:Hcseigenvalue) with (2.3).
We can now give the full mapping from a fqH basis state to a state in the CS basis
$`|\{m_j;n_i\}^Q=|\left\{\sigma \right\}_J,Q+NpM`$
$`\mathrm{with}\left\{\sigma \right\}=(\left\{m\right\}+N^M)\left\{\nu ^{}\right\},`$ (3.24)
where the sum of the partitions is $`\left\{\mu \right\}+N^M=(\mu _1+N,\mu _2+N,..,\mu _M+N)`$, and the cup product $`\left\{\lambda \right\}\left\{\rho \right\}`$ denotes the partition obtained from sorting the parts $`(\lambda _1,\mathrm{},\lambda _S,\rho _1,\mathrm{},\rho _R)`$ in descending order.
The mapping from the CS basis back to the fqH basis is slightly more complicated. We start from a CS state $`|\left\{\sigma \right\}_J,q`$, with $`S=l(\left\{\sigma \right\})`$, to which we associate a multi-$`J`$ state as in (3.17). We note that the quantity $`\sigma _jqpj`$ decreases with increasing label $`j`$ which allows us to fix $`j`$ such that
$$\sigma _jqpj0,\sigma _{j+1}qp(j+1)<0.$$
(3.25)
The condition on $`\sigma _j`$ guarantees that the associated $`J`$-mode is an allowed mode on a fqH vacuum $`|Q`$ with $`(p1)Q0`$. We now consider the state that remains when all $`J`$-modes left from and at the position $`j`$ are removed. With the use of duality this state can be rewritten,
$$|\left\{\nu ^{}\right\}_J,q+pj|\left\{\nu \right\}_\varphi ,q+pj$$
(3.26)
with
$$\nu _i^{}=\sigma _{i+j}\mathrm{for}i=1\mathrm{}Sj.$$
(3.27)
The $`\varphi `$-operators in the leading state of $`|\left\{\nu \right\}_\varphi ,q+pj`$ act on the vacuum with charge $`\stackrel{~}{Q}=q+pj\sigma _{j+1}`$ vacuum. The second condition in (3.25) guarantees that $`\stackrel{~}{Q}(p1)`$. If one has $`\stackrel{~}{Q}0`$, one identifies $`Q=\stackrel{~}{Q}`$ as the vacuum charge of the fqH basis state. If, however, $`\stackrel{~}{Q}`$ is larger than zero the state $`|\stackrel{~}{Q}`$ was created from $`|0`$ using $`\varphi `$-operators with the highest allowed mode index. From this argument we obtain that the state $`|\left\{\sigma \right\}_J,q`$ from the CS basis can be rewritten as $`|\{m_i;n_i\}^Q`$ with the following rules for selecting $`Q`$, $`M`$ and $`N`$ and the mode indices $`\{m_i\}`$ and $`\{n_i\}`$
$$M=j,N=\mathrm{max}(q+pj,\sigma _{j+1}),Q=q+pMN$$
(3.28)
$`m_i`$ $`=\sigma _{j+1i}N`$ $`\mathrm{for}i=1\mathrm{}j`$
$`n_i`$ $`=\nu _{\sigma _{j+1}+1i}`$ $`\mathrm{for}i=1\mathrm{}\sigma _{j+1}`$ (3.29)
$`=0`$ $`\mathrm{for}i=\sigma _{j+1}+1\mathrm{}N,`$
with $`j`$ and $`\left\{\nu \right\}`$ as specified in (3.25) and (3.27).
### 3.5 Norms for the fqH basis
Of importance for later calculations are the norms of the states $`|\{m_j;n_i\}^Q`$ of the fqH basis. These norms can be evaluated by using the fqH-CS correspondence, together with the norms in the CS basis, as specified in (3.13) and (3.15). The general result is
$$N_{\{m_j;n_i\}}=||\{m_j;n_i^Q|^2=(j_{\left\{\nu ^{}\right\}}^p)^2j_{\left\{\sigma ^{}\right\}}^{\frac{1}{p}},$$
(3.30)
where $`\left\{\nu \right\}`$ and $`\left\{\mu \right\}`$ are the tableaus corresponding to $`\left\{n_i\right\}`$ and $`\left\{m_j\right\}`$, respectively, and the tableau $`\left\{\sigma \right\}`$ is specified in (3.24). The norm can be factorized into the norms of the $`J`$ and $`\varphi `$ parts separately times an extra factor associated with the added partition $`N^M`$
$$N_{\{m_j;n_i\}}=j_{\left\{\nu ^{}\right\}}^pj_{\left\{\mu ^{}\right\}}^{\frac{1}{p}}\underset{(i,j)M^N}{}\frac{(\mu _ji+M)+p(\nu _ij+N+1)}{(\mu _ji+1+M)+p(\nu _ij+N)}.$$
(3.31)
The expressions (3.11) are special cases of this general formula.
In the case where $`m_j,n_i1`$ the expressions simplify and one finds the following factorized form
$$N_{\{m_j;n_i\}}\underset{j=1}{\overset{M}{}}\frac{m_{j}^{}{}_{}{}^{p1}}{\mathrm{\Gamma }(p)}\underset{i=1}{\overset{N}{}}\frac{n_{i}^{}{}_{}{}^{\frac{1}{p}1}}{\mathrm{\Gamma }(\frac{1}{p})}.$$
(3.32)
## 4 Form factors
### 4.1 Vacuum form factors
We start by considering the simplest non-vanishing form factors of the basic electron operators $`J(z)`$, $`J^{}(z)`$ against the multi-particle states in the fqH basis
$`0|J_{\frac{p}{2}+m}|\{n_p,\mathrm{},n_2,n_1\}_N=[N_{\{n_p,\mathrm{},n_2,n_1\}}]^{\frac{1}{2}}f_J(n_p,\mathrm{},n_1)\delta _{m,n_p+\mathrm{}+n_1},`$
$`0|J_{+\frac{p}{2}+m}^{}|\{m_1\}_N=[N_{\{m_1\}}]^{\frac{1}{2}}f_J^{}(m_1)\delta _{m,m_1},`$ (4.1)
where the subscript $`N`$ indicates that the state has been properly normalized.
One immediately finds
$$f_J^{}(m_1)=1.$$
(4.2)
We briefly explain the exact evaluation of the form factor $`f_J(n_p,\mathrm{},n_1)`$ as defined in (4.1). Let us consider the special case $`p=2`$ first. In that case the operator $`J(z)`$ has conformal dimension $`1`$ and may be identified with one of the currents of the affine Kac-Moody algebra $`\widehat{su(2)}_1`$ (see appendix A). By exploiting the OPE
$$\varphi (w_1)\varphi (w_2)=(w_1w_2)^{+{\scriptscriptstyle \frac{1}{2}}}\left[J^{}(w_2)+𝒪(w_1w_2)\right]$$
(4.3)
one obtains
$$J^{}(w_2)=_{C_{w_2}}\frac{dw_1}{2\pi i}(w_1w_2)^{\frac{3}{2}}\varphi (w_1)\varphi (w_2).$$
(4.4)
Using the expansion formula (B.7) we obtain
$$J^{}(w_2)|\mathrm{\hspace{0.17em}0}=\underset{n_2,n_1}{}P_{\{n_2,n_1\}}^{\frac{1}{2}}(w_2,w_2)|\{n_2,n_1\}_\varphi ,q=2$$
(4.5)
and it follows that
$$0|J_{1+m}|\{n_2,n_1\}_N=[N_{\{n_2,n_1\}}]^{\frac{1}{2}}P_{\{n_2,n_1\}}^{\frac{1}{2}}(1,1)\delta _{m,n_2+n_1}.$$
(4.6)
For general $`p`$ one obtains a similar result in terms of Jack polynomials with label $`\frac{1}{p}`$. Using explicit expressions for $`P_{\{\mathrm{}\}}^{\frac{1}{p}}(1,\mathrm{},1)`$, we obtain
$$f_J(n_p,\mathrm{},n_1)=\underset{i=0}{\overset{p1}{}}\frac{\mathrm{\Gamma }(\frac{1}{p})}{\mathrm{\Gamma }(1\frac{i}{p})}\underset{i<j}{}\frac{\mathrm{\Gamma }(n_jn_i+\frac{ji+1}{p})}{\mathrm{\Gamma }(n_jn_i+\frac{ji}{p})}.$$
(4.7)
In the large volume limit, where all $`n_i1`$, this becomes
$$f_J(n_p,\mathrm{},n_1)=\frac{[\mathrm{\Gamma }(\frac{1}{p})]^p}{_{i=0}^{p1}\mathrm{\Gamma }(1\frac{i}{p})}\underset{i<j}{}(n_jn_i)^{\frac{1}{p}}.$$
(4.8)
The form (4.8) of the form factor can be viewed as a limit in (chiral) CFT of a result on correlation functions for the ‘classical’ model of quantum mechanics with inverse square exchange. This result was conjectured by Haldane and later proven in .
An important insight is that there are no other non-vanishing form factors of the state $`J_{\frac{p}{2}m}^{}|0`$ with the elements of the fqH basis. In other words, the spectral weight of this state is completely accounted for by states having precisely the minimal number of $`p`$ quasi-holes. On the basis of this observation, it has been proposed that the $`T=0`$ particle system underlying the fqH basis be viewed as an ‘ideal gas of fractional statistics particles’. In formula, this completeness is expressed by the following identity for the $`T=0`$ 2-point function of the edge electron operator $`J_t`$
$`{\displaystyle \frac{1}{m+1}}0|J_{m+1}J_{m1}^{}|0`$ (4.9)
$`={\displaystyle \frac{1}{m+1}}{\displaystyle \underset{n_2n_1,n_2+n_1=m}{}}0|J_{m+1}|\{n_2,n_1\}_N{}_{N}{}^{}\{n_2,n_1\}|J_{m1}^{}|0`$
$`={\displaystyle \frac{1}{m+1}}{\displaystyle \underset{n_2n_1,n_2+n_1=m}{}}{\displaystyle \frac{(n_2n_1)}{(n_2n_1)^{\frac{1}{2}}}}={\displaystyle _0^{\frac{1}{2}}}𝑑x{\displaystyle \frac{(12x)}{[x(1x)]^{{\scriptscriptstyle \frac{1}{2}}}}}=1,`$
where $`x=\frac{m_1}{m}`$ and we inserted asymptotic expressions valid for $`m,n_2,n_11`$. This result is in agreement with a direct computation using algebraic properties of the $`J`$ modes.
We remark here that, as we shall see in the next sections, the structure of more general form factors, shows ‘many body effects’ and is not easily reconciled with a notion of an ideal gas of fractional statistics particles.
### 4.2 More general form factors
We now consider more general form factors for the edge electron ‘annihilation operator’ $`J^{}(z)`$. The simplest form factor with a 2-particle in-state is
$`{}_{N}{}^{}\{m_1^{}\}|J_{\frac{3p}{2}+m}^{}|\{m_2,m_1\}_{N}^{}=`$ (4.10)
$`\left[{\displaystyle \frac{N_{\{m_2,m_1\}}}{N_{m_1^{}}}}\right]^{\frac{1}{2}}f_{J|JJ}(m_1^{},m_2,m_1)\delta _{m,m_2+m_1m_1^{}}.`$
Using the expansion formula (B.8) we derive the following general result
$$f_{J|JJ}(m_1^{},m_2,m_1)=\underset{s=0}{\overset{p}{}}C_s^{(p)}p_{m_1^{}m_1s}^{(p)}$$
(4.11)
where
$$P_{\{m_2,m_1\}}^p(z_1,z_2)=z_1^{m_2}z_2^{m_1}\underset{l=0}{\overset{m_2m_1}{}}p_l^{(p)}\left(\frac{z_2}{z_1}\right)^l.$$
(4.12)
Specializing to $`p=1,2,3`$, we have the following explicit results
$`p=1:f_{J|JJ}(m_1^{},m_2,m_1)=\delta _{m_1^{}=m_1}\delta _{m_1^{}=m_2+1}`$
$`p=2:f_{J|JJ}(m_1^{},m_2,m_1)=`$
$`\delta _{m_1^{}=m_1}+\delta _{m_1^{}=m_2+2}{\displaystyle \frac{2}{m_2m_1+1}}\mathrm{\Theta }(m_1<m_1^{}<m_2+2)`$
$`p=3:f_{J|JJ}(m_1^{},m_2,m_1)=`$
$`\delta _{m_1^{}=m_1}\delta _{m_1^{}=m_2+3}{\displaystyle \frac{6(m_2+m_12m_1^{})}{(m_2m_1+1)(m_2m_1+2)}}\mathrm{\Theta }(m_1<m_1^{}<m_2+3).`$
We also consider the case where the in-state contains one $`\varphi `$ and one $`J`$ quantum
$${}_{N}{}^{}\{n_1^{}\}|J_{\frac{p}{2}+m}^{}|\{m_1;n_1\}_{N}^{}=\left[\frac{N_{\{m_1;n_1\}}}{N_{n_1^{}}}\right]^{\frac{1}{2}}f_{\varphi |J\varphi }(n_1^{},m_1,n_1)\delta _{m,m_1+n_1n_1^{}}.$$
(4.14)
Using the expansion formula
$$:J(z)\varphi (w):|0=\underset{m_1;n_1}{}r_{m_1,n_1}(z,w)|\{m_1;n_1\}$$
(4.15)
with
$$r_{m,n}(z,w)=z^{m+1}w^n\frac{m+1+pn}{m+p(n+1)}z^mw^{n+1},$$
(4.16)
we derive
$$f_{\varphi |J\varphi }(n_1^{},m_1,n_1)=\left[\delta _{n_1^{}=n_1}\frac{p1}{m_1+p(n_1+1)}\mathrm{\Theta }(n_1^{}>n_1)\right].$$
(4.17)
While the results for $`p=1`$ are a direct consequence of the Wick theorem, the expressions for $`p1`$ show that the ‘ideal gas interpretation’ is no longer applicable for general $`p`$: both form factors (4.10) and (4.14) can be non-vanishing when the energy of the electron annihilation operator $`J^{}`$ does not match the incoming electron energies $`m_2`$ or $`m_1`$, and where the energy difference is transferred to a second ‘spectator particle’. Furthermore, $`f_{J|JJ}`$ and $`f_{\varphi |J\varphi }`$ are not the only non-vanishing form factors of $`J^{}`$ with two incoming particles. For example, there are non-vanishing overlaps between a state created by applying $`J^{}`$ on the a 2-electron state and states containing more quasi-particles than just a single electron. The additional quasi-particles can be visualized as (neutral) density waves or excitons, which are composed of a single electron and $`p`$ quasi-holes. In the next sub-section, we explore the selection rules that determine in a more general setting the possible out-states for which the form factor of $`J^{}`$ with a given in-state is non-vanishing.
### 4.3 A form factor selection rule
In this section we put a bound on the possible out-state that arise upon acting with an electron creation or annihilation operator on a given in-state. We perform this analysis in the CS basis, where the systematics of Jack polynomials come to our help. Using the mapping of section 3.4, the results can be translated to the fqH basis.
We focus on the form factors
$$\left\{\mu \right\}_J,q|J_{m+q+p\frac{p}{2}}|\left\{\nu \right\}_J,q+p,$$
(4.18)
with (unnormalized) Jack states as in (3.12).
The power of Jack Polynomial technology in analyzing this form factor comes from the fact a product of vertex-operators can be written as a sum over products of a ‘coordinates’ Jack polynomials $`P_{\left\{\lambda \right\}}^p\left(\{z_i\}\right)`$ and ‘bosonic modes’ Jack operators $`J_{\left\{\lambda \right\}}^{\frac{1}{p}}`$. Both the coordinate and the bosonic mode Jack polynomials can be manipulated using the results for Jack polynomials in mathematical literature . In Appendix B, we used this to rearrange the part of the vertex operators that survives after applying them to the vacuum. Here we apply the same method which enables us to analyze the action of a single mode operator $`J_{m+q+p\frac{p}{2}}`$ on a state created by a Jack operator. We note that for $`m0`$ this mode operator creates an additional edge electron.
A product of $`N+1`$ edge electron operators acting on a vacuum $`|\stackrel{~}{q}`$ can be expanded in Jack polynomials and Jack operators,
$`J(w){\displaystyle \underset{i=1}{\overset{N}{}}}J(z_i)|\stackrel{~}{q}=`$
$`w^{\stackrel{~}{q}}{\displaystyle \underset{i=1}{\overset{N}{}}}(wz_i)^p{\displaystyle \underset{i<j}{}}(z_iz_j)^p{\displaystyle \underset{i=1}{\overset{N}{}}}z_i^{\stackrel{~}{q}}`$
$`\times \left({\displaystyle \underset{\{\lambda \}}{}}(1)^{|\lambda |}P_{\left\{\lambda \right\}}^p\left(w\right)J_{\left\{\lambda ^{}\right\}}^{\frac{1}{p}}\right)\left({\displaystyle \underset{\{\rho \}}{}}(1)^{|\nu |}P_{\left\{\rho \right\}}^p\left(\{z_i\}\right)J_{\left\{\rho ^{}\right\}}^{\frac{1}{p}}\right)|\stackrel{~}{q}pNp.`$
The state $`|\left\{\nu \right\}_J,\stackrel{~}{q}Np`$ with $`N=l(\left\{\nu \right\})`$ can be extracted from $`_{i=1}^NJ(z_i)|\stackrel{~}{q}`$ by applying the operator
$`O_{\left\{\nu \right\},\stackrel{~}{q}}^p\left(\{z_i\}\right)=`$
$`(p_{\left\{\nu \right\},l(\left\{\nu \right\})}^p)^1\left({\displaystyle \underset{i=1}{\overset{l(\left\{\nu \right\})}{}}}{\displaystyle \frac{dz_i}{2\pi i}\frac{1}{z_i^{\stackrel{~}{q}+1}}}\right)P_{\left\{\nu \right\}}^p\left(\{z_i^1\}\right)\mathrm{\Delta }^p\left(\left\{z_i^1\right\}\right),`$
where we made use of the inner product (B.6) on coordinate dependent Jack polynomials. The norm $`p_{\left\{\nu \right\},l(\left\{n\right\}}^p`$ of the Jack polynomials will drop out of the final result. We can write
$`J_{m+q+p\frac{p}{2}}|\left\{\nu \right\}_J,\stackrel{~}{q}pl(\left\{\nu \right\})=`$ (4.21)
$`M_{mqp}(w)J(w)O_{\left\{\nu \right\},\stackrel{~}{q}}^p\left(\{z_i\}\right){\displaystyle \underset{i=1}{\overset{l(\left\{\nu \right\})}{}}}J(z_i)|\stackrel{~}{q}`$
where
$$M_m(w)=\frac{dw}{2\pi i}\frac{1}{w^{m+1}}.$$
(4.22)
From the definition it is clear that $`O_{\left\{\nu \right\},\stackrel{~}{q}}^p\left(\{z_i\}\right)`$ commutes with $`J(w)`$, which allows us to interchange the order in the expression above and to use the expansion (LABEL:eq:JackExp) of $`J(w)_{i=1}^NJ(z_i)|\stackrel{~}{q}`$ in terms of Jack polynomials. By taking the inner product of the resulting expression with $`\left\{\mu \right\}_J,q|`$, where $`q=\stackrel{~}{q}(N+1)p`$, we obtain an expression for the form factor (4.18).
To decide for which choices of $`\left\{\mu \right\}`$, $`\left\{\nu \right\}`$ and $`m`$ the form factor can be non-vanishing we proceed as follows. We first rewrite the product
$$\underset{i=1}{\overset{N}{}}(wz_i)^p=\underset{\{n_i\}}{}C_{\{n_i\}}^{(p)}w^{pN}\underset{i=1}{\overset{N}{}}\left(\frac{z_i}{w}\right)^{n_i}$$
(4.23)
where $`n_i=0,1,\mathrm{},p`$. We insert this into the expansion given in (LABEL:eq:JackExp), and we use that $`P_{\left\{\lambda \right\}}^p\left(w\right)`$ is zero when $`\left\{\lambda \right\}`$ is not of the form $`\left\{\lambda \right\}=(\lambda _1,0,0,\mathrm{})`$. The contour integration contained in $`M_{mpq}(w)`$ selects the following value for $`\lambda _1`$
$$\lambda _1=|n|+m.$$
(4.24)
Varying $`|n|=_in_i`$ over all allowed values $`|n|=0,\mathrm{},pN`$, we find that $`\lambda _1`$ has to satisfy the inequalities
$$m\lambda _1m+pN.$$
(4.25)
We write $`\left\{\mu \right\}_J|\left\{\lambda \right\}_J\left\{\rho \right\}_J`$ for the inner product $`\left\{\mu \right\}_J,q|J_{\left\{\lambda ^{}\right\}}^{\frac{1}{p}}J_{\left\{\rho ^{}\right\}}^{\frac{1}{p}}|q`$. Using a result by Stanley (Proposition 5.3 from ) we learn that this inner product is non-zero if and only if $`\rho \mu `$ and $`\mu /\rho `$ is a horizontal $`\lambda _1`$-strip. The relation $`\rho \mu `$ indicates that for all $`i`$ we have $`\rho _i\mu _i`$. The skew tableau $`\mu /\rho `$ is the tableau containing all boxes which are in the tableau $`\mu `$ but not in the tableau $`\rho `$. If every column of the skew tableau contains at most one box it is called a horizontal strip and if furthermore the total number of boxes in it is $`\lambda _1`$ it is called a horizontal $`\lambda _1`$-strip.
A particular consequence is that $`\lambda _1`$ satisfies
$$0\lambda _1\mu _1.$$
(4.26)
Combining this inequality with (4.25), we conclude that
$$\mathrm{max}(0,m)\lambda _1\mu _1,$$
(4.27)
which implies that
$$\mu _1<mJ_{mqp+\frac{p}{2}}^{}|\{\mu \}_J,q=0,$$
(4.28)
in agreement with explicit step-functions in the form factors (LABEL:FFjres) and (4.17).
The form factor is now written as
$`q,\left\{\mu \right\}_J|J_{m+q+p\frac{p}{2}}|\left\{\nu \right\}_J,q+p=`$
$`(p_{\left\{\nu \right\},l(\left\{\nu \right\})}^p)^1{\displaystyle \underset{\{n\},\{\lambda _1\},\{\rho \}}{}}\delta _{|n|+m,\lambda _1}C_{\{n\}}^{(p)}\left\{\mu \right\}_J|\left\{\lambda _1\right\}_J\left\{\rho \right\}_J`$
$`\times \left({\displaystyle \underset{i=1}{\overset{N}{}}}M_0(z_i)\right)m_{\{n\}}(\{z_i\})\mathrm{\Delta }^p\left(\left\{z_i^1\right\}\right)\mathrm{\Delta }^p(\{z_i\})P_{\left\{\rho \right\}}^p\left(\{z_i\}\right)P_{\left\{\nu \right\}}^p\left(\left\{z_i^1\right\}\right),`$
where the summations extend over $`i,j=1,\mathrm{},N`$ and $`\{n\}=(n_1,\mathrm{},n_N)`$ with $`n_i0,1,\mathrm{},p`$.
The last part of this expression is the inner product on products of Jack polynomials with a finite number of arguments $`\{z_i\}`$. In Appendix C we discuss restrictions on the tableaus $`\{n\}`$, $`\left\{\nu \right\}`$ and $`\left\{\rho \right\}`$ that follow from imposing that this final inner product be non-zero. Combining all ingredients, one arrives at the following
> Form factor selection rule
>
> The form factor
>
> $$\left\{\mu \right\}_J,q|J_{m+q+p\frac{p}{2}}|\left\{\nu \right\}_J,q+p$$
>
> can only be non-zero if $`\left\{\mu \right\}`$, $`\left\{\nu \right\}`$, $`m`$ satisfy the following conditions
>
> * $`|\nu |+m=|\mu |`$
> * + $`\nu _j\mu _{j+1}`$ for all $`j`$
> + $`\nu _ip`$ for $`i>l(\left\{\mu \right\})`$
> * $`m+_{il(\left\{\mu \right\})}\nu _i\mu _1`$
> * $`_{i=1}^j\nu _i_{i=1}^j(\mu _i+p)`$.
>
> These conditions imply that the tableau $`\left\{\nu \right\}`$ should have at most $`p`$ legs and $`l(\left\{\mu \right\})`$ arms, see fig. 4.1.
We refer to Appendix C for a complete proof of this result.
We remark that the above selection rule can be viewed as a generalization of a selection rule that was used by Lesage, Pasquier and Serban for the evaluation of the zero-temperature density-density correlation function in the Calogero-Sutherland model. These authors found that the (neutral) density operator $`\rho `$ when acting on the vacuum creates an ‘exiton’ with $`p`$ quasi-holes and a single electron, corresponding to a Young tableau with $`p`$ legs and a single arm. In the processes described by the form factor discussed in this section a similar structure is found. Starting from a multi-$`J`$ state described by a tableau $`\left\{\mu \right\}`$, the operator labeled by $`m`$ annihilates one of the $`J`$-quanta. If there is a mismatch between the modes of the operator and of the quantum that is annihilated, the remaining momentum is carried away by a density fluctuation, which roughly speaking corresponds to one extra arm and the $`p`$ legs that can be present in the tableau $`\left\{\nu \right\}`$. If one starts from a state which, in the fqH basis, has a number of $`J`$-quanta and $`N`$ quasi-holes, with maximal mode $`n_N`$, one finds that upon annihilating a $`J`$ mode up to $`n_N+1`$ exitons can be created.
### 4.4 Relation with $`S`$-matrix approach
In this section, we consider the structure of the quasi-particle form factors from the point of view of an associated $`S`$-matrix structure.
Via the TBA procedure, the distribution functions for fractional exclusion statistics are linked to an $`S`$-matrix with the following dependence on particle rapidities $`\theta =\theta _2\theta _1`$
$$S_{ab}(\theta )=\mathrm{exp}[2\pi i(\delta _{ab}G_{ab})\mathrm{\Theta }(\theta )].$$
(4.30)
Although the quasi-particle states that we have considered are part of the discrete spectrum of a finite size system, it is natural to identify the quasi-particle states with a set of asymptotic particle states in a scattering theory with 2-body $`S`$-matrix of this type, with diagonal statistics matrix $`G_{11}=p`$, $`G_{22}=\frac{1}{p}`$.
Via the well-known form factor axioms, this identification leads to specific properties of the form factors. In particular, we expect factors
$$(ϵ_iϵ_j)^p,(\stackrel{~}{ϵ}_i\stackrel{~}{ϵ}_j)^{\frac{1}{p}},$$
(4.31)
in form factors with particles $`J(ϵ_i)`$ and $`\varphi (\stackrel{~}{ϵ}_j)`$ in the in-state, and annihilation poles between particles in the in- and out-states.
The explicit result (4.8) for the vacuum form factor $`f_J`$ has the expected zero’s $`(n_in_j)^{\frac{1}{p}}`$. For the more general form factors discussed in section 4.2 the structure is less clear. We observe however that, upon heuristically replacing
$$\delta _{m_2,m_1}\frac{1}{(ϵ_2ϵ_1)},\mathrm{\Theta }(m_2m_1)\mathrm{log}(ϵ_2ϵ_1)$$
(4.32)
we have (for $`p=1,2,3`$)
$$[_{ϵ_1^{}}]^{p1}\left[f_{J|JJ}(ϵ_1^{},ϵ_2,ϵ_1)\right]\frac{(ϵ_2ϵ_1)^p}{(ϵ_1^{}ϵ_2)^p(ϵ_1^{}ϵ_2)^p}.$$
(4.33)
It will be most interesting to investigate whether the asymptotic limit of the form factors considered and computed in this paper can be obtained by means of an axiomatic approach.
## 5 Form factor expansion at finite temperature
### 5.1 General remarks
In a system of non-interacting electrons, transport properties such as $`I`$-$`V`$ and noise characteristics are obtained by computing the relevant amplitudes for transmission and reflection of single particles, and then performing a statistical average using a one-particle Fermi-Dirac distribution function. An important goal, that we had in mind when setting up the quasi-particle formulation of fqH edges, is to arrive at a similar description of transport processes in these interacting, highly non-Fermi liquid, systems.
As a first attempt in this direction, one may try to simply replace free electron amplitudes by corresponding amplitudes for fqH quasi-particles, and simultaneously replace the Fermi-Dirac distribution by an appropriate distribution function for fractional statistics. While, as we shall argue, this idea is essentially correct, we stress that a correct implementation is subtle and involves the important concept of a so-called form factor expansion.
In this section, we shall focus on the following finite temperature Green’s functions in the CFT for the $`\nu =\frac{1}{p}`$ fqH edge
$$h(ϵ)=\psi _{\nu =\frac{1}{p}}^{}(ϵ)\psi _{\nu =\frac{1}{p}}(ϵ)_T,H(ϵ)=\psi _{\nu =\frac{1}{p}}(ϵ)\psi _{\nu =\frac{1}{p}}^{}(ϵ)_T,$$
(5.1)
where the operators $`\psi _{\nu =\frac{1}{p}}^{}(ϵ)`$ and $`\psi _{\nu =\frac{1}{p}}(ϵ)`$ are the continuum limits of the edge electron operators $`J_s`$ and $`J_s^{}`$ considered in this paper.
In the next subsection, we recall how this Green’s function is used for the computation of the $`I`$-$`V`$ characteristics for the tunneling of electrons into a fqH edge. After that, we give the general form of the form factor expansions for finite temperature correlation functions. We shall then zoom in on the case $`m=2`$, and explain how the finite-$`T`$ Green’s function $`h(ϵ)`$ can be approximated in a form factor expansion.
### 5.2 Kinetic equation for electron tunneling
As explained in , the Green’s functions (5.1) can be used to computed the finite temperature $`I`$-$`V`$ characteristics for the tunneling of electrons from a Fermi Liquid (FL) reservoir into a $`\nu =\frac{1}{p}`$ fqH edge. Starting from the tunneling hamiltonian
$$H_{int}tdϵ[\psi _{\mathrm{FL}}^{}(ϵ)\psi _{\nu =\frac{1}{3}}(ϵ)+\mathrm{h}.\mathrm{c}.],$$
(5.2)
one can show that, in lowest order perturbation theory, the current-voltage characteristics are given by
$$I(V,T)et^2_{\mathrm{}}^{\mathrm{}}𝑑ϵ\left[f(ϵeV)H(ϵ)F(ϵeV)h(ϵ)\right],$$
(5.3)
where $`f(ϵ)`$ and $`F(ϵ)`$ are the standard Fermi-Dirac distributions for particles and holes. Using the conformal mapping from a plane to a cylinder, or employing an imaginary time approach, one finds the following exact expression for the case $`\nu =\frac{1}{3}`$
$$H(ϵ)=\frac{ϵ^2+\frac{\pi ^2}{\beta ^2}}{e^{\beta ϵ}+1},h(ϵ)=\frac{ϵ^2+\frac{\pi ^2}{\beta ^2}}{1+e^{\beta ϵ}}.$$
(5.4)
They lead to $`I`$-$`V`$ characteristics
$$I(V,T)et^2\beta ^3\left(\frac{\beta eV}{2\pi }+\left(\frac{\beta eV}{2\pi }\right)^3\right),$$
(5.5)
in agreement with the result obtained in different approaches . The $`I`$-$`V`$ characteristics (5.5) show cross-over from a linear (thermal) regime into a power-law behavior at high voltages and thus presents a clear fingerprint of the Luttinger liquid features of the fqH edge. The experimental results of are in agreement with these predictions. (See for a further theoretical analysis of the data.)
### 5.3 Form factor expansion
As a proto-type study for a form factor expansion based on CFT quasi-particles, we now analyze the Green’s function $`h(ϵ)`$, for $`p=2`$ in that spirit. Obviously, an exact result is easily obtained
$$h(ϵ)=\frac{ϵ}{e^{\beta ϵ}1}.$$
(5.6)
The Bose-Einstein denominator in this expression has its origin in the fact that the operators $`J`$, $`J^{}`$ satisfy bosonic commutation relations. In the spirit of the quasi-particle formulation of this paper, we wish to treat the $`J`$, $`J^{}`$-quanta as quasi-particles with exclusion statistics $`g=2`$, and see if we can recover the Green’s function $`h(ϵ)`$ in such an approach.
The Green’s function $`h(ϵ)`$ can be viewed as a one-point function for the operator $`N_\psi (ϵ)=\psi _{\nu =\frac{1}{p}}^{}\psi _{\nu =\frac{1}{p}}(ϵ)`$. In the formulation on the finite system of size $`L`$, this operator is represented as $`N_J(m)=aJ_mJ_m^{}`$, with $`ϵ=am`$, with $`a=\frac{2\pi }{L\rho _0}`$ the energy level spacing in the finite size system. This one-point function is formally expressed as
$$\frac{\underset{\mathrm{\Psi }}{}\mathrm{\Psi }|N_J(m)|\mathrm{\Psi }}{_\mathrm{\Psi }\mathrm{\Psi }|\mathrm{\Psi }}.$$
(5.7)
The sum runs over a basis the full Hilbert space of the edge CFT, and we can opt for the fqH quasi-particle basis discussed in this paper. The idea is now that the matrix elements $`\mathrm{\Psi }|N_J(m)|\mathrm{\Psi }`$ are dominated by processes where only a few of the quasi-particles that are present in a concrete basis state $`|\{m_i;n_j\}`$ participate (we restrict our attention to states in the $`Q=0`$ sector of the fqH basis).
For the case at hand, the lowest contributions comes from 1-particle states $`|\{m_1\}`$, for which one computes the form factor
$`D^{(1,0)}(m;m_1)={}_{N}{}^{}\{m_1\}|J_{1m}J_{+1+m}^{}|\{m_1\}_{N}^{}=`$ (5.8)
$`(m+1)\delta _{m,m_1}+2\left(1{\displaystyle \frac{m+1}{m_1+1}}\right)\mathrm{\Theta }(m<m_1).`$
The expected presence of an edge electron of energy $`m_1`$ is given by the distribution function $`\overline{n}_2(ϵ_1=am_1)`$. This leads to the following contribution to the Green’s function
$$h^{(1,0)}(ϵ)=a\underset{m_1}{}D^{(1,0)}(m,m_1)\overline{n}_2(am_1).$$
(5.9)
If we now consider the form factor of $`N_J(m)`$ against a two-electron state, we find (see next subsection) that it is not simply the sum of two 1-particle contributions. The left-over part is what we call the irreducible 2-electron form factor
$`D^{(2,0)}(m;m_1,m_2)={}_{N}{}^{}\{m_1,m_2\}|J_{3m}J_{+3+m}^{}|\{m_2,m_1\}_{N}^{}`$ (5.10)
$`{}_{N}{}^{}\{m_1\}|J_{3m}J_{+3+m}^{}|\{m_1\}_{N}^{}{}_{N}{}^{}\{m_2\}|J_{3m}J_{+3+m}^{}|\{m_2\}_{N}^{}.`$
It leads to an additional contribution $`h^{(2,0)}(m)`$ to the Green’s function
$$h^{(2,0)}(ϵ)=a\underset{m_1,m_2}{}D^{(2,0)}(m;m_1,m_2)\overline{n}_2(am_1)\overline{n}_2(am_2).$$
(5.11)
Similarly, we define
$`D^{(1,1)}(m;m_1,n_1)=`$
$`{}_{N}{}^{}\{n_1,m_1\}|J_{1m}J_{+1+m}^{}|\{m_1,n_1\}_{N}^{}{}_{N}{}^{}\{m_1\}|J_{1m}J_{+1+m}^{}|\{m_1\}_{N}^{}`$
and
$$h^{(1,1)}(ϵ)=a\underset{m_1,n_1}{}D^{(1,1)}(m;m_1,n_1)\overline{n}_2(am_1)\overline{n}_{\frac{1}{2}}(an_1).$$
(5.13)
Continuing in this manner, we build up the following expansion
$`h={\displaystyle \underset{M,N}{}}h^{(M,N)}(ϵ),`$
$`h^{(M,N)}(ϵ)=a{\displaystyle \underset{\{m_i;n_j\}}{}}D^{(M,N)}(m;\{m_i;n_j\}){\displaystyle \underset{i}{}}\overline{n}_2(am_i){\displaystyle \underset{j}{}}\overline{n}_{\frac{1}{2}}(an_j).`$
(5.14)
We remark that an expansion of precisely this type has been proposed by LeClair and Mussardo , see also . This work was done in the context of integrable qft’s, that are fully characterized by a factorized $`S`$-matrix. In such a context, the irreducible form factors are constrained by the form factor axioms, and the distribution functions have their origin in a TBA procedure. Although clearly in the same spirit, the analysis that we present here is very different at the technical level. We obtain the relevant form factor by explicit computation in a theory that is regularized by the finite size of the fqH edge, and we have identified the relevant distribution functions by analyzing the state counting of the (discrete) spectrum of the finite-size system. We thus do not rely on an underlying (massless) $`S`$-matrix point of view. Despite these differences, it seems clear that the two approaches are closely related: in subsection 4.5 we briefly indicated that our form factor have symmetry properties that are expected on the basis of a ‘purely statistical $`S`$-matrix’. We leave this interesting issue for further study.
### 5.4 Irreducible form factors
To evaluate explicitly the leading terms in the form factor expansion (5.14) for $`h(ϵ)`$, we need to evaluate the relevant irreducible form factors. While it is clear that these form factors have very special mathematical properties, we here compute them by a simple brute force computation, relying on the explicit form of the two-particle states (3.8) and (3.9), and on the algebraic properties of the operators $`J`$, $`J^{}`$ and $`\varphi `$ (see appendix A).
#### 5.4.1 Two electrons
For the irreducible two electron form factor we find
$`D^{(2,0)}(m;m_2,m_1)=`$ (5.15)
$`\delta _{mm_2}{\displaystyle \frac{2(m_2+3)}{m_2m_1+3}}+\delta _{mm_1+2}{\displaystyle \frac{2(m_1+1)}{m_2m_1+1}}`$
$`+{\displaystyle \frac{4}{(m_2m_1+3)}}{\displaystyle \frac{1}{(m_2m_1+1)}}{\displaystyle \frac{1}{(m_1+1)(m_2+3)}}`$
$`\times [\mathrm{\Theta }(m<m_12)P(m;m_1,m_2)`$
$`+\mathrm{\Theta }(m<m_2<m+m_1)Q(m;m_1,m_2)`$
$`+\mathrm{\Theta }(m<m_2)R(m;m_1,m_2)],`$
with
$`P(m;m_1,m_2)=`$
$`(m_2m_1+3)(m_1m2)(2m_1m_23)`$
$`+(m_1m2)(m_1m3)(3m_2+{\displaystyle \frac{5}{3}}m_1+{\displaystyle \frac{1}{3}}m{\displaystyle \frac{26}{3}})`$
$`+(m+3)[2(m_2m_1+3)(2m_1m1)`$
$`2m_1(m_1+1)+(m+3)(m_2+m_1m+1)]`$
$`Q(m;m_1,m_2)=`$
$`(m_1m_2+m+1)[(m_2m_1+3)^2+2(m_2m_1+3)(m_1m_2+m)`$
$`+{\displaystyle \frac{2}{3}}(m_1m_2+m)(m_1m_2+m1)]`$
$`R(m;m_1,m_2)=`$
$`(m_2m)(m_1+1)(m_2m_1+3)+{\displaystyle \frac{1}{3}}m_1(m_1+1)(m_1+3m_23m+2).`$
The polynomials $`P`$, $`Q`$ and $`R`$ enjoy special properties, which include
$$(P+Q+R)(m;m_1,m_2)=\frac{1}{3}(m_1m_21)(m_1m_22)(m_1m_23).$$
(5.17)
#### 5.4.2 One electron and one quasi-hole
The irreducible form factor with one electron and one hole is found to be
$`D^{(1,1)}(m;m_1,n_1)=`$ (5.18)
$`\delta _{m_1,m}{\displaystyle \frac{m_1+1}{m_1+2n_1+1}}`$
$`+\mathrm{\Theta }(m<m_1){\displaystyle \frac{1}{C_{n_1}^{(\frac{1}{2})}(m_1+2n_1+2)(m_1+2n_1+1)(m_1+1)}}`$
$`\times \left[C_{n_1m_1+m}^{(\frac{1}{2})}S(m;m_1,n_1)+C_{n_1}^{(\frac{1}{2})}T(m;m_1,n_1)\right],`$
with
$`S(m;m_1,n_1)=`$
$`(m_1+2n_1+1)^2+(m+n_1m_1)({\displaystyle \frac{8}{3}}4(m_1+2n_1+2))+{\displaystyle \frac{4}{3}}(m+n_1m_1)^2`$
$`T(m;m_1,n_1)=`$ (5.19)
$`2(m_1m)((m_1+2n_1+1)^21)+2(2n_1+1)(m_1m1)`$
$`+2({\displaystyle \frac{2}{3}}n_1+1)(2n_1+1).`$
### 5.5 Evaluating the series
With the information collected in the previous subsections, we can evaluate the 1-particle and 2-particle contributions $`h^{(1,0)}`$, $`h^{(2,0)}`$ and $`h^{(1,1)}`$ to the Green’s function $`h(ϵ)`$.
The expressions (5.9), (5.11) and (5.13) for $`h^{(2,0)}`$ and $`h^{(1,1)}`$ are discrete sums, which we wish to study in the limit $`a0`$. In this limit, one may view the expressions as Riemann sums and evaluate them using continuous integrals; however, one needs to be careful because the integrands as they stand are have singularities, and the sums are not term-by-term convergent. One may check however that by carefully redistributing some of the terms, one obtains convergent sums that can be approximated by the corresponding continuous integrals. Proceeding in this manner, and using a numerical integrator, we obtained the results plotted in figure 5.1.
We observe that the form factor series converge in the following sense: while the 1-particle terms agree with the exact result for $`ϵ`$ greater than about $`3k_BT`$, the approximation including 2-particle terms reaches the exact curve at $`ϵ`$ around $`2k_BT`$. For energies $`ϵk_BT`$, the thermal factors do not efficiently suppress many particle contributions, and the convergence of the form factor expansion is expected to be slow.
We remark that the asymptotic behavior for $`ϵk_BT`$ of the 2-particle terms is
$$h^{(2,0)}(ϵ)c_2e^{\beta ϵ},h^{(1,1)}(ϵ)c_{\frac{1}{2}}e^{\beta ϵ}$$
(5.20)
with
$$c_2=2_0^{\mathrm{}}𝑑ϵ_1\overline{n}_2(ϵ_1),c_{\frac{1}{2}}=_0^{\mathrm{}}𝑑\stackrel{~}{ϵ}_1\overline{n}_{\frac{1}{2}}(\stackrel{~}{ϵ}_1).$$
(5.21)
Remarkably, the duality relation (2.8) leads to the relation
$$c_2=c_{\frac{1}{2}}$$
(5.22)
meaning that the Boltzmann tails of the 2-particle terms precisely cancel. This ‘conspiracy’ was needed as, numerically, it is seen that the deviation between the exact curve $`h(ϵ)`$ and the 1-particle term $`h^{(1,0)}(ϵ)`$ is far smaller than the individual Boltzmann tails of $`h^{(2,0)}`$ and $`h^{(1,1)}`$.
## 6 Conclusions
Summarizing the results collected in this paper, we have made some first steps on the way to realizing a computational scheme where the $`T`$-dependence of physical observables in a fqH system (charge transport properties in particular) is computed with direct reference to fractional statistics of the fundamental quasi-particles. We expect that on the basis of the formalism presented here, meaningful claims about the observability of the fractional statistics of CFT edge quasi-particles can be formulated. We leave this most interesting aspect for further study.
We remark the the continuum (CFT) limit of the CS model provides an ideal testing ground for form factor expansions for finite temperature correlation functions, such as discussed in section 5.3 and in the literature . This is because on the one hand the theory is explicitly regularized by the finite extent of the spatial direction and, on the other, the finite temperature Green’s functions are known from standard CFT considerations.
We thank A.W.W. Ludwig for many insightful comments and collaboration in the early stages of this project, and F.H.L. Essler and F.A. Smirnov for discussions. This research has benefited from the NATO Collaborative Research Grant SA.5-2-05(CRG.951303) and from support from the foundation FOM of the Netherlands. We thank the Erwin Schrödinger Institute (Vienna) and the Centre de Recherches Mathématiques (Montreal) for hospitality during the course of this work.
## Appendix A Algebraic properties of $`\nu =\frac{1}{2}`$ edge operators
The charged edge operators $`J=J^{}`$, $`J^{}=J^+`$ in the edge theory at $`\nu =\frac{1}{2}`$ are part of a $`SU(2)_1`$ affine symmetry algebra. Together with the charge density $`Q=i\sqrt{p}\phi `$ they satisfy the commutation relations
$`[J_{m_2}^+,J_{m_1}^{}]=m_1\delta _{m_2+m_1}+Q_{m_2+m_1}`$
$`[Q_{m_2},J_{m_1}^\pm ]=\pm 2J_{m_2+m_1}^\pm ,[Q_{m_2},Q_{m_1}]=2m_1\delta _{m_2+m_1}.`$ (A.1)
The fractionally charged edge quasi-particles $`\varphi ^\pm `$ transform in the spin-$`\frac{1}{2}`$ representation of the $`SU(2)`$ symmetry
$$[J_m^\pm ,\varphi _s^\pm ]=0,[J_m^\pm ,\varphi _s^{}]=\pm \varphi _{m+s}^\pm ,[Q_m,\varphi _s^\pm ]=\pm \varphi _{m+s}^\pm .$$
(A.2)
Among themselves, the modes of $`\varphi ^\pm `$ satisfy so-called generalized commutation relations, which have been studied in the context of the spinon formulation of the $`SU(2)_1`$ CFT .
## Appendix B Jack polynomials and Jack operators
In this appendix we briefly introduce the Jack polynomials that are used in sections 3 and 4 of this paper. We essentially follow the conventions of Iso , but we introduce different notations for the coordinate dependent Jack polynomials $`P_{\left\{\mu \right\}}^\beta \left(\{z_i\}\right)`$ and the bosonic mode Jack operators $`J_{\left\{\mu \right\}}^\beta \left(\left\{\frac{a_n}{\sqrt{\beta }}\right\}\right)`$ .
We start by specifying an inner product on the ring of symmetric polynomials,
$$p_{\left\{\lambda \right\}}|p_{\left\{\mu \right\}}_\beta =\delta _{\left\{\lambda \right\},\left\{\mu \right\}}\beta ^{l(\left\{\lambda \right\})}z_\lambda $$
(B.1)
where $`p_{\left\{\lambda \right\}}(z_i)=_{j=1}^{l(\left\{\lambda \right\})}p_{\lambda _j}(\{x_i\})\mathrm{with}p_{\lambda _j}(\{x_i\})=_ix_i^{\lambda _j}`$ is the power sum set by a Young tableau $`\left\{\lambda \right\}=(\lambda _1,\lambda _2,\mathrm{},\lambda _l)`$, $`z_{\left\{\lambda \right\}}`$ is $`_{i1}i^{l_i}l_i!`$ with $`l_j`$ the number of entries in $`\left\{\lambda \right\}`$ which satisfy $`\lambda _i=j`$ and $`\beta `$ a rational number.
The coordinate Jack polynomials $`P_{\left\{\lambda \right\}}^\beta \left(\{z_i\}\right)`$ are symmetric functions in the coordinates $`\{z_i\}`$ labeled by a Young tableau $`\{\lambda \}`$ and a rational number $`\beta `$. They are defined by the following properties
orthogonality:
$`P_{\left\{\lambda \right\}}^\beta \left(\{z_i\}\right)|P_{\left\{\nu \right\}}^\beta \left(\{z_i\}\right)_\beta =\delta _{\left\{\lambda \right\},\left\{\nu \right\}}j_{\left\{\nu \right\}}^\beta `$
triangularity:
$`P_{\left\{\lambda \right\}}^\beta \left(\{z_i\}\right)=_{\left\{\mu \right\}}v_{\lambda ,\mu }(\beta )m_{\left\{\mu \right\}}`$ where $`v_{\lambda ,\mu }(\beta )=0`$ unless $`\left\{\mu \right\}\left\{\lambda \right\}`$
normalization:
the coefficient $`v_{\lambda ,\lambda }=1`$ .
In this definition, $`m_{\left\{\lambda \right\}}(\{z_i\})`$ are the monomial symmetric functions $`_\sigma _iz_i^{\lambda _{\sigma (i)}}`$ where $`_\sigma `$ denotes the sum over all permutations of the indices $`i`$. The partial ordering $``$ on partitions is the so-called dominance ordering on partitions of equal weight ($`|\lambda |=|\mu |`$): $`\left\{\lambda \right\}\left\{\mu \right\}_i=1^j\lambda _i_{i=1}^j\mu _j`$ for all $`j`$. The function $`j_{\left\{\nu \right\}}^\beta `$ in the inner product can be shown to be given by
$`j_{\left\{\nu \right\}}^\beta `$ $`=`$ $`{\displaystyle \underset{(i,j)\{\nu \}}{}}{\displaystyle \frac{\beta (\nu _j^{}i)+\nu _ij+1}{\beta (\nu _j^{}i+1)+\nu _ij}}.`$ (B.2)
In a notation where Jack polynomials are written as functions of power sums $`p_n`$, they satisfy a duality between $`\beta =p`$ and $`\beta =\frac{1}{p}`$
$$P_{\left\{\lambda ^{}\right\}}^p\left(\left\{\frac{p_n}{p}\right\}\right)=(1)^{|\lambda |}j_{\left\{\lambda \right\}}^pP_{\left\{\lambda \right\}}^{\frac{1}{p}}\left(\{p_n\}\right),$$
(B.3)
where $`\left\{\lambda ^{}\right\}`$ is the Young tableau dual to $`\left\{\lambda \right\}`$. It follows that
$$j_{\left\{\nu \right\}}^pj_{\left\{\nu ^{}\right\}}^{\frac{1}{p}}=1.$$
(B.4)
The following elementary property of the Jack polynomials
$$\underset{i,j}{}(1x_iy_j)=\underset{\left\{\lambda \right\}}{}(1)^{|\lambda |}P_{\left\{\lambda \right\}}^\beta \left(\{x_i\}\right)P_{\left\{\lambda ^{}\right\}}^{\frac{1}{\beta }}\left(\{y_j\}\right)$$
(B.5)
can be used to rewrite expressions involving vertex-operators.
For integer $`\beta `$ an alternative inner product on the Jack polynomials $`P_{\left\{\mu \right\}}^\beta `$ depending on only a finite set of coordinates $`\{z_i\}=\{z_1,\mathrm{},z_n\}`$ is given by
$`P_{\left\{\nu \right\}}^\beta \left(\{z_i\}\right)|P_{\left\{\mu \right\}}^\beta \left(\{z_i\}\right)=`$ (B.6)
$`\left({\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{dz_i}{2\pi i}\frac{1}{z_i}}\right)\mathrm{\Delta }^\beta (\{z_i^1\})\mathrm{\Delta }^\beta (\{z_i\})P_{\left\{\nu \right\}}^\beta \left(\{z_i^1\}\right)P_{\left\{\mu \right\}}^\beta \left(\{z_i\}\right),`$
where $`\mathrm{\Delta }^\beta (\{x_i\})=_{i<j}(x_ix_j)^\beta `$ denotes a generalized Vandermonde determinant. Although it is also possible to define this inner product for fractional $`\beta `$ , we will use it in this form for integer $`\beta `$. The Jack polynomials are orthogonal w.r.t. this alternative inner product.
For the product of $`N`$ quasi-hole vertex operators $`\varphi (z_i)`$, the following expression can be derived
$`\varphi (z_1)\mathrm{}\varphi (z_N)|q=`$ (B.7)
$`\mathrm{\Delta }^{\frac{1}{p}}(\{z_i\}){\displaystyle \underset{\left\{\lambda \right\}}{}}(1)^{|\lambda |}P_{\left\{\lambda ^{}\right\}}^{\frac{1}{p}}\left(\{z_j\}\right)J_{\left\{\lambda \right\}}^p\left(\left\{{\displaystyle \frac{a_n}{\sqrt{p}}}\right\}\right){\displaystyle \underset{j=1}{\overset{N}{}}}z_j^{\frac{q}{p}}|q+N,`$
where we wrote $`J_{\left\{\lambda \right\}}^p\left(\left\{\frac{a_n}{\sqrt{p}}\right\}\right)`$ for a Jack polynomials in which power sums $`p_n`$ are replaced by bosonic modes, writing $`a_n=(\phi )_n`$. \[We refer to such expressions as Jack operators.\] Similarly,
$`J(z_1)\mathrm{}J(z_N)|q=`$ (B.8)
$`\mathrm{\Delta }^p(\{z_i\}){\displaystyle \underset{\left\{\lambda \right\}}{}}(1)^{|\lambda |}P_{\left\{\lambda ^{}\right\}}^p\left(\{z_j\}\right)J_{\left\{\lambda \right\}}^{\frac{1}{p}}\left(\{\sqrt{p}a_n\}\right){\displaystyle \underset{j=1}{\overset{N}{}}}z_j^q|qNp.`$
For brevity, we sometimes drop the explicit reference to the bosonic modes and write
$$J_{\left\{\lambda \right\}}^{\frac{1}{p}}J_{\left\{\lambda \right\}}^{\frac{1}{p}}\left(\{\sqrt{p}a_n\}\right),J_{\left\{\lambda \right\}}^pJ_{\left\{\lambda \right\}}^p\left(\{\frac{1}{\sqrt{p}}a_n\}\right).$$
## Appendix C Proof of selection rules
We present a proof of the form factor selection rules of section 4.3.
(a.) This is a consequence of energy conservation and can be found from the product of three delta functions,
$$\delta _{|n|+m,\lambda _1}\delta _{\lambda _1+|\rho |,|\mu |}\delta _{|n|+|\rho |,|\nu |}$$
(C.1)
present implicitly in eq. (LABEL:eq:theorem).
(b.) Proposition 2.4 in states that it is possible to rewrite any Jack polynomial as a linear combination of products of Jack polynomials labeled with horizontal strips $`𝒥_{\left\{\lambda \right\}}^p=_iP_{(\lambda _i)}^p`$,
$$P_{\left\{\lambda \right\}}^p=\underset{\left\{\sigma \right\}\left\{\lambda \right\}}{}\stackrel{~}{q}_{\left\{\lambda \right\}\left\{\sigma \right\}}𝒥_{\left\{\sigma \right\}}^p.$$
(C.2)
An important difference between this expansion and the expansion of a Jack polynomial in monomial symmetric functions is that the sum runs over tableaus $`\left\{\sigma \right\}`$ satisfying $`\left\{\sigma \right\}\left\{\lambda \right\}`$ instead of $`\left\{\sigma \right\}\left\{\lambda \right\}`$. From repeated application of proposition 5.3 in (see also section 4.3) to a product of two Jack polynomials, where one is expanded using the expansion above, it follows that the tableau labeling the non-expanded Jack polynomial is contained in every tableau labeling a Jack polynomial appearing in the product. Exchanging the roles we see that also the tableau labeling the other Jack polynomial is contained in these tableaus.
Combining this knowledge with the fact that by triangularity the monomial symmetric functions $`m_{\left\{n\right\}}`$ can be expanded in Jack polynomials,
$$m_{\left\{n\right\}}=\underset{\left\{\tau \right\}\left\{n\right\}}{}\stackrel{~}{v}_{\left\{n\right\}\left\{\tau \right\}}P_{\left\{\tau \right\}}^p,$$
(C.3)
and applying this to the coordinate inner product in eq. (LABEL:eq:theorem) we find that $`\left\{\rho \right\}`$ is contained in $`\left\{\nu \right\}`$. The operator inner product shows that $`\left\{\mu \right\}`$ differs at most a horizontal strip from $`\left\{\rho \right\}`$ and thus $`\left\{\rho \right\}`$ contains $`\left\{\stackrel{~}{\mu }\right\}=(\mu _2,\mathrm{},\mu _M)`$. We can conclude that $`\left\{\nu \right\}`$ contains $`\left\{\stackrel{~}{\mu }\right\}`$ and we have obtained (b.1.).
We can extract extra information from examining this construction once more, under the addition of a horizontal strip the length of a column can grow with one box. If we now multiply the Jack polynomials $`P_{\left\{tau\right\}}^p`$ appearing in the expansion of the monomial symmetric function $`m_{\left\{n\right\}}`$ with the $`𝒥_{\left\{\sigma \right\}}^p`$ appearing in the expansion of the Jack polynomial $`P_{\left\{\rho \right\}}^p`$ we find that the maximal difference in column length between $`\left\{\nu \right\}`$ and $`\left\{\tau \right\}`$ is $`l(\left\{\rho \right\})`$ from which we can conclude that the only columns in $`\left\{\nu \right\}`$ which have a length exceeding $`l(\left\{\mu \right\})l(\left\{\rho \right\})`$ are those columns for which the corresponding tableau labeling the monomial symmetric function has a column of non-zero length. Because only monomial symmetric functions with at most $`p`$ non-zero columns appear we obtain (b.2.).
(c.) This is a simple consequence of (a.) and (b.1.)
(d.) By definition the Jack polynomials can be expanded in monomial symmetric functions, so we have
$$P_{\left\{\rho \right\}}^p=\underset{\left\{\sigma \right\}\left\{\rho \right\}}{}v_{\left\{\rho \right\}\left\{\sigma \right\}}m_{\left\{\sigma \right\}}.$$
(C.4)
If we now use this expansion $`P_{\left\{\rho \right\}}^p`$ and then multiply the resulting $`m_{\left\{\sigma \right\}}`$ in the expansion by $`m_{\left\{n\right\}}`$, then the products in the expansion will be linear combinations of $`m_{\left\{\sigma ^{}\right\}}`$ satisfying $`\left\{\sigma ^{}\right\}\left\{\sigma \right\}+\left\{n\right\}`$,
$`m_{\left\{n\right\}}P_{\left\{\rho \right\}}^p`$ $`=`$ $`{\displaystyle \underset{\left\{\sigma \right\}\left\{\rho \right\}}{}}v_{\left\{\rho \right\}\left\{\sigma \right\}}m_{\left\{n\right\}}m_{\left\{\sigma \right\}}`$ (C.5)
$`=`$ $`{\displaystyle \underset{\left\{\sigma ^{}\right\}\left\{\rho \right\}+\left\{n\right\}}{}}u_{\left\{n\right\}\left\{\rho \right\}\left\{\sigma ^{}\right\}}m_{\left\{\sigma ^{}\right\}}.`$
Expanding the $`m_{\left\{\sigma ^{}\right\}}`$ as we did in the proof of (b.) the sum over products of monomial symmetric functions can be rewritten in terms of Jack polynomials again,
$$m_{\left\{n\right\}}P_{\left\{\rho \right\}}^p=\underset{\left\{\sigma \right\}\left\{\rho \right\}+\left\{n\right\}}{}w_{\left\{n\right\}\left\{\rho \right\}\left\{\sigma \right\}}P_{\left\{\sigma \right\}}^p,$$
(C.6)
and we find that $`\left\{\nu \right\}`$ is smaller in the sense of dominance order than a tableau $`\left\{\sigma \right\}`$ of the form $`\left\{\rho \right\}+\left\{n\right\}`$. Since $`n_ip`$ and $`\rho _i\mu _i`$, the result (d.) follows.
Figure 5.1: One-particle Green’s function $`h(ϵ)`$ for $`p=2`$ as a function of energy. The solid curve is the exact result; the data points are the numerical results for: $`h^{(1,0)}`$ (diamonds), $`h^{(2,0)}`$ (circles) and $`h^{(1,1)}`$ (plusses). The sum of all contributions with up to two particles is represented by squares.
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# Commensurate-incommensurate transition in two-coupled chains of nearly half-filled electrons
## 1 Introduction
Organic conductors Bechgaard salts, which consist of an array of one-dimensional chains, show various ordered states such as spin-Peierls state, spin density wave state and superconducting state at low temperatures jerome\_revue\_1d . Besides these states of broken symmetry, the normal state above the transition temperature exhibits a remarkable electronic state associated with a charge gap. This charge gap is supposed to originate from the electronic interactions and the commensurate band-filling ($`1/4`$ filling) of the compound Giamarchi\_physica . Indeed in one dimension, commensurate systems are Mott insulators. Such an interpretation of the charge gap has received support from recent optical experiments dressel\_optical\_tmtsf ; Schwartz and transverse conductivity measurements Jerome\_EPJ . In addition, since these compounds are quasi-one-dimensional systems there is a competition between the one dimensional charge gap, that tends to localize the electrons, and the interchain hopping tending to make the system three (two) dimensional. This competition could lead to a confinement-deconfinement transition responsible for the difference of behavior between the TMTTF and the TMTSF salts Giamarchi\_physica ; Gruner\_science . Experimentally the confinement (deconfinement) is found in TMTTF (TMTSF) salts whose interchain hopping is smaller (larger) than the charge gap Gruner\_science . From a theoretical point of view, the transition should take place when the interchain hopping renormalized by interactions is comparable with the single chain gap Bourbonnais ; Giamarchi\_physica ; Schwartz . There is thus still considerable debate on how such a transition takes place, and what are the relevant energy scales.
Unfortunately studying an infinite number of coupled chains is extremely difficult. It is thus worth to investigate these issues on a finite number of coupled chains, i.e. on ladder systems dagotto\_ladder\_review . By studying explicitly two-coupled chains with a half-filled band, it has been shown that the interchain hopping becomes irrelevant, in the sense of renormalization group, for a charge gap larger than the interchain hopping Suzumura ; Tsuchiizu . It is reasonable to identify this relevance (irrelevance) of interchain hopping with the deconfinement (confinement) for an infinite number of chains since the interchain hopping is renormalized to zero in the limit of low energy for the irrelevant case. Of course there are some differences between the simplified ladder system and the infinite number of chains: (i) for the ladder the confinement-deconfinement transition is in fact a crossover as far as the ground state is concerned LeHur ; Donohue ; (ii) in the ladder both the relevant and irrelevant cases lead to an insulating state due to a gap in the total charge excitation. Thus in the ladder to obtain a metallic state, as observed in the experimental (infinite number of chains) compounds, it would be necessary to explicitly dope the system and go away from the commensurate filling. This could be a way to mimic in this simplified model the small deviation of the commensurate filling due to warping of the Fermi surface perpendicular to chain direction Giamarchi\_physica .
Even if it is known that upon doping the ladder becomes metallic and the confined region is suppressed Suzumura\_P , it is yet very unclear what are the full properties of the system upon doping and what are the characteristics of such a metal-insulator transition. The purpose of the present paper is to examine these questions in the ladder system. We show that the metal-insulator transition in the ladder can be accurately described by a commensurate-incommensurate transition, and we determine its characteristics. The commensurate-incommensurate transition is well known in the classical case for one-dimension Fukuyama ; Okwamoto or quasi-one-dimensional Itakura system. However the quantum case which corresponds to interacting electron systems, is known only for one-dimensional case Haldane ; Schulz1983 where its connections to the Mott transition have been investigated in details Giamarchi\_PRB ; Mori ; Giamarchi\_physica . In the ladder, although both the half-filled commensurate system (spin ladder) and the extremely incommensurate case have been widely investigated Fabrizio ; KR ; Nagaosa ; Balents ; Schulz , comparatively little has been done on the metal-insulator transition close to half-filling. The ladder system with the umklapp scattering has been examined in the chain basis, where the incommensurate phase shows no gap in the total charge fluctuation KR . Studies using either a mapping on a $`SO(8)`$ symmetric model lin\_so8 ; konik\_exact\_commensurate\_ladder , onto a hard core boson system schulz\_mitwochain or in the large interchain hopping limit Ledermann show a drastic modification of the universal properties of the metal-insulator transition compared to the single-chain case. The universal properties close to half-filling have been checked numerically by DMRG Siller . We use here the bosonization technique and renormalization group to study the full problem as a function of the doping and the strength of the interchain hopping. This allows us to obtain the full phase diagram and in particular the interplay between the confinement-deconfinement at commensurate filling and the metal-insulator transition upon doping the ladder.
The present paper is organized as follows. In Sec. II, formulation is given in terms of a bosonized phase Hamiltonian and renormalization group equations are derived by assuming scaling invariance for response functions. In Sec. III, the phase diagram of the commensurate state and the incommensurate state is calculated by integrating the renormalization group equations. The charge gap is also estimated and the critical properties of the transition are given. In Sec. IV a summary and a discussion of the results can be found. Technical details can be found in the Appendix.
## 2 Formulation
### 2.1 System at half-filling
We consider two-coupled chains of a quarter-filled Hubbard model with a dimerization given by Tsuchiizu\_PHB
$``$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{\sigma =,}{}}{\displaystyle \underset{l=1,2}{}}\left[t+(1)^jt_\mathrm{d}\right]\left(c_{j,\sigma ,l}^{}c_{j+1,\sigma ,l}+\text{h.c.}\right)`$ (1)
$`2t_{}{\displaystyle \underset{j}{}}{\displaystyle \underset{\sigma =,}{}}\left(c_{j,\sigma ,1}^{}c_{j,\sigma ,2}+\text{h.c.}\right)`$
$`+U{\displaystyle \underset{j,l}{}}n_{j,,l}n_{j,,l},`$
where $`\sigma `$ ($`=,`$ or $`+,`$) and $`l`$ $`(=1,2)`$ denote the spin and chain index, respectively, and $`t_\mathrm{d}`$ is the dimerization in the one-dimensional chains. After the diagonalization of the $`t_\mathrm{d}`$-term, the kinetic term is written as
$`_0`$ $`=`$ $`{\displaystyle \underset{k,\sigma ,l}{}}\epsilon _k[d_{k,\sigma ,l}^{}d_{k,\sigma ,l}u_{k,\sigma ,l}^{}u_{k,\sigma ,l}]`$
$`2t_{}{\displaystyle \underset{k,\sigma }{}}[d_{k,\sigma ,1}^{}d_{k,\sigma ,2}+\mathrm{h}.\mathrm{c}.]`$
$`2t_{}{\displaystyle \underset{k,\sigma }{}}[u_{k,\sigma ,1}^{}u_{k,\sigma ,2}+\mathrm{h}.\mathrm{c}.],`$
where $`\epsilon _k=2\sqrt{t^2\mathrm{cos}^2ka+t_\mathrm{d}^2\mathrm{sin}^2ka}`$ Penc with the lattice constant $`a`$ and $`d_{k,\sigma ,l}`$ ($`u_{k,\sigma ,l}`$) is the fermion operator for the lower (upper) band on the $`l`$-th chain. Note that the umklapp scattering is induced by the dimerization, which has an effect of reducing the quarter-filled band into an effectively half-filled one. Since we consider only the lower band, we use half-filling instead of quarter-filling in the present paper. By diagonalizing the interchain hopping term with the use of $`a_{k,\sigma ,\zeta }=(\zeta d_{k,\sigma ,1}+d_{k,\sigma ,2})/\sqrt{2}`$ ($`\zeta =\pm `$), we obtain $`_0=_{k,\sigma ,\zeta }\epsilon (k,\zeta )a_{k,\sigma ,\zeta }^{}a_{k,\sigma ,\zeta }`$ where the energy dispersion is given by
$`\epsilon (k,\zeta )`$ $`=`$ $`2\sqrt{t^2\mathrm{cos}^2ka+t_\mathrm{d}^2\mathrm{sin}^2ka}2t_{}\zeta .`$ (3)
Here we introduce a linear dispersion, $`ϵ(k,\zeta )v_\mathrm{F}(pkk_{F\zeta })`$ where $`p`$ is the index of the branch $`p=+()`$ for the right moving (left moving) electrons and the new Fermi point is given by $`k_{\mathrm{F}\zeta }=k_\mathrm{F}\zeta 2t_{}/v_\mathrm{F}`$. We have neglected the $`t_{}`$-dependence of the Fermi velocity. After the bosonization of electrons around the new Fermi point $`k_{\mathrm{F}\zeta }`$, we introduce the phase variables defined by Tsuchiizu
$`\theta _{\nu \pm }(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{q0}{}}{\displaystyle \frac{\pi i}{qL}}\mathrm{e}^{\frac{\alpha }{2}|q|iqx}{\displaystyle \underset{k,\sigma ,\zeta }{}}f(\sigma ,\zeta )`$ (4)
$`\times `$ $`\left(a_{k+q,+,\sigma ,\zeta }^{}a_{k,+,\sigma ,\zeta }\pm a_{k+q,,\sigma ,\zeta }^{}a_{k,,\sigma ,\zeta }\right),`$
where $`f(\sigma ,\zeta )=1`$ (for $`\nu =\rho `$), $`\sigma `$ (for $`\nu =\sigma `$), $`\zeta `$ (for $`\nu =C`$) and $`\sigma \zeta `$ (for $`\nu =S`$), respectively. The phase variables $`\theta _{\rho +}`$ and $`\theta _{\sigma +}`$ ($`\theta _{\mathrm{C}+}`$ and $`\theta _{\mathrm{S}+}`$), express fluctuations of the total (transverse) charge density and spin density. They satisfy the boson commutation relation given by $`[\theta _{\nu +}(x),\theta _\nu ^{}(x^{})]=i\pi \delta _{\nu ,\nu ^{}}\mathrm{sgn}(xx^{})`$. In terms of these phase variables, the field operator, $`\psi _{p,\sigma ,\zeta }=L^{1/2}_k\mathrm{e}^{ikx}`$ $`a_{k,p,\sigma ,\zeta }`$, is given by Tsuchiizu
$`\psi _{p,\sigma ,\zeta }(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \alpha }}}\mathrm{exp}\left(ipk_{\mathrm{F}\zeta }x+i\mathrm{\Theta }_{p,\sigma ,\zeta }\right)\mathrm{exp}\left(i\pi \mathrm{\Xi }_{p,\sigma ,\zeta }\right),`$
where $`\alpha `$ is of the order of the lattice constant and $`\mathrm{\Theta }_{p,\sigma ,\zeta }=p/(2\sqrt{2})[\theta _{\rho +}+p\theta _\rho +\sigma (\theta _{\sigma +}+p\theta _\sigma )+\zeta (\theta _{\mathrm{C}+}+p\theta _\mathrm{C})+\sigma \zeta (\theta _{\mathrm{S}+}+p\theta _\mathrm{S})]`$. The phase factor, $`\pi \mathrm{\Xi }_{p,\sigma ,\zeta }`$, in Eq. (LABEL:eqn:field), is introduced to ensure the anticommutation relation for $`\psi _{p,\sigma ,\zeta }`$ with different $`p`$, $`\sigma `$ and $`\zeta `$.
### 2.2 System close to half-filling
In order to consider the system, which is slightly away from half-filling, we consider the following additional term,
$`_\mu `$ $`=`$ $`\mu {\displaystyle \underset{j,\sigma ,l}{}}c_{j,\sigma ,l}^{}c_{j,\sigma ,l}`$ (6)
$`=`$ $`\mu {\displaystyle \frac{\sqrt{2}}{\pi }}{\displaystyle 𝑑x_x\theta _{\rho +}},`$
where $`\mu `$ is the chemical potential and $`\mu =0`$ corresponds to the half-filling. We apply the transformation $`\sqrt{2}\theta _{\rho +}\sqrt{2}\theta _{\rho +}+q_0x`$ with $`q_0=4\mu K_\rho /v_\rho `$, which leads to a misfit term $`q_0x`$ in the cosine term expressing the umklapp scattering. The phase Hamiltonian is written as,
$``$ $`=`$ $`{\displaystyle \underset{\nu =\rho ,\sigma ,\mathrm{C},\mathrm{S}}{}}{\displaystyle \frac{v_\nu }{4\pi }}{\displaystyle 𝑑x\left[\frac{1}{K_\nu }\left(\theta _{\nu +}\right)^2+K_\nu \left(\theta _\nu \right)^2\right]}`$ (7)
$`+{\displaystyle \frac{g_\rho }{4\pi ^2\alpha ^2}}{\displaystyle 𝑑x\left[\mathrm{cos}\left(\sqrt{2}\theta _{\mathrm{C}+}\frac{8t_{}}{v_\mathrm{F}}x\right)+\mathrm{cos}\sqrt{2}\theta _\mathrm{C}\right]}`$
$`\times \left[\mathrm{cos}\sqrt{2}\theta _{\mathrm{S}+}\mathrm{cos}\sqrt{2}\theta _\mathrm{S}\right]`$
$`+{\displaystyle \frac{g_\sigma }{4\pi ^2\alpha ^2}}{\displaystyle 𝑑x\left[\mathrm{cos}\left(\sqrt{2}\theta _{\mathrm{C}+}\frac{8t_{}}{v_\mathrm{F}}x\right)\mathrm{cos}\sqrt{2}\theta _\mathrm{C}\right]}`$
$`\times \left[\mathrm{cos}\sqrt{2}\theta _{\mathrm{S}+}+\mathrm{cos}\sqrt{2}\theta _\mathrm{S}\right]`$
$`+{\displaystyle \frac{g_u}{2\pi ^2\alpha ^2}}{\displaystyle 𝑑x\mathrm{sin}\left(\sqrt{2}\theta _{\rho +}+q_0x\right)}`$
$`\times [\mathrm{cos}(\sqrt{2}\theta _{\mathrm{C}+}{\displaystyle \frac{8t_{}}{v_\mathrm{F}}}x)+\mathrm{cos}\sqrt{2}\theta _\mathrm{C}`$
$`\mathrm{cos}\sqrt{2}\theta _{\mathrm{S}+}+\mathrm{cos}\sqrt{2}\theta _\mathrm{S}]`$
$`+{\displaystyle \frac{g_{}}{2\pi ^2\alpha ^2}}{\displaystyle 𝑑x\mathrm{cos}\sqrt{2}\theta _{\sigma +}}`$
$`\times [\mathrm{cos}(\sqrt{2}\theta _{\mathrm{C}+}{\displaystyle \frac{8t_{}}{v_\mathrm{F}}}x)\mathrm{cos}\sqrt{2}\theta _\mathrm{C}`$
$`\mathrm{cos}\sqrt{2}\theta _{\mathrm{S}+}\mathrm{cos}\sqrt{2}\theta _\mathrm{S}],`$
where $`v_{\rho (\sigma )}=v_\mathrm{F}[1+()U/\pi v_\mathrm{F}]^{1/2}`$, $`v_{\mathrm{C}(\mathrm{S})}=v_\mathrm{F}`$, $`K_{\rho (\sigma )}=[1+()U/\pi v_\mathrm{F}]^{1/2}`$, $`K_{\mathrm{C}(\mathrm{S})}=1`$, $`g_\rho =g_\sigma =g_{}=Ua`$ and a coupling constant for the umklapp scattering is given by $`g_3=Ua(2t_\mathrm{d}/t)/[1+(t_\mathrm{d}/t)^2]`$ Penc .
### 2.3 Renormalization group equations
By utilizing a renormalization group method, we analyze Eq. (7) where the nonlinear terms in Eq. (7) are rewritten as
$$\frac{1}{2\pi ^2\alpha ^2}g_{\nu p,\nu ^{}p^{}}𝑑x\mathrm{cos}\sqrt{2}\overline{\theta }_{\nu p}\mathrm{cos}\sqrt{2}\overline{\theta }_{\nu ^{}p^{}},$$
(8)
where $`\overline{\theta }_{\rho +}=\theta _{\rho +}+q_0x/\sqrt{2}`$, $`\overline{\theta }_{C+}=\theta _{\rho +}+4\sqrt{2}t_{}x/v_\mathrm{F}`$ and $`\overline{\theta }_{\nu p}=\theta _{\nu p}`$ otherwise. The coupling constants are given by $`g_{\rho +,C\pm }=g_{\rho +,S\pm }=g_u`$, $`\pm g_{\sigma +,C\pm }=g_{\sigma +,S\pm }=g_{}`$, $`\pm g_{C\pm ,S\pm }=(g_\rho +g_\sigma )/2`$ and $`\pm g_{C\pm ,S}=(g_\rho g_\sigma )/2`$, where each coupling constant is treated in the renormalization group method. The renormalization group equations are derived from the response functions, which are given by $`T_\tau \mathrm{exp}[i\theta _{\rho +}(x_1,\tau _1)]\mathrm{exp}[i\theta _{\rho +}(x_2,\tau _2)]`$ Giamarchi\_JPF , By making use of the scaling of the cutoff ($`\alpha \alpha ^{}=\alpha \mathrm{e}^{dl}`$), we obtain the equations as (Appendix A),
$`{\displaystyle \frac{d}{dl}}\stackrel{~}{t}_{}`$ $`=`$ $`\stackrel{~}{t}_{}{\displaystyle \frac{1}{8}}G_{\rho +,C+}^2K_CF_1(8\stackrel{~}{t}_{};q_0\alpha )`$
$`{\displaystyle \frac{1}{8}}\left(G_{\sigma +,C+}^2+G_{C+,S+}^2+G_{C+,S}^2\right)K_CJ_1(8\stackrel{~}{t}_{}),`$
$`{\displaystyle \frac{d}{dl}}q_0\alpha `$ $`=`$ $`q_0\alpha G_{\rho +,C+}^2K_\rho F_1(q_0\alpha ;8\stackrel{~}{t}_{})`$
$`\left(G_{\rho +,C}^2+G_{\rho +,S+}^2+G_{\rho +,S}^2\right)K_\rho J_1(q_0\alpha ),`$
$`{\displaystyle \frac{d}{dl}}K_\rho `$ $`=`$ $`{\displaystyle \frac{1}{2}}K_\rho ^2[G_{\rho +,C+}^2F_0(8\stackrel{~}{t}_{};q_0\alpha )`$ (11)
$`+(G_{\rho +,C}^2+G_{\rho +,S+}^2+G_{\rho +,S}^2)J_0(q_0\alpha )],`$
$`{\displaystyle \frac{d}{dl}}K_\sigma `$ $`=`$ $`{\displaystyle \frac{1}{2}}K_\sigma ^2[G_{\sigma +,C+}^2J_0(8\stackrel{~}{t}_{})+G_{\sigma +,C}^2`$ (12)
$`+G_{\sigma +,S+}^2+G_{\sigma +,S}^2],`$
$`{\displaystyle \frac{d}{dl}}K_C`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p=\pm }{}}[K_C^2J_0(8\stackrel{~}{t}_{})\delta _{p,+}\delta _{p,}](G_{Cp,S+}^2`$ (13)
$`+G_{Cp,S}^2+G_{\sigma +,Cp}^2)`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p=\pm }{}}[K_C^2F_0(8\stackrel{~}{t}_{};q_0\alpha )\delta _{p,+}`$
$`J_0(q_0\alpha )\delta _{p,}]G_{\rho +,Cp}^2,`$
$`{\displaystyle \frac{d}{dl}}K_S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p=\pm }{}}(K_S^2\delta _{p,+}\delta _{p,})[G_{\rho +,Sp}^2J_0(q_0\alpha )`$ (14)
$`+G_{\sigma +,Sp}^2+G_{C+,Sp}^2J_0(8\stackrel{~}{t}_{})+G_{C,Sp}^2],`$
$`{\displaystyle \frac{d}{dl}}G_{\nu +,C\pm }`$ $`=`$ $`\left(2K_\nu K_C^\pm \right)G_{\nu +,C\pm }`$ (15)
$`G_{\nu +,S+}G_{C\pm ,S+}G_{\nu +,S}G_{C\pm ,S},`$
$`{\displaystyle \frac{d}{dl}}G_{\nu +,S\pm }`$ $`=`$ $`\left(2K_\nu K_S^\pm \right)G_{\nu +,S\pm }`$ (16)
$`G_{\nu +,C+}G_{C+,S\pm }J_0(8\stackrel{~}{t}_{})`$
$`G_{\nu +,C}G_{C,S\pm },`$
$`{\displaystyle \frac{d}{dl}}G_{Cp,Sp^{}}`$ $`=`$ $`\left(2K_C^pK_S^p^{}\right)G_{Cp,Sp^{}}`$ (17)
$`G_{\rho +,Cp}G_{\rho +,Sp^{}}J_0(q_0\alpha )`$
$`G_{\sigma +,Cp}G_{\sigma +,Sp^{}},`$
where $`\nu =\rho ,\sigma `$ and $`p=\pm `$. The quantities $`F_0(x;y)`$ and $`F_1(x;y)`$ are defined by
$`F_0(x;y)`$ $``$ $`{\displaystyle \frac{1}{2}}[J_0(|x+y|)+J_0(|xy|)],`$
$`F_1(x;y)`$ $``$ $`{\displaystyle \frac{1}{2}}[J_1(|x+y|)\mathrm{sgn}(x+y)`$
$`+J_1(|xy|)\mathrm{sgn}(xy)],`$
and $`J_n`$ is the $`n`$-th order Bessel function. In the above equations, the $`l`$-dependence is not written explicitly. The initial conditions are given by $`K_\nu (0)=K_\nu `$, $`G_{\nu p,\nu ^{}p^{}}(0)=g_{\nu p,\nu ^{}p^{}}/2\pi v_\mathrm{F}`$, $`\stackrel{~}{t}_{}(0)=t_{}/W`$ with $`Wv_\mathrm{F}\alpha ^1`$, and $`q_0(0)=q_0`$, respectively . The cutoff $`\alpha `$ can be related to the lattice constant by $`\alpha =2a/\pi `$ Tsuchiizu\_PHB , which leads to $`W\pi t/\sqrt{2}`$ for small $`t_\mathrm{d}/t`$. We take $`a`$=1. The renormalization group equations for the velocity $`v_\nu `$ are discarded and the velocity $`v_\rho `$ and $`v_\sigma `$ are set to $`v_\mathrm{F}`$.
In the limit of $`t_{}0`$ these equations reduce to those of a single chain given by
$`{\displaystyle \frac{d}{dl}}q_0\alpha `$ $`=`$ $`q_0\alpha 4G_u^2J_1(q_0\alpha ),`$ (18)
$`{\displaystyle \frac{d}{dl}}G_\rho `$ $`=`$ $`2G_u^2J_0(q_0\alpha ),`$ (19)
$`{\displaystyle \frac{d}{dl}}G_u`$ $`=`$ $`2G_\rho G_u,`$ (20)
$`{\displaystyle \frac{d}{dl}}G_\sigma `$ $`=`$ $`2G_{}^2,`$ (21)
$`{\displaystyle \frac{d}{dl}}G_{}`$ $`=`$ $`2G_\sigma G_{},`$ (22)
where $`K_\rho =1G_\rho `$ and $`K_\sigma =1G_\sigma `$. In the above equations, $`G_uG_{\rho +,C\pm }=G_{\rho +,S\pm }`$, $`G_{}\pm G_{\sigma +,C\pm }=G_{\sigma +,S\pm }`$, $`(G_\rho +G_\sigma )\pm 2G_{C\pm ,S\pm }`$ and $`(G_\rho G_\sigma )2G_{C\pm ,S}`$. It is straightforward to derive Eqs. (18)-(22) from the one-dimensional (1D) Hamiltonian given by (Appendix A)
$`_{1\mathrm{D}}`$ $`=`$ $`{\displaystyle \frac{v_\rho }{4\pi }}{\displaystyle 𝑑x\left[\frac{1}{K_\rho }\left(_x\theta _+\right)^2+K_\rho \left(_x\theta _{}\right)^2\right]}`$ (23)
$`+{\displaystyle \frac{v_\sigma }{4\pi }}{\displaystyle 𝑑x\left[\frac{1}{K_\sigma }\left(_x\varphi _+\right)^2+K_\sigma \left(_x\varphi _{}\right)^2\right]}`$
$`{\displaystyle \frac{\mu }{\pi }}{\displaystyle 𝑑x_x\theta _+}`$
$`+{\displaystyle \frac{g_u}{2\pi ^2\alpha ^2}}{\displaystyle 𝑑x\mathrm{cos}2\theta _+}`$
$`+{\displaystyle \frac{g_{}}{2\pi ^2\alpha ^2}}{\displaystyle 𝑑x\mathrm{cos}2\varphi _+},`$
where $`\theta _\pm `$ and $`\varphi _\pm `$ are the phase variables expressing the charge and spin fluctuations, respectively Suzumura\_PTP .
## 3 Commensurate and incommensurate states
We examine the commensurate-incommensurate transition by solving the renormalization group equations numerically. The dimerization is taken as $`t_\mathrm{d}/t=0.05`$ where such a choice does not change qualitatively the results as seen later. The scaling quantity $`l`$ is related to energy $`\omega `$ and/or temperature by $`l=\mathrm{ln}(W/\omega )=\mathrm{ln}(W/T)`$.
Equation (LABEL:eq:dl-q0) shows that the quantity $`q_0(l)`$ in the absence of the interaction increases as $`q_0(l)=q_0\mathrm{e}^l`$, while the increase of $`q_0(l)`$ is suppressed by the presence of $`G_{\rho +,C\pm }`$ and $`G_{\rho +,S\pm }`$ coming from the umklapp scattering $`g_u`$. In Fig. 1, $`q_0(l)`$ are shown with several choices of $`\mu /t`$ for $`U/t=5`$, $`t_\mathrm{d}/t=0.05`$ and $`t_{}/t=0.1`$. For $`\mu /t=0.15`$ (curve (1)), the misfit term $`q_0(l)`$ increases rapidly and becomes relevant. In this case, the umklapp scattering can be neglected since the Bessel functions $`J_0(q_0\alpha )`$ and $`J_1(q_0\alpha )`$ in Eqs. (LABEL:eq:dl-tperp)-(17) become small due to large $`q_0\alpha `$. The quantity $`K_\rho (l)`$ at the fixed point takes a finite value due to the relevant $`q_0\alpha `$, and then the incommensurate state is obtained. For smaller value of the chemical potential given by $`\mu /t=0.1`$ (curve(4)), the quantity $`q_0(l)\alpha `$ has a maximum and reduces to zero with increasing $`l`$, showing the irrelevant $`q_0(l)\alpha `$. Such a behavior of $`q_0(l)\alpha `$ leads to the commensurate state, which gives an insulating state due to the zero value of $`K_\rho (l)`$ in the limit of large $`l`$. We note that the commensurate state corresponds to that of a half-filled case. The commensurate-incommensurate transition occurs at the critical value given by $`\mu /t=\mu _c/t(0.125)`$ (curve (3)).
### 3.1 Half-filled case
It is known that a transition from the relevant interchain hopping to the irrelevant one occurs with increasing $`t_{}`$ at half-filling Suzumura ; Tsuchiizu . We show that such a result is also obtained for $`\mu 0`$, which still leads to the commensurate state. In Fig. 2, the quantity $`t_{}(l)/t`$ is calculated for the commensurate state with $`\mu /t=0.1`$. With decreasing the interchain hopping as $`t_{}/t=0.1`$(a), 0.057(b), 0.056(c) and 0.02(d), the hopping changes from a divergent behavior to a convergent one with zero value. One finds the relevant interchain hopping for large $`t_{}`$ (curve (a)) but irrelevant one for small $`t_{}`$ (curve (d)). The behavior which separates these two cases is obtained at a critical value of $`t_{}=t_{,c}`$ (curve(c)). For a relevant interchain hopping, one obtains an insulating phase with spin gap excitations which is called “D-Mott” phase lin\_so8 with “C0S0” phase. The notation C$`n`$S$`m`$ denotes a state with $`n`$ massless charge mode and $`m`$ massless spin mode. Such a spin gap is also known in the $`SO(8)`$ Gross-Neveu model lin\_so8 , which includes an additional term given by $`g_{\rho +,\sigma +}\mathrm{sin}\sqrt{2}\theta _{\rho +}\mathrm{cos}\sqrt{2}\theta _{\sigma +}`$. For an irrelevant interchain hopping, one can also expect a state with the spin gap due to the following fact. The phase has all charge excitations gapped, and because of such generated Heisenberg exchange ($`Jt_{}^2/\mathrm{\Delta }_\rho `$) LeHur is equivalent to a spin ladder. Such a system is known to have all spin excitations gapped. We note that, for two-coupled chains, these states undergo a crossover since both states exhibit C0S0 phase and there is no clear distinction between these two states LeHur ; Donohue . However we may identify the irrelevant (relevant) $`t_{}`$ as confinement (deconfinement) since such a result could be expected for the case of many chains Suzumura\_JPCS and higher dimension.
### 3.2 Commensurate-incommensurate transition
Based on Figs. 1 and 2, we examine quantities $`q`$ and $`t_{}^{\mathrm{eff}}`$, which express effective quantities of $`q_0`$ and $`t_{}`$ in the presence of the interaction and the chemical potential. These quantities are estimated from $`q\alpha =\mathrm{exp}[l_q]\times 4`$ with $`q_0(l_q)\alpha =4`$, and $`t_{}^{\mathrm{eff}}=t\mathrm{exp}[l_1]`$ with $`t_{}(l_1)/t=1`$. One finds that $`q=q_0`$ and $`t_{}^{\mathrm{eff}}=t_{}`$ for the noninteracting case. In the present analysis, the effective quantity $`q`$ is related to the carrier density $`n`$ by $`n=q/\pi `$ (Appendix A). In Fig. 3, the $`\mu `$-dependence of $`q`$ is shown with fixed $`t_{}/t=0`$, 0.05 and 0.2 for $`U/t=5`$ and $`t_\mathrm{d}/t=0.05`$. The quantity $`q`$, which is suppressed by the umklapp scattering, becomes zero in the commensurate phase. The thin dotted line denotes $`q(=q_0)`$ in the absence of the umklapp scattering. The critical value $`\mu _c`$ for the commensurate-incommensurate transition is not monotonical as a function of $`t_{}`$ since $`\mu _c`$ for $`t_{}/t=0`$ is larger (smaller) than that for $`t_{}/t=0.05`$ (0.2). The explicit $`t_{}`$-dependence of $`\mu _c`$ is evaluated in Fig. 4. In the inset of Fig. 3, the effective interchain hopping normalized by $`t_{}`$ is shown for $`t_{}/t=0.05`$ (dotted curve) and 0.2 (solid curve) where $`t_{}^{\mathrm{eff}}/t_{}`$ is symmetric with respect to $`\mu =0`$. For $`t_{}/t=0.05`$, $`t_{}^{\mathrm{eff}}/t_{}`$ becomes zero for small $`\mu /t(\stackrel{<}{}\mathrm{\hspace{0.17em}0.11})`$ and exhibits the irrelevant interchain hopping while the interchain hopping for $`t_{}/t=0.2`$ shows always relevant one.
From the results of Fig. 3, we examine the boundary between the commensurate state (C<sub>I</sub> and C<sub>II</sub>) and the incommensurate state (IC) and also the boundary between the relevant interchain hopping and irrelevant one. In Fig. 4, the phase diagram of these states is shown for $`U/t=5`$ and $`t_\mathrm{d}/t=0.05`$. The solid curve denotes the boundary between the incommensurate state and the commensurate state. The incommensurate state corresponds to the metallic state. The commensurate state is the (Mott-)insulating state where the dashed curve in the commensurate state denotes the boundary between the relevant interchain hopping (“D-Mott phase”) LeHur and the irrelevant one (“Confined phase”). The critical value shown by the arrow on the $`t_{}`$-axis is given by $`t_{,c}^0/t0.072`$ Tsuchiizu\_PHB ; Suzumura\_JPCS . The critical value shown by the arrow on the $`\mu `$-axis is given by $`\mu _c^0/t0.154`$, at which the solid curve merges with the dashed curve. The interchain hopping is always relevant in the incommensurate phase since the umklapp scattering becomes irrelevant in the limit of low energy as seen from Eq. (7). The incommensurate phase shows no gap in the total charge fluctuation due to the relevant $`\mu `$ KR , while the gaps still exist for other fluctuations due to the relevant $`t_{}`$. The incommensurate state corresponds to a “C1S0” $`d`$-wave superconductivity Schulz and is called “Luther-Emery liquid” Ledermann . The solid curve corresponding to the critical value $`\mu _c`$ shows non-monotonical behavior as a function of $`t_{}`$. With increasing $`t_{}`$, $`\mu _c`$ has a minimum at $`t_{}/t0.07`$ and becomes larger than $`\mu _c^0`$. The decrease of $`\mu _c`$ for small $`t_{}/t(\stackrel{<}{}\mathrm{\hspace{0.17em}0.07})`$ originates in the fact that the effect of the umklapp scattering becomes weakened for finite $`t_{}`$ due to the misfit $`(8t_{}/v_\mathrm{F})x`$ in $`g_u`$-term of Eq. (7).
For large $`t_{}`$ and small $`g_u`$ ($`g_u`$ is controlled by $`t_\mathrm{d}`$ which is small here) we can explain the behavior by a qualitative analysis of the renormalization group equations. The strong interchain hopping between the two interacting chains opens a gap in all sectors except for the total charge sector. We may define a scale $`l_1`$ where the couplings in the $`\sigma `$, $`C`$, $`S`$ sectors have reached a value of order one. This scale will depend on $`t_{}`$ and the value of the bare interaction but not on $`g_u`$. Up to $`l_1`$, $`g_u`$ will renormalize by some finite multiplicative constant that will not affect the asymptotic dependence of the charge gap. However above $`l_1`$ the scaling dimension of $`g_u`$ is $`2K_\rho `$, instead of $`22K_\rho `$ if $`t_{}=0`$. Indeed $`\mathrm{cos}(\sqrt{2}\theta _{S+})`$ and $`\mathrm{cos}(\sqrt{2}\theta _C)`$ acquire a mean value (with opposite signs) so that the umklapp term reduces to:
$$g_u\mathrm{sin}\sqrt{2}\theta _{\rho +}\left(\mathrm{cos}\sqrt{2}\theta _C\mathrm{cos}\sqrt{2}\theta _{S+}\right).$$
(24)
As a first consequence this means that the power law dependence of the gap for asymptotically small umklapp scattering is enhanced by interchain hopping from $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}g_u^{1/(22K_\rho )}`$ to $`\mathrm{\Delta }_\rho g_u^{1/(2K_\rho )}`$. The second physical consequence is for the commensurate-incommensurate transition. In the absence of chemical potential of Eq. (6), the operator $`g_u`$ in Eq. (24) would be relevant when $`K_\rho ^{\mu =0}<2`$. As is standard for the commensurate-incommensurate transition Schulz ; Giamarchi\_physica , the addition of a chemical potential of Eq. (6) leads to a universal value of the Luttinger exponent close to zero doping of $`K_\rho ^{}=K_\rho ^{\mu =0}/2=1`$. It is also easy to see that the charge excitations connecting two minima of the potential, Eq. (24), correspond to a charge $`+2`$. We thus recover quite generally the results that were established in the various previous limits konik\_exact\_commensurate\_ladder ; schulz\_mitwochain ; Ledermann ; Siller .
Next, we estimate the magnitude of the charge gap on the plane of $`\mu /t`$ and $`t_{}/t`$ of Fig. 4. In a way similar to the previous calculations in the case for $`\mu =0`$ Suzumura ; Tsuchiizu , the charge gap $`\mathrm{\Delta }_\rho `$ is calculated from the renormalization group flow where $`\mathrm{\Delta }_\rho /t=W\mathrm{exp}[l_\mathrm{\Delta }]`$ and $`l_\mathrm{\Delta }`$ is defined by $`K_\rho (l_\mathrm{\Delta })1/2`$. We set $`\mathrm{\Delta }_\rho =0`$ for $`K_\rho (\mathrm{})>1/2`$. In Fig. 5, the $`t_{}`$-dependence of $`\mathrm{\Delta }_\rho `$ is shown for $`\mu /t=0`$, 0.1 and 0.14. When $`t_{}`$ increases, the charge gap with $`\mu =0`$ decreases and has a minimum at $`t_{}/t=0.154`$. A similar dependence is found in the $`t_{}`$-dependence of $`\mu _c`$ in Fig. 4 where the location for the minimum is the same within the numerical accuracy. The identification of the charge gap for the commensurate case $`\mathrm{\Delta }_\rho (\mu =0)`$ with $`\mu _c`$ can be expected for the commensurate-incommensurate transition on general grounds. Indeed the charge gap is the smallest energy one must pay to inject an excitation which carries a charge. When the chemical potential reaches the charge gap, charged excitations are injected in the system, the system is incommensurate. The non-monotonical dependences of $`\mathrm{\Delta }_\rho `$ on $`t_{}`$ indicate a crossover from the irrelevant interchain hopping to the relevant one with increasing $`t_{}`$. For $`\mu /t=0.1`$, $`\mathrm{\Delta }_\rho `$ is suppressed but is still similar to that of $`\mu =0`$. However the $`t_{}`$-dependence of $`\mathrm{\Delta }_\rho `$ for $`\mu /t=0.14`$ is qualitatively different from others. The metallic state with $`\mathrm{\Delta }_\rho /t=0`$ appears in the interval region of $`0.028<t_{}/t<0.2`$. This region corresponds to the incommensurate state in Fig. 4. We note that numerical integration of the renormalization group flow shows that $`\mathrm{\Delta }_\rho `$ as a function of $`\mu `$ decreases monotonically to zero at the commensurate-incommensurate transition, in agreement with the annihilation of the gap expected for a commensurate-incommensurate transition.
Finally we examine the $`U`$ and $`t_\mathrm{d}`$ dependences of the critical values, $`\mu _c^0`$ and $`t_{,c}^0`$, which are shown by the arrows in Fig. 4. In Fig. 6 the $`U`$-dependences of $`\mu _c^0`$ and $`t_{,c}`$ are shown by the solid curve and the dashed curve respectively for $`t_\mathrm{d}/t=0.05`$. In the inset, the $`t_\mathrm{d}`$-dependences of the corresponding quantities are shown by solid and dashed curves for $`U/t=5`$. These $`U`$ and $`t_\mathrm{d}`$-dependences show that $`\mu _c^0/t_{,c}^02.2`$ for the present choice of parameters. Thus a rescaled phase diagram, which is independent of $`U`$ and $`t_\mathrm{d}`$, can be obtained from Fig. 4 by using the rescaled variables $`\mu /\mu _c^0`$ and $`t_{}/t_{,c}^0`$. Our renormalization group procedure correctly gives a value for $`\mu _c^0`$ identical, within the scale of Fig. 4, with $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}`$. The quantity $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}`$ denotes $`\mathrm{\Delta }_\rho `$ of the single chain where $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}W(g_u/W)^{1/(22K_\rho )}`$ for small $`g_u`$ Schwartz . We note that the boundary between a relevant $`t_{}`$ and an irrelevant $`t_{}`$ is determined more accurately by the competition between $`t_{}^{\mathrm{eff},0}`$ and $`\mathrm{\Delta }_\rho ^{1\mathrm{D}}`$ where $`t_{}^{\mathrm{eff},0}`$ denotes the effective interchain hopping energy for $`t_\mathrm{d}=\mu =0`$ Tsuchiizu\_PRG ; Suzumura\_JPCS . The quantity $`t_{}^{\mathrm{eff},0}`$ is given analytically by $`t_{}^{\mathrm{eff},0}=t_{}(t_{}/t)^{\alpha _0/(1\alpha _0)}`$ with $`\alpha _0=(K_\rho +K_\rho ^1+K_\sigma +K_\sigma ^14)/4`$ Bourbonnais1 ; Bourbonnais2 ; Bourbonnais .
## 4 Summary and discussion
In the present paper, we investigated the properties of two-coupled chains with interchain electron hopping and a filling close to half-filling. By using bosonization and a renormalization group method we obtained the full phase diagram of the system. There is a metal-insulator transition for a critical value $`\mu _c`$ of the chemical potential, describable by a commensurate-incommensurate phase transition on the bosonized Hamiltonian. The critical value $`\mu _c`$, shown in Fig. 4, exhibits a non-monotonical dependence on the perpendicular hopping $`t_{}`$. The minimum of $`\mu _c`$ occurs for values of $`t_{}`$ close to the ones for which a confinement deconfinement crossover takes place for the commensurate case. For large $`t_{}`$ the relevance of interchain hopping reinforces the commensurate character of the system leading to an enhanced commensurability gap. The crossover line separating the regions of irrelevant and relevant interchain hopping (confinement-deconfinement line) merges with the boundary between the commensurate state and the incommensurate state at $`\mu =\mu _c^0`$, the critical chemical potential for a single chain ($`t_{}=0`$). We found that the phase diagram of Fig. 4 becomes almost independent of parameters such as interactions and dimerization when expressed in terms of the rescaled variables $`\mu /\mu _c^0`$ ($`t_{}/t_{,c}^0`$). In addition, given the form for the Hamiltonian, we could show that at the limit of small doping, the Luttinger liquid parameter takes the universal value $`K_\rho ^{}=1`$, thereby confirming the results lin\_so8 ; konik\_exact\_commensurate\_ladder ; schulz\_mitwochain ; Ledermann ; Siller obtained on specific limits of the model.
Finally let us discuss the interchain exchange interactions in the commensurate confined phase (irrelevant interchain hopping). Even when irrelevant, the interchain hopping generates two particles and particle hole hopping Bourbonnais . Within the present formalism, the two particle interchain hopping can be taken into account by starting from the chain. The Hamiltonian corresponding to the umklapp scattering, $`_u`$, in Eq. (7) is rewritten as
$`_u`$ $`=`$ $`g_u^{1\mathrm{D}}{\displaystyle \underset{p,l}{}}{\displaystyle 𝑑x\psi _{p,,l}^{}\psi _{p,,l}\psi _{p,,l}^{}\psi _{p,,l}}`$
$`+J_{uz}{\displaystyle \underset{p}{}}{\displaystyle }dx[\psi _{p,,1}^{}\psi _{p,,1}\psi _{p,,2}^{}\psi _{p,,2}+\mathrm{h}.\mathrm{c}.]`$
$`+J_u{\displaystyle \underset{p}{}}{\displaystyle }dx[\psi _{p,,1}^{}\psi _{p,,1}\psi _{p,,2}^{}\psi _{p,,2}+\mathrm{h}.\mathrm{c}.]`$
$`+g_u^{pt}{\displaystyle \underset{p}{}}{\displaystyle }dx[\psi _{p,,1}^{}\psi _{p,,1}^{}\psi _{p,,2}\psi _{p,,2}+\mathrm{h}.\mathrm{c}.],`$
where $`\psi _{p,\sigma ,l}=L^{1/2}_k\mathrm{e}^{ikx}d_{k,p,\sigma ,l}`$ with $`l=1,2`$ denoting the chain index and $`p=+()`$ corresponding to the right (left) moving electrons. The coupling constants are given by $`g_u^{1\mathrm{D}}=(g_{\rho +,C+}+g_{\rho +,C}g_{\rho +,S+}+g_{\rho +,S})/4`$, $`J_{uz}=(g_{\rho +,C+}g_{\rho +,C}g_{\rho +,S+}g_{\rho +,S})/4`$, $`J_u=(g_{\rho +,C+}g_{\rho +,C}+g_{\rho +,S+}+g_{\rho +,S})/4`$ and $`g_u^{pt}=(g_{\rho +,C+}+g_{\rho +,C}+g_{\rho +,S+}g_{\rho +,S})/4`$, where $`g_u^{1\mathrm{D}}`$ and $`J_{uz}`$ ($`J_u`$) denote the intrachain umklapp scattering and the interchain umklapp exchange, and $`g_u^{pt}`$ denotes the pair tunneling between chains. Other terms in the Hamiltonian can be rewritten in a similar form. The initial values of the interchain exchange and that of pair tunneling are zero since the initial values of the umklapp scattering in the renormalization group equations are given by $`g_{\rho +,C+}=g_{\rho +,C}=g_{\rho +,S+}=g_{\rho +,S}=g_u`$. Both of these interactions are generated by the renormalization.
It would be a interesting problem to investigate, by taking the higher order terms in the renormalization group, the consequences of these terms on the full phase diagram investigated in this paper.
###### Acknowledgements.
One of the authors (Y.S.) is thankful for the financial support from Université Paris–Sud and also for the kind hospitality during his stay at Ecole Normale Supérieure. M.T. and Y.S. thank T. Itakura for useful discussions. This work was partially supported by a Grant-in-Aid for Scientific Research from the Ministry of Education, Science, Sports and Culture (Grant No.09640429), Japan.
## Appendix A Renormalization group equations
In this section, we derive the renormalization group equation for $`\mu `$.
First, we treat the system with a single chain where the phase Hamiltonian is given by
$`_{1\mathrm{D}}`$ $`=`$ $`{\displaystyle \frac{v_\rho }{4\pi }}{\displaystyle 𝑑x\left[\frac{1}{K_\rho }\left(_x\theta _+\right)^2+K_\rho \left(_x\theta _{}\right)^2\right]}`$ (26)
$`+{\displaystyle \frac{g_u}{2\pi ^2\alpha ^2}}{\displaystyle 𝑑x\mathrm{cos}\left(2\theta _++q_0x\right)}.`$
The quantities $`v_\rho `$ and $`K_\rho `$ are the same as two-coupled chains, and $`[\theta _+(x),\theta _{}(x^{})]=i\pi \mathrm{sgn}(xx^{})`$ Suzumura\_PTP and $`q_0=(4K_\rho /v_\rho )\mu `$. Here the expectation value of the carrier density, $`n`$, can be evaluated as
$`n`$ $`=`$ $`{\displaystyle \frac{2K_\rho }{\pi v_\rho }}\mu +{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{T}{L}}{\displaystyle 𝑑x𝑑\tau _x\theta _+},`$ (27)
where the factor $`2K_\rho /\pi v_\rho `$ in the first term of r.h.s. corresponds to the compressibility in the absence of the umklapp scattering. The second term of Eq. (27) can be calculated as follows,
$`{\displaystyle 𝑑x𝑑\tau _x\theta _+}`$ (28)
$`=`$ $`{\displaystyle \frac{1}{Z}}\mathrm{Tr}\left[{\displaystyle \frac{}{\lambda }}\mathrm{exp}\left({\displaystyle 𝑑\tau _{1\mathrm{D}}}+\lambda {\displaystyle 𝑑x𝑑\tau _x\theta _+}\right)|_{\lambda =0}\right]`$
$`=`$ $`{\displaystyle \frac{4G_u}{\alpha ^2}}K_\rho {\displaystyle 𝑑x𝑑\tau x\mathrm{sin}\left(2\theta _++q_0x\right)},`$
where $`Z=\mathrm{Tr}\mathrm{exp}(𝑑\tau _{1\mathrm{D}})`$ and $`G_u=g_u/(2\pi v_\rho )`$. In Eq. (28), the new phase variable $`\stackrel{~}{\theta }_+(=\theta _+2\pi K_\rho \lambda x/v_\rho )`$ has been introduced and rewritten as $`\stackrel{~}{\theta }_+\theta _+`$. Then Eq. (27) leads
$`n`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}q_0+{\displaystyle \frac{4}{\pi \alpha ^2}}G_uK_\rho {\displaystyle \frac{T}{L}}{\displaystyle 𝑑x𝑑\tau x\mathrm{sin}\left(2\theta _++q_0x\right)}.`$
Here it is worthwhile noting that Eq. (LABEL:eq:nb) is compared with
$`\mu `$ $`=`$ $`{\displaystyle \frac{\pi v_\rho }{2K_\rho }}n{\displaystyle \frac{2v_\rho }{\alpha ^2}}G_u{\displaystyle \frac{T}{L}}{\displaystyle 𝑑x𝑑\tau x\mathrm{sin}\left(2\theta _++2\pi nx\right)},`$
which is obtained by using the Legendre transformation, i.e., by calculating the value of the chemical potential at fixed carrier density $`n`$ Giamarchi\_JPF ; Giamarchi\_PRB .
After a straightforward calculation of Eq. (LABEL:eq:nb), one finds, up to the lowest order of $`G_u`$,
$`n`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}q_0{\displaystyle \frac{2}{\pi \alpha }}G_u^2K_\rho {\displaystyle \frac{dr}{\alpha }\left(\frac{r}{\alpha }\right)^{24K_\rho }J_1(q_0r)}.`$ (31)
By assuming that the quantity $`n`$ is simply scaled as $`n(l)=n\mathrm{e}^l`$ Giamarchi\_JPF with the transform $`\alpha \alpha \mathrm{e}^l`$, the renormalization group equation for $`q_0`$ is obtained as,
$`{\displaystyle \frac{d}{dl}}q_0\alpha `$ $`=`$ $`q_0\alpha 4G_u^2K_\rho J_1(q_0\alpha ).`$ (32)
The renormalization group equations for $`K_\rho `$ and $`G_u`$ can be obtained in a way similar to Ref. Giamarchi\_PRB as
$`{\displaystyle \frac{d}{dl}}K_\rho `$ $`=`$ $`2G_u^2K_\rho ^2J_0(q_0\alpha ),`$ (33)
$`{\displaystyle \frac{d}{dl}}G_u`$ $`=`$ $`(22K_\rho )G_u.`$ (34)
By integrating this renormalization group equation, the effective quantity of $`q_0`$ can be estimated from $`q\alpha =c\mathrm{exp}(l_q)`$ with $`q_0(l_q)\alpha =c`$, where $`c`$ is numerical constant of the order of unity. One finds that the quantity $`q`$ can be related to carrier density $`n`$ by $`n=q/(2\pi )`$ from Eq. (31).
Next we calculate the case of two-coupled chains. The expectation value of the carrier density, $`n`$, can be evaluated as
$`n`$ $`=`$ $`{\displaystyle \frac{4K_\rho }{\pi v_\rho }}\mu +{\displaystyle \frac{\sqrt{2}}{\pi }}{\displaystyle \frac{T}{L}}{\displaystyle 𝑑x𝑑\tau _x\theta _{\rho +}}.`$ (35)
where the first term of r.h.s. becomes twice as that in Eq. (27), since here we consider two chains. From a procedure similar to the single chain, Eq. (LABEL:eq:nb) is replaced by
$`n`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}q_0+{\displaystyle \frac{4}{\pi \alpha ^2}}G_{\rho +,C+}K_\rho {\displaystyle \frac{T}{L}}{\displaystyle 𝑑x𝑑\tau x}`$ (36)
$`\times \mathrm{sin}\left(\sqrt{2}\theta _{\rho +}+q_0x\right)\mathrm{cos}\left(\sqrt{2}\theta _{C+}{\displaystyle \frac{8t_{}}{v_\mathrm{F}}}x\right)`$
$`+`$ $`{\displaystyle \frac{4}{\pi \alpha ^2}}G_{\rho +,C}K_\rho {\displaystyle \frac{T}{L}}{\displaystyle 𝑑x𝑑\tau x}`$
$`\times \mathrm{sin}\left(\sqrt{2}\theta _{\rho +}+q_0x\right)\mathrm{cos}\sqrt{2}\theta _C`$
$`+`$ $`{\displaystyle \frac{4}{\pi \alpha ^2}}G_{\rho +,S+}K_\rho {\displaystyle \frac{T}{L}}{\displaystyle 𝑑x𝑑\tau x}`$
$`\times \mathrm{sin}\left(\sqrt{2}\theta _{\rho +}+q_0x\right)\mathrm{cos}\sqrt{2}\theta _{S+}`$
$`+`$ $`{\displaystyle \frac{4}{\pi \alpha ^2}}G_{\rho +,S}K_\rho {\displaystyle \frac{T}{L}}{\displaystyle 𝑑x𝑑\tau x}`$
$`\times \mathrm{sin}\left(\sqrt{2}\theta _{\rho +}+q_0x\right)\mathrm{cos}\sqrt{2}\theta _S.`$
The coupling constants that appear in Eq. (36) are those including the misfit $`q_0x`$ in the cosine potential of Eq. (7). The calculation in the lowest order of perturbation yields,
$`n`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}q_0+{\displaystyle \frac{1}{\pi \alpha }}G_{\rho +,C+}^2K_\rho `$
$`\times {\displaystyle }{\displaystyle \frac{dr}{\alpha }}\left({\displaystyle \frac{r}{\alpha }}\right)^{22K_\rho 2K_C}F_1(q_0r;8t_{}r/v_\mathrm{F})`$
$`+{\displaystyle \frac{1}{\pi \alpha }}G_{\rho +,C}^2K_\rho {\displaystyle \frac{dr}{\alpha }\left(\frac{r}{\alpha }\right)^{22K_\rho 2/K_C}J_1(q_0r)}`$
$`+{\displaystyle \frac{1}{\pi \alpha }}G_{\rho +,S+}^2K_\rho {\displaystyle \frac{dr}{\alpha }\left(\frac{r}{\alpha }\right)^{22K_\rho 2K_S}J_1(q_0r)}`$
$`+{\displaystyle \frac{1}{\pi \alpha }}G_{\rho +,S}^2K_\rho {\displaystyle \frac{dr}{\alpha }\left(\frac{r}{\alpha }\right)^{22K_\rho 2/K_S}J_1(q_0r)},`$
where $`F_1(x;y)[J_1(|x+y|)\mathrm{sgn}(x+y)+J_1(|xy|)\mathrm{sgn}(xy)]/2`$. The infinitesimal transform of the cutoff $`\alpha \alpha ^{}=\alpha e^{dl}`$ in Eq. (LABEL:eq:n2c) leads the renormalization equation, Eq. (LABEL:eq:dl-q0). The renormalization group equation for $`t_{}`$, Eq. (LABEL:eq:dl-tperp), can be obtained in a similar way. The equations for the other coupling constants, Eqs. (11)-(17) are also obtained in a way similar to Ref. Tsuchiizu .
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# Phenomenology of Maximal and Near–Maximal Lepton Mixing
## I Introduction
The data from both atmospheric and solar neutrino experiments have given a rather convincing evidence for non-zero neutrino masses and mixing. Furthermore, some features of the neutrino flavour parameters are surprising and quite different from the quark flavour parameters. In particular, one or two of the three mixing angles in the MNS (Maki-Nakagawa-Sakata) mixing matrix for leptons are large. Specifically, the simplest interpretation of the atmospheric neutrino measurements gives
$$\mathrm{sin}^22\theta _{\mu \tau }0.81,\mathrm{\Delta }m_{32}^2(25)\times 10^3\text{eV}^2.$$
(1)
There exist several solutions of the solar neutrino problem involving oscillations of electron neutrinos into some combination, $`\nu _x`$, of $`\mu `$ and $`\tau `$neutrinos with large mixing angle with parameters in the range
$$\mathrm{sin}^22\theta _{ex}0.71,\mathrm{\Delta }m_{12}^2(0.24)\times 10^4\text{eV}^2\text{or}(0.0520)\times 10^8\text{eV}^2.$$
(2)
The mixing angles involved in the explanation of the solar and atmospheric neutrino data are not just order one. They are actually near–maximal, that is, $`\mathrm{sin}^22\theta `$ close to $`1`$. If indeed one of the mixing angles is near–maximal, it would provide a strong support to the idea that the corresponding neutrinos are in a pseudo-Dirac state. Such a scenario would have very interesting implications for theoretical flavour models. These implications have been recently studied in Ref. . A precise knowledge of the mixing and, in particular lower and/or upper bounds on small deviations from maximal mixing provides an excellent probe of the related flavour physics. It is the purpose of this work to study the experimental ways in which the region of near–maximal mixing can be probed.
It is important to notice that there exists also a viable solution of the solar neutrino problem that does not involve large mixing, the so called small mixing angle solution (SMA). Clearly, identification of the SMA solution will immediately exclude the possibility discussed in this paper. Also discovery of the sterile neutrinos will require a change of the whole picture. Here we consider only three light active neutrinos and take the $`3\times 3`$ MNS matrix to be unitary.
Our main interest lies in the study of near–maximal mixing involving $`\nu _e`$. We define a small parameter $`ϵ`$ which describes the deviation from maximal mixing as:
$$\mathrm{sin}^2\theta _{ex}\frac{1}{2}(1ϵ),|ϵ|1.$$
(3)
Our convention is such that $`\mathrm{\Delta }m_{21}^2m_2^2m_1^2>0`$. Then $`ϵ>0`$ ($`<0`$) corresponds to the case that the lighter neutrino mass eigenstate, $`\nu _1`$, has a larger (smaller) component of $`\nu _e`$ and the heavier one, $`\nu _2`$, has a larger (smaller) component of $`\nu _x`$.
Which value of deviation from maximal mixing is expected? In the case of theoretical models where the pseudo-Dirac structure is naturally induced, one expects that the deviation from maximal mixing is suppressed by a small parameter that is related to an approximate horizontal symmetry. If the same symmetry is also responsible for the smallness and hierarchy of the quark sector parameters, then it is quite plausible that $`|ϵ|𝒪(\lambda )`$, where $`\lambda =0.22`$ is the Cabibbo angle in the Wolfenstein parametrization. In various flavor models, the deviation from maximal mixing is related to other physical parameters. For example, in a large class of models, we have
$$|ϵ|\text{ }>2m_e/m_\mu 0.01.$$
(4)
The $`U(1)\times U(1)`$ models described in refs. give $`|ϵ|(\mathrm{\Delta }m_{12}^2/\mathrm{\Delta }m_{23}^2)^{1/3}`$, while models of $`L_eL_\mu L_\tau `$ symmetry give $`|ϵ|\mathrm{\Delta }m_{12}^2/\mathrm{\Delta }m_{23}^2`$. Quite generally we have $`|ϵ|\text{ }>\mathrm{\Delta }m_{12}^2/\mathrm{\Delta }m_{23}^2`$, which could be more restrictive than Eq. (4) if the solution to the solar neutrino problem lies at the upper end of the $`\mathrm{\Delta }m_{12}^2`$ range given in Eq. (2). In a large class of models we also have $`|ϵ|`$ of the same order of magnitude as $`|U_{e3}|`$, the mixing of $`\nu _e`$ in the third mass eigenstate, which can again be tested in the future .
Large values of $`ϵ`$, $`|ϵ|\text{ }>0.3`$, will testify against at least the simplest versions of these theories. Therefore we consider both positive and negative values of $`ϵ`$ in the range
$$|ϵ|\text{ }<0.3$$
(5)
which corresponds to $`\mathrm{sin}^22\theta _{ex}\text{ }>0.91`$. As concerns the mass-squared difference we cover the whole range below the reactor bound:
$$10^{11}\text{eV}^2\text{ }<\mathrm{\Delta }m_{21}^2\text{ }<10^3\text{eV}^2.$$
(6)
Most of our discussion takes place in an effective two neutrino generation framework. This is justified if $`U_{e3}`$ is zero or sufficiently small. We quantify this condition and consider the effect of a non-zero $`U_{e3}`$ in the Sect. VI. We find there that, for matter oscillations, the leading corrections to the case of maximal mixing are of $`𝒪(ϵ,|U_{e3}|^2)`$, so that a reduction to a two neutrino analysis is justified for $`|U_{e3}|^2ϵ`$. For vacuum oscillations, the corrections are of $`𝒪(ϵ^2,|U_{e3}|^2)`$, and a two neutrino analysis is valid for $`|U_{e3}|ϵ`$.
On the experimental side, we note that if $`\mathrm{\Delta }m_{\mathrm{atm}}^2>2\times 10^3\text{eV}^2`$, then the CHOOZ experiment limit requires $`|U_{e3}|^20.05`$ . For such values of $`|U_{e3}|^2`$, a two generation analysis of matter (vacuum) effects is always valid for $`|ϵ|>0.1(0.3)`$. On the theoretical side, in a large class of flavour models, $`|U_{e3}||ϵ|`$ . In such a framework, a two generation analysis of matter effects is a good approximation while vacuum oscillations should be considered in the three generation framework.
In the limit $`|U_{e3}|=0`$ (which reduces the problem to a two neutrino framework), the deviation from maximal mixing can be determined as in (3). In addition the mixing of $`\nu _\mu `$ and $`\nu _\tau `$ can also be maximal, as favored by the atmospheric neutrino data. In this case we have the mixing structure
$$|U_{e3}|^2=0,|U_{e1}|^2=|U_{e2}|^2=\frac{1}{2},|U_{\mu 3}|^2=|U_{\tau 3}|^2=\frac{1}{2},$$
(7)
which corresponds to the so called bi-maximal mixing scheme . The analysis for the solar neutrino phenomenology is independent of $`U_{\mu 3}`$ and therefore the results discussed in Secs. III, IV, V, VII B and VII C apply also to the bi-maximal mixing case. Only in Sec. VII A where we discuss atmospheric neutrinos, does $`U_{\mu 3}`$ play a role, and there we take it to be large and possibly maximal.
The plan of this paper is as follows. In Sec. II we give some basic physics considerations and useful expressions for the survival probability and for various observables. In Sec. III we describe the present status of maximal mixing from solar neutrino experiments. The results of a global fit to all available solar neutrino data are given in Sec. III A. The dependence of these results on various aspects of the analysis is described in Sec. III B. In Sec. IV we study the predictions that follow from near–maximal mixing for individual, existing measurements: total rates, Argon production rate, Germanium production rate, the Day–Night asymmetry in elastic scattering events, the zenith angle distribution of elastic scattering events, and the shape of the recoil energy spectrum. In Sec. V we suggest tests of maximal mixing from future experiments. We describe the conditions for having unambiguous tests in Sec. V A. Then we examine individual experiments: GNO and Super–Kamiokande, SNO, Borexino and low energy solar neutrino experiments. The effect of extending to three neutrino scenario is commented in Sec. VI. In Sec. VII we discuss the effect of maximal mixing in atmospheric neutrinos, supernova neutrinos and neutrinoless double beta decay. We present our conclusions in Sec. VIII.
## II Physics at near–maximal mixing
### A The survival probability for solar neutrinos
In this subsection we present general expressions for the survival probability of solar electron neutrino in a two generation framework valid in the full range of $`\mathrm{\Delta }m^2`$ which we use in our numerical calculations.
The survival amplitude for a solar $`\nu _e`$ neutrino of energy $`E`$ at a detector in the Earth can be written as:
$$A_{ee}=\underset{i=1}{\overset{2}{}}A_{ei}^SA_{ie}^E\mathrm{exp}[im_i^2(Lr)/2E].$$
(8)
Here $`A_{ei}^S`$ is the amplitude of the transition $`\nu _e\nu _i`$ ($`\nu _i`$ is the $`i`$-mass eigenstate) from the production point to the Sun surface, $`A_{ie}^E`$ is the amplitude of the transition $`\nu _i\nu _e`$ from the Earth surface to the detector, and the propagation in vacuum from the Sun to the surface of the Earth is given by the exponential. $`L`$ is the distance between the center of the Sun and the surface of the Earth, and $`r`$ is the distance between the neutrino production point and the surface of the Sun. The corresponding survival probability $`P_{ee}`$ is then given by:
$`P_{ee}`$ $`=`$ $`P_1P_{1e}+P_2P_{2e}+2\sqrt{P_1P_2P_{1e}P_{2e}}\mathrm{cos}\xi `$ (9)
$`=`$ $`P_1+(12P_1)P_{2e}+2\sqrt{P_1(1P_1)P_{2e}(1P_{2e})}\mathrm{cos}\xi .`$ (10)
Here $`P_i|A_{ei}^S|^2`$ is the probability that the solar neutrinos reach the surface of the Sun as $`|\nu _i`$ and we use $`P_1+P_2=1`$, while $`P_{ie}|A_{ie}^E|^2`$ is the probability of $`\nu _i`$ arriving at the surface of the Earth to be detected as a $`\nu _e`$, and we use $`P_{1e}+P_{2e}=1`$. The phase $`\xi `$ is given by
$$\xi =\frac{\mathrm{\Delta }m^2(Lr)}{2E}+\delta ,$$
(11)
where $`\delta `$ contains the phases due to propagation in the Sun and in the Earth and can be safely neglected since it is always much smaller than the preceeding term .
From Eq. (9) one can recover more familiar expressions for $`P_{ee}`$:
(1) For $`\mathrm{\Delta }m^2/E5\times 10^{17}`$ eV, the matter effect supresses flavour transitions in both the Sun and the Earth. Consequently, the probabilities $`P_1`$ and $`P_{2e}`$ are simply the projections of the $`\nu _e`$ state onto the mass eigenstates: $`P_1=\mathrm{cos}^2\theta `$, $`P_{2e}=\mathrm{sin}^2\theta `$. In this case we are left with the standard vacuum oscillation formula:
$$P_{ee}^{\mathrm{vac}}=1\mathrm{sin}^22\theta \mathrm{sin}^2(\mathrm{\Delta }m^2(Lr)/4E)$$
(12)
which describes the oscillations on the way from the surface of the Sun to the surface of the Earth. The probability is symmetric under $`\theta \frac{\pi }{2}\theta `$.
(2) For $`\mathrm{\Delta }m^2/E10^{14}`$ eV, the last term in Eq. (9) vanishes and we recover the incoherent MSW survival probability. For $`\mathrm{\Delta }m^2/E10^{14}10^{12}`$ eV<sup>2</sup>, this term is zero because $`\nu _e`$ adiabatically converts to $`\nu _2`$ and $`P_1=0`$. For $`\mathrm{\Delta }m^2/E\text{ }>10^{12}`$ eV<sup>2</sup>, both $`P_1`$ and $`P_2`$ are nonzero and the term vanishes due to averaging of $`\mathrm{cos}\xi `$.
(3) In the intermediate range, $`5\times 10^{17}\mathrm{\Delta }m^2/E10^{14}`$ eV, adiabaticity is violated and the $`\mathrm{cos}\xi `$ coherent term should be taken into account. The result is similar to vacuum oscillations but with small matter corrections. We define this case as quasi-vacuum oscillations .
The results presented in the following sections have been obtained using the general expression for the survival probability in Eq. (9) with $`P_1`$ and $`P_{2e}`$ found by numerically solving the evolution equation in the Sun and the Earth matter. For $`P_i`$ we use the electron number density of BP2000 model . For $`P_{ie}`$ we integrate numerically the evolution equation in the Earth matter using the Earth density profile given in the Preliminary Reference Earth Model (PREM) .
### B The mixing angle in matter
While, as explained above, our results are obtained by a numerical calculation, it is useful to find approximate analytical expressions for the neutrino survival probability and for various observables. This is done in this and in the following subsections. The analytical expressions help us to qualitatively understand the behaviour of the different observables, particularly for the case of near–maximal mixing.
The probabilities and observables depend on the mixing angle in matter $`\theta _m`$ via $`\mathrm{cos}2\theta _m`$ which enters the probability of the adiabatic conversion and via $`\mathrm{sin}^22\theta _m`$ which determines, e.g., the depth of oscillations in a uniform medium.
We can write $`\mathrm{cos}2\theta _m`$ in terms of the neutrino oscillation parameters and the electron density in medium as:
$$\mathrm{cos}2\theta _m=\frac{1+\eta \mathrm{cos}2\theta }{(12\eta \mathrm{cos}2\theta +\eta ^2)^{1/2}}.$$
(13)
Here
$$\eta \frac{l_0}{l_\nu }=0.66\left(\frac{\mathrm{\Delta }m^2/E}{10^{15}\text{eV}}\right)\left(\frac{10^2\text{g cm}^3}{\rho Y_e}\right)$$
(14)
is the ratio between the refraction length, $`l_0`$, and the neutrino oscillation length in vacuum, $`l_\nu `$:
$`l_0{\displaystyle \frac{2\pi m_N}{\sqrt{2}G_F\rho Y_e}},`$ $`l_\nu {\displaystyle \frac{4\pi E}{\mathrm{\Delta }m^2}}.`$ (15)
Here $`\rho `$ is the matter density and $`Y_e`$ is the number of electrons per nucleon.
Around maximal mixing $`\mathrm{cos}2\theta _m`$ can be expanded as
$$\mathrm{cos}2\theta _m=\frac{1}{\sqrt{1+\eta ^2}}\left(1\frac{\eta ^3}{1+\eta ^2}ϵ\right).$$
(16)
In the limit of weak matter effects, $`\eta 1`$, and in the matter dominance case, $`\eta 1`$, we get:
$$\mathrm{cos}2\theta _m=\{\begin{array}{cc}ϵ\hfill & \eta 1\hfill \\ 1\hfill & \eta 1\hfill \end{array}.$$
(17)
The dependence of $`\mathrm{cos}2\theta _m`$ on $`ϵ`$ is smooth. It is stronger for $`\eta 1`$ and highly suppressed for $`\eta 1`$. For precisely maximal mixing we have $`\mathrm{cos}2\theta _m=1/\sqrt{1+\eta ^2}`$, which decreases from zero in vacuum to $`1`$ in the matter dominance case.
The expression for $`\mathrm{sin}^22\theta _m`$ for near maximal mixing is given by
$$\mathrm{sin}^22\theta _m=\frac{\eta ^2}{1+\eta ^2}\left(1+\frac{2\eta }{1+\eta ^2}ϵ\right),$$
(18)
which leads to
$$\mathrm{sin}^22\theta _m=\{\begin{array}{cc}1+\frac{2}{\eta }ϵ\hfill & \eta 1\hfill \\ \eta ^2(1+2\eta ϵ)\hfill & \eta 1\hfill \end{array}.$$
(19)
In both cases the $`ϵ`$ corrections to the value at maximal mixing are strongly suppressed. In vacuum, the correction is quadratic in $`ϵ`$: $`\mathrm{sin}^22\theta =1ϵ^2`$.
### C Survival probability
We first consider the survival probability for electron neutrinos without the regeneration effect in the Earth. It describes the $`\nu _e`$ flux arriving at the Earth during the day. In daytime, $`P_{2e}=\mathrm{sin}^2\theta `$. Consequently, Eq. (9) gives, in the region where the oscillating term in Eq. (9) is absent,
$$P_D=\frac{1}{2}+ϵ\left(P_1\frac{1}{2}\right).$$
(20)
The neutrino evolution in the Sun described by the probability $`P_1`$ can be approximated by the well known expression
$$P_1=\frac{1}{2}+\left(\frac{1}{2}P_c\right)\mathrm{cos}2\theta _S.$$
(21)
Here $`\theta _S`$ is the matter mixing angle at the production point:
$$\mathrm{cos}2\theta _S\mathrm{cos}2\theta _m(\eta _S),\eta _S\eta (\rho _SY_{eS}),$$
(22)
where $`\rho _S`$ and $`Y_{eS}`$ are, respectively, the matter density and the number of electrons per nucleon at the production point. Eq. (21) is assumed to be averaged over the production region in the Sun. $`P_c`$ is the jump probability which describes the violation of adiabaticity. For an exponential density profile it takes the following form :
$$P_c=\frac{e^{\gamma \mathrm{sin}^2\theta }e^\gamma }{1e^\gamma }=\frac{e^{(\gamma /2)(1ϵ)}e^\gamma }{1e^\gamma },$$
(23)
where $`\gamma `$ is the ratio of the density scale height $`l_\rho `$ and the neutrino oscillation length:
$`\gamma {\displaystyle \frac{4\pi ^2l_\rho }{l_\nu }}=1.05\left({\displaystyle \frac{\mathrm{\Delta }m^2/E}{10^{15}\text{eV}}}\right)\left({\displaystyle \frac{l_\rho }{r_0}}\right),`$ $`l_\rho {\displaystyle \frac{\rho }{d\rho /dr}}.`$ (24)
The length scale $`r_0=R_{}/10.54`$ is related to the exponential approximation to the solar density profile, $`\rho =\rho _0\mathrm{exp}(r/r_0)`$.
Notice that, originally, Eq. (23) was derived for a mixing angle $`\theta <\frac{\pi }{4}`$ where resonant enhancement is possible. However both Eq. (13) and Eq. (23) can be analytically continued into the second octant, $`\theta >\pi /4`$, and used to compute the corresponding survival probability for $`ϵ<0`$ .
Inserting $`P_1`$ of Eq. (21) into Eq. (20), we get
$$P_D=\frac{1}{2}+ϵ\left(\frac{1}{2}P_c\right)\mathrm{cos}2\theta _S.$$
(25)
Let us study the properties of $`P_D`$. In Fig. 1 we plot $`P_D`$ as a function of $`\mathrm{\Delta }m^2/4E`$ for different values of the deviation from maximal mixing. We show in the figure the results obtained by the numerical calculation as well as from the corresponding analytical approximation (23) for exponential density profile. We learn the following points from Eq. (25) and Fig. 1:
1) For maximal mixing ($`ϵ=0`$), $`P_D=1/2`$ independently of the adiabaticity violation (encoded in $`P_c`$), matter effects, etc..
2) For near-maximal mixing ($`ϵ0`$), solar matter effects lead to an energy dependent probability in the MSW region $`4\times 10^{16}\mathrm{\Delta }m^2/E10^{10}`$ eV.
3) For $`\mathrm{\Delta }m^2/E10^{15}`$ eV, the adiabaticity condition, $`\gamma 1`$, is fulfilled (see Eq. (24)). Consequently $`P_c=0`$ and Eq. (25) gives
$$P_D=\frac{1}{2}\left(1+ϵ\mathrm{cos}2\theta _S\right)(\mathrm{\Delta }m^2/E10^{15}\text{eV}).$$
(26)
From Fig. 1 we see that the analytical expression gives a good description of the propagation in the Sun for this region.
4) For $`\eta _S1`$ (weak matter effect), we have $`\mathrm{cos}2\theta _S=ϵ`$ (see Eq. (17)) and the probability (26) reduces to the vacuum oscillation probability:
$$P_D=\frac{1}{2}\left(1+ϵ^2\right)(\mathrm{\Delta }m^2/E10^{11}\text{eV}).$$
(27)
5) For $`\eta _S1`$ (the resonance layer is far enough from the production point which is close to the center of the sun), we have $`\mathrm{cos}2\theta _S=1`$ and
$$P_D=\frac{1}{2}\left(1ϵ\right)(10^{15}\text{eV}\mathrm{\Delta }m^2/E10^{11}\text{eV}).$$
(28)
This is the region of strong adiabatic conversion, when $`\nu _e`$ produced in the center of the Sun practically coincides with the matter eigenstate $`\nu _{2m}`$ at the production point and during its propagation inside the Sun. Therefore, it emerges from the Sun and reaches Earth as $`\nu _2`$.
6) For $`\mathrm{\Delta }m^2/E<10^{15}`$ eV, adiabaticity is violated. We see from Fig. 1 that the analytical results obtained for an exponential density profile differ from the results of the numerical calculations. In particular, the analytical result shows a “slower” transition to the vacuum oscillation regime or, in other words, it overestimates the size of the matter effects in the quasi–vacuum oscillation region. The same value of the survival probability appears for about two times smaller $`\mathrm{\Delta }m^2/E`$. Similar conclusions have been drawn in refs. .
7) For $`ϵ0`$, the effects of the adiabatic edge situated at $`\mathrm{\Delta }m^2/E=(10^{12}10^{10})`$ eV, and of the non-adiabatic edge situated at $`\mathrm{\Delta }m^2/E=(10^{16}3\times 10^{15})`$ eV, become important.
8) For $`\mathrm{\Delta }m^2/E10^{17}`$ eV, as we noted in Sec. II A, the effect is reduced to the vacuum oscillation between the surfaces of the Sun and the Earth and the average survival probability, $`P_D=1\frac{1}{2}\mathrm{sin}^22\theta `$, (shown in Fig. 1) coincides with that in Eq. (27). However, in this region averaging of oscillations does not occur and for the survival probability we should use
$$P_{ee}1(1ϵ^2)\mathrm{sin}^2\frac{\varphi }{2},$$
(29)
where $`\varphi =\mathrm{\Delta }m^2(Lr)/4E`$ is the oscillation phase which does not depend on $`ϵ`$. In principle matter effects strongly suppress the oscillations inside the Sun and the Earth. However, the modification of the observables is small, since the size of the Sun (and the Earth) is much smaller than the oscillation length in vacuum. The quadratic dependence of the probability on $`ϵ`$ is again restored. Moreover, time variations of signals and the distortion of the energy spectrum originate from the dependence of the phase on $`L`$ and $`E`$. Therefore according to (29), both the variations and the distortion are proportional to $`(1ϵ^2)`$. That is, the dependence of all observables on $`ϵ^2`$ near maximal mixing is very weak.
According to Eq. (28) and Fig. 1, in the MSW region (more precisely at the bottom of the suppression pit), the survival probability depends on $`ϵ`$ linearly. For $`\mathrm{\Delta }m^2/E`$ below and above the MSW region the probability converges to the vacuum oscillation probability. The deviation of the latter from the probability at maximal mixing depends on $`ϵ`$ quadratically and the dependence on the sign of $`ϵ`$ disappears. From this we infer that the sensitivity of $`P_D`$ and consequently of observables to $`ϵ`$ is much higher in the MSW region.
Notice that there are two transition regions (between MSW and the vacuum oscillation regions) where the effects are mainly due to vacuum oscillations with small matter corrections. We call these regions ‘Quasi-Vacuum Oscillation regions’: we denote by QVO<sub>L</sub> (QVO<sub>S</sub>) the one with $`\mathrm{\Delta }m^2`$ larger (smaller) than in the MSW region. In the QVO<sub>L</sub> and QVO<sub>S</sub> regions, the linear dependence of the probability and the observables transforms into quadratic dependence.
### D Regeneration effects in the Earth
For a neutrino arriving at night time, Earth matter effects should be taken into account. To gain a qualitative understanding of the Earth effects, we make the crude approximation of uniform density, such that the mixing angle in the Earth is constant, $`\theta _E`$, along the neutrino trajectory. In this case, the neutrino propagation has an oscillatory character. We assume that the oscillations are averaged out, and therefore
$`P_{2e}(\text{vacuum})`$ $`=`$ $`\mathrm{sin}^2\theta ,`$ (30)
$`P_{2e}(\text{uniform density})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\mathrm{sin}^2\theta +\mathrm{sin}^2(2\theta _E\theta )\right].`$ (31)
It is convenient to introduce a regeneration factor which describes the Earth matter effect:
$$f_{\text{reg}}P_{2e}(\text{matter})P_{2e}(\text{vacuum})=\frac{1}{2}\mathrm{sin}2\theta _E\mathrm{sin}(2\theta _E2\theta ).$$
(32)
Denoting by $`\eta _E`$ the parameter $`\eta `$ in the Earth:
$$\eta _E\eta (\rho _EY_{eE}),$$
(33)
where $`\eta `$ is defined in Eq. (14), we find from Eq. (32):
$$f_{\text{reg}}=\frac{\eta _E\mathrm{sin}^22\theta }{2(12\eta _E\mathrm{cos}2\theta +\eta _E^2)}.$$
(34)
Note that $`f_{\text{reg}}`$ is always (for any value $`\mathrm{\Delta }m^2/E`$ and $`ϵ`$) positive, i.e. the matter effect of Earth always enhances the survival probability $`P_{ee}`$.
The parameter $`f_{\text{reg}}`$ can be expanded around $`ϵ=0`$:
$$f_{\text{reg}}=\frac{\eta _E}{2(1+\eta _E^2)}\left(1+\frac{2\eta _E}{1+\eta _E^2}ϵ\right).$$
(35)
The regeneration factor has a maximum at $`\eta _E1`$ which corresponds to $`\mathrm{\Delta }m^2/E(23)\times 10^{13}`$ eV. This determines the strong regeneration region in the $`\mathrm{\Delta }m^2/E`$ scale. It is situated in the middle of the the solar MSW region. Strong regeneration effects are already excluded by the SuperKamiokande result on Day-Night asymmetry . Consequently, the strong regeneration region separates two (allowed) parts within the MSW region in which the regeneration effects are small:
1. HIGH $`\mathrm{\Delta }m^2`$ region, where $`\eta _E1`$ and $`f_{\text{reg}}1/\eta _E`$;
2. LOW $`\mathrm{\Delta }m^2`$ region, where $`\eta _E1`$ and $`f_{\text{reg}}\eta _E/2`$.
We will quantify the borders of these regions in Sec. III. In both cases the regeneration effect is suppressed by a small parameter and it disappears when moving away from the strong regeneration region.
Let us stress that the SuperKamiokande limit on regeneration effects holds for the energy range $`E=515`$ MeV to which this experiment is sensitive. This corresponds to $`\mathrm{\Delta }m^2=(23)\times 10^6`$ eV<sup>2</sup>. However, strong regeneration effects are not excluded for other energies, in particular, for low energy neutrinos. In this case $`\eta _E`$ can be of the order one and $`f_{\text{reg}}`$ at its maximum.
To first order in $`ϵ`$, we obtain from Eq. (35):
$$f_{\text{reg}}=\{\begin{array}{cc}\frac{1}{2\eta _E}\left(1+\frac{2}{\eta _E}ϵ\right)\hfill & (\mathrm{HIGH})\hfill \\ \frac{\eta _E}{2}(1+2\eta _Eϵ)\hfill & (\mathrm{LOW})\hfill \end{array}.$$
(36)
In both the HIGH and LOW regions the dependence of the regeneration factor on $`ϵ`$ is further suppressed by the small parameter min$`\{2/\eta _E,2\eta _E\}`$.
The probability $`P_{ee}`$ at night time, $`P_N`$, can be written (for the region where the oscillating term in Eq. (9) is absent) as follows:
$$P_N=P_1+(12P_1)\left(\mathrm{sin}^2\theta +f_{\text{reg}}\right)=\frac{1}{2}+\frac{1}{2}(12P_c)\mathrm{cos}2\theta _S(\mathrm{cos}2\theta 2f_{\text{reg}}).$$
(37)
The average daily survival probability is given for $`\eta _S1`$ ($`\mathrm{cos}2\theta _S1`$) by
$$\overline{P}\frac{1}{2}(P_D+P_N)=\frac{1}{2}[1+(12P_1)(f_{\text{reg}}\mathrm{cos}2\theta )].$$
(38)
The Day–Night asymmetry is given by
$$A_{\mathrm{N}\mathrm{D}}\frac{P_NP_D}{\overline{P}}=\frac{2f_{\text{reg}}}{1/(12P_1)\mathrm{cos}2\theta +f_{\text{reg}}}.$$
(39)
Let us now consider the dependence of $`A_{ND}`$ and $`\overline{P}`$ on $`ϵ`$ in the HIGH and in the LOW regions keeping the lowest order terms in $`ϵ`$ and in min$`\{\eta _E,1/\eta _E\}`$.
1. In the HIGH region, the adiabaticity condition is satisfied and we can safely put $`P_c=0`$. Consequently, $`P_1=\mathrm{cos}^2\theta _S`$. Then, Eq. (38) simplifies:
$$\overline{P}=\frac{1}{2}[1\mathrm{cos}2\theta _S(f_{\text{reg}}\mathrm{cos}2\theta )].$$
(40)
With near-maximal mixing, we obtain:
$$\overline{P}=\frac{1}{2}\left(1\frac{ϵ1/2\eta _E}{\sqrt{1+\eta _S^2}}\right)+𝒪(\frac{1}{\eta _E\eta _S},\frac{ϵ}{\eta _E}).$$
(41)
For $`\eta _S^21`$ (small $`\mathrm{\Delta }m^2`$) we find:
$$\overline{P}\frac{1}{2}\left(1ϵ+\frac{1}{2\eta _E}+\frac{ϵ}{\eta _E^2}\right),$$
(42)
where the last term is a small correction which comes from the regeneration factor. For the Day–Night asymmetry we obtain:
$$A_{\mathrm{N}\mathrm{D}}=\frac{1}{\eta _E\sqrt{1+\eta _S^2}}\left[1+\frac{(\sqrt{1+\eta _S^2}\eta _S^3)ϵ}{1+\eta _S^2}+𝒪(\frac{1}{\eta _E^2},\frac{ϵ}{\eta _E})\right].$$
(43)
Notice that the $`ϵ`$-dependent effect changes sign between large and small $`\eta _S`$. For $`\eta _S^21`$ we get:
$$A_{\mathrm{N}\mathrm{D}}\frac{1}{\eta _E}(1+ϵ).$$
(44)
The asymmetry increases linearly with $`ϵ`$ for small enough $`\eta _S^2`$.
2. In the LOW region we can take $`\eta _S=0`$ and, consequently, $`P_1=P_c`$. Then, Eq. (38) simplifies:
$$\overline{P}=\frac{1}{2}\left[1+(12P_c)(f_{\text{reg}}ϵ)\right].$$
(45)
In LOW region we have $`\gamma 1`$ (see Eq. (24)) and, therefore, $`P_c1`$ (see Eq. (23)). In any case, the $`P_c`$-dependent term in Eq. (45) is suppressed by a small factor and can be neglected. We get:
$$\overline{P}=\frac{1}{2}\left(1ϵ+\frac{\eta _E}{2}\right).$$
(46)
For the Day-Night asymmetry we obtain taking $`P_c1`$:
$$A_{\mathrm{N}\mathrm{D}}\frac{2f_{\text{reg}}}{1+2P_cϵ+f_{\text{reg}}},$$
(47)
and in the limit $`\eta _E1`$:
$$A_{\mathrm{N}\mathrm{D}}=\eta _E\left(12e^{\gamma /2}\frac{\eta _E}{2}+ϵ\right).$$
(48)
A few comments are in order. The Day–Night asymmetry is suppressed by $`\eta _E`$. The $`ϵ`$-dependent effect is one of several small corrections to the leading result. For large parts of the LOW region it is the leading correction, but for large $`\mathrm{\Delta }m^2/E`$, the subleading regeneration effect could be comparable while for small $`\mathrm{\Delta }m^2/E`$, the non-adiabatic correction could give the main correction.
As we mentioned before, a strong regeneration effect with $`\eta 1`$ and $`f_{\text{reg}}1`$ is not excluded for low energy neutrinos. In particular in the LOW region strong regeneration can show up for the beryllium and pp-neutrino components of the spectrum. In this case, the approximation $`\eta _E1`$ does not work and one should use the complete expression for the regeneration factor and the asymmetry.
As seen from Eqs. (42), (46), (44) and (48), the dominant dependence on $`ϵ`$ of both $`\overline{P}`$ and $`A_{\mathrm{N}\mathrm{D}}`$ arises from the dependence of the solar survival probability on $`ϵ`$. The dependence which follows from the regeneration factor is further supressed by min$`\{\eta _E,1/\eta _E\}`$. (For $`A_{\mathrm{N}\mathrm{D}}`$ the $`ϵ`$ dependence follows from the dependence on $`\overline{P}`$ in the denominator.)
From Eqs. (44) and (48) we find that the Day-Night asymmetry strongly depends on $`\mathrm{\Delta }m^2`$ due to the $`\eta _E`$ or $`1/\eta _E`$ factors. In the HIGH region $`A_{\mathrm{N}\mathrm{D}}\mathrm{\Delta }m^2`$ while in LOW region $`A_{\mathrm{N}\mathrm{D}}1/\mathrm{\Delta }m^2`$. The dependence of the asymmetry on $`ϵ`$ is much weaker: In both regions $`A_{\mathrm{N}\mathrm{D}}(1+ϵ)`$. Consequently, the measurements of $`A_{\mathrm{N}\mathrm{D}}`$ are very sensitive to $`\mathrm{\Delta }m^2`$ while the sensitivity to $`ϵ`$ is substantially lower.
Let us comment on the range of validity of the approximate treatment of the Earth effects . In the HIGH region, for $`E`$ MeV and $`\mathrm{\Delta }m^210^6\text{eV}^2`$, the oscillation length $`l_\nu `$ is much shorter than the size of the Earth and neutrinos undergo many oscillations inside the Earth. The constant electron number density approximation gives a good description of $`A_{\mathrm{N}\mathrm{D}}`$ which involves integration over the zenith angle. In the LOW region, for $`\mathrm{\Delta }m^210^6\text{eV}^2`$, the oscillation length is approximately equal to the refraction length $`l_0`$ and the latter is comparable to the size of the Earth (independently of $`\mathrm{\Delta }m^2`$). The details of the $`N_e`$ profile do not play an important role, and the effect is determined by the average density along the neutrino trajectory.
Regeneration effect leads also to seasonal variations of signals . These variations, however, are less sensitive to the oscillation parameters.
### E Distortion of the energy spectrum
The distortion of the solar neutrino energy spectrum can be characterized by the distortion parameter defined as:
$$s_\nu \frac{E}{\overline{P}}\frac{d\overline{P}}{dE}.$$
(49)
Averaged over the appropriate energy interval this parameter is proportional to the shift of the first moment of the spectrum or to the slope parameter used in the literature.
To understand the distortion of the spectrum (energy dependence of the averaged probability), we remind the reader that $`\eta 1/E`$ and use Eq. (40) for the HIGH and QVO<sub>L</sub> regions and Eq. (45) for the LOW region.
1. The HIGH and QVO<sub>L</sub> regions:
(i) For large $`\mathrm{\Delta }m^2`$, the effect of the adiabatic edge of the suppression pit which is encoded in the $`\eta _S`$ dependence is important:
$$s_\nu ϵ\frac{\eta _S^2}{(1+\eta _S^2)^{3/2}}.$$
(50)
The distortion parameter is proportional to $`ϵ`$. The slope (shift of the first moment) is positive (negative) for $`ϵ<0`$ ($`ϵ>0`$). The distortion decreases rapidly with $`\eta _S`$.
(ii) For small $`\mathrm{\Delta }m^2`$, the earth regeneration which is related to the $`\eta _E`$ dependence is important:
$$s_\nu \frac{1}{2\eta _E}\left(1+ϵ\frac{1}{2\eta _E}+\frac{4ϵ}{\eta _E}\right).$$
(51)
The distortion decreases with increase of $`\eta _E`$ or $`\mathrm{\Delta }m^2`$. The sensitivity to $`ϵ`$ is much weaker than in the previous case and it follows mainly from the dependence on the average probability in the denominator.
2. The LOW region:
(i) For large $`\mathrm{\Delta }m^2`$, the regeneration effect is important:
$$s_\nu \frac{\eta _E}{2}(1+ϵ).$$
(52)
The slope increases with $`\eta _E`$. The dependence on $`ϵ`$ is weak.
(ii) For small $`\mathrm{\Delta }m^2`$, the effect of the nonadiabatic edge of the solar suppression pit gives the dominant effect:
$$s_\nu =2ϵE\frac{dP_c}{dE}ϵ\gamma e^{\gamma /2}.$$
(53)
The distortion is proportional to $`ϵ`$. The slope is positive (negative) for $`ϵ>0`$ ($`ϵ<0`$). The effect is suppressed for relatively weak violation of the adiabaticity.
As we have mentioned previously, for non-averaged vacuum oscillations $`s_\nu (1ϵ^2)`$, that is, the dependence of distortion on $`ϵ^2`$ is very weak.
### F Summary
Let us summarize the results of our analytical studies. We have found simple analytical expressions for various observables (rates of events given by the survival probability, Day-Night asymmetry, distortion of the spectrum etc.) in terms of $`\mathrm{\Delta }m^2`$ and the deviation from maximal mixing $`ϵ`$. These approximate analytical expressions reproduce correctly the functional dependence of the observables on $`\mathrm{\Delta }m^2`$ and $`ϵ`$ and allow us to understand all the features of the exact numerical calculations.
The effects and the dependence of observables on $`ϵ`$ change with $`\mathrm{\Delta }m^2`$. Accordingly, we define several regions of $`\mathrm{\Delta }m^2`$ which have different physical pictures. As $`\mathrm{\Delta }m^2`$ decreases from its upper (CHOOZ) bound, we have the following regions (we quantify borders of these regions in the next section):
* quasi-vacuum oscillation region with large $`\mathrm{\Delta }m^2`$ (QVO<sub>L</sub>),
* MSW region with high $`\mathrm{\Delta }m^2`$ (HIGH),
* MSW region with low $`\mathrm{\Delta }m^2`$ (LOW),
* quasi-vacuum oscillation region with small $`\mathrm{\Delta }m^2`$ (QVO<sub>S</sub>),
* region of non-averaged vacuum oscillations (VO). The high $`\mathrm{\Delta }m^2`$ part of the VO region will be called VAC<sub>L</sub>.
As concerns the $`ϵ`$ dependence of observables in these regions, we find two main conclusions:
1. Maximal mixing, $`ϵ=0`$, is not a special point as far as the phenomenology is concerned (in contrast with theory). No divergencies or discontinuities appear in the dependence of observables on $`ϵ`$. The dependence on $`ϵ`$ is smooth and, in many cases, very weak. To mention a few examples: The day time survival probability is 1/2 at $`ϵ=0`$. Earth regeneration effects, however, enhance the survival probability. In the MSW region the slope of the energy spectrum distortion is proportional to $`ϵ`$ in certain regions of $`\mathrm{\Delta }m^2`$, and consequently changes sign at $`ϵ=0`$. The Day-Night asymmetry is proportional to $`(1+ϵ)`$.
2. The character of the $`ϵ`$-dependence of observables is different in the vacuum oscillations and MSW regions. In the regions of vacuum oscillation all the effects depend on $`ϵ`$ quadratically. More precisely, the $`ϵ`$-dependent factors are of two types:
$`1+ϵ^2`$ $`(\mathrm{averaged}\mathrm{oscillations}),`$ (54)
$`1ϵ^2`$ $`(\mathrm{non}\mathrm{averaged}\mathrm{oscillations}).`$ (55)
The dependence is symmetric with respect to interchanging $`ϵϵ`$. In the MSW regions (both HIGH and LOW), observables depend on $`ϵ`$ linearly. Obviously the dependence is non-symmetric with respect to $`ϵϵ`$. In the region of strong adiabatic conversion (bottom of the pit) we get for the survival probaility, the Day-Night asymmetry and the distortion parameter:
$$P_D(1ϵ),A_{\mathrm{N}\mathrm{D}},s_\nu (1+ϵ).$$
(56)
At the edges of the MSW region (edges of the suppression pit) we find
$$s_\nu ϵ.$$
(57)
The quasi-vacuum oscillation regions, QVO<sub>L</sub> and QVO<sub>S</sub>, are transition regions between the MSW and vacuum oscillation regions, where the linear dependence of observables transform into quadratic dependence.
Thus, for $`ϵ1`$ the sensitivity of experiments to deviation from maximal mixing is much higher in the MSW regions. It will be difficult to measure $`ϵ`$ near maximal mixing if $`\mathrm{\Delta }m^2`$ turns out to be in the VO or QVO regions.
## III Present Status of Maximal Mixing
In this section we describe the present status of maximal as well as near–maximal mixing. We use for this purpose the latest available results on solar neutrinos from Homestake , SAGE , GALLEX+GNO , and the 1117 days of data sample of Super–Kamiokande experiments.
We calculate the acceptability of maximal and near–maximal mixing as a function of $`\mathrm{\Delta }m^2`$ in the whole allowed range, that is, below the upper bound $`\mathrm{\Delta }m^210^3`$ eV<sup>2</sup> from the reactor experiments . The goal of this study is to find excluded regions of the oscillation parameters $`\mathrm{\Delta }m^2ϵ`$, as well as the regions of these parameters which are allowed and most plausible. We quantify our statements in terms of confidence level at which a given region is accepted (probability of realization) or excluded. We study the dependence of our conclusions on uncertainties in the solar neutrino fluxes (the SSM uncertainties) as well as on the procedure employed in the analysis. Some model- and the procedure-independent statements are formulated. The extent to which the results of this Section hold in a three generation framework with $`|U_{e3}|0`$ is discussed in Sec. VI.
### A Global fit: allowed and forbidden regions
The results of a global fit to all existing experimental data on solar neutrinos are shown in Figs. 2$``$5. The analysis includes rates in Chlorine , Gallium , and Super–Kamiokande experiments, as well as the zenith angle dependence and the shape of the recoil electron spectrum in Super–Kamiokande.
In Fig. 2 we plot the contours of constant confidence level (iso–contours) in the $`\mathrm{\Delta }m^2ϵ`$ plane. Points inside a given contour are accepted at a lower confidence level than on the contour itself. In the “global” analysis we combine the information on the Day–Night variation of the event rates and the recoil energy spectrum at Super–Kamiokande by using their independently measured spectra during the day and during the night. With this the total number of independent experimental inputs in the global analysis is 38 which includes 3 rates, and 35 data points for the Super–Kamiokande day and night recoil energy spectra ($`2\times 18`$ bins minus 1 overall normalization). We do not include in the analysis the new lower energy bin as its systematic uncertainty is still under study by the Super–Kamiokande Collaboration . We use the solar neutrino fluxes from the Standard Solar Model (SSM) of Ref. (BP98). The contours have been defined by the shift in $`\chi ^2`$, $`\mathrm{\Delta }\chi ^2`$, with respect to the global minimum in the whole plane of the oscillation parameters. The minimum lies in the LMA solution region:
$$\chi ^2=33.4\text{for}36\text{d.o.f}.$$
(58)
which corresponds to a probability of 59%. For details on the statistical analysis we refer to Ref. .
In Fig. 3 we show the dependence of $`\mathrm{\Delta }\chi ^2`$ on $`\mathrm{\Delta }m^2`$ for three different values of $`ϵ`$ ($`0.3,0,+0.3`$). This figure corresponds to three $`\mathrm{\Delta }m^2`$-profiles (cuts) of the confidence level from Fig. 2.
What is the impact of individual experimental results on the global fit? Panels (a)$``$(f) of Fig. 4 show the contours of constant $`\chi _i^2`$ for each individual observable $`i`$. As mentioned above the total $`\chi ^2`$ in panel (f) is obtained by combining the $`\chi ^2`$ of the individual rates (including the correlation of their theoretical errors) with the corresponding $`\chi ^2`$ for the Super–Kamiokande night and day recoil energy spectra.
Our calculations allow us to define regions of oscillation parameters that are excluded and other that are allowed. It follows from Fig. 2 that there are two main regions of $`\mathrm{\Delta }m^2`$ which are excluded at a very high (more than $`99.99\%`$) confidence level:
1. The regeneration region: for maximal mixing we find
$$\mathrm{\Delta }m^2=(0.68)\times 10^6\text{eV}^2(ϵ=0).$$
(59)
The excluded region increases with $`ϵ`$ for positive values of $`ϵ`$. In particular, for $`ϵ=0.3`$, the excluded range becomes $`\mathrm{\Delta }m^2(0.410)\times 10^6\text{eV}^2`$. The excluded region also expands for negative values of $`ϵ`$ at $`ϵ\text{ }<0.25`$.
In the regeneration region the solar neutrino observables are strongly modified by the $`\nu _e`$-regeneration in Earth. As we have discussed in Sec. II, the regeneration is always positive, thus leading to an increase in the $`\nu _e`$-flux. Correspondingly, the counting rates in all experiments increase. There are two main physical effects in the regeneration region that are inconsistent with observations and therefore lead to the exclusion:
(i) A large Day–Night asymmetry, $`A_{\mathrm{N}\mathrm{D}}\text{ }>0.2`$ (see Figs. 4(d) and 6), is in contradiction to the Super–Kamiokande result and plays a dominant role for positive $`ϵ`$.
(ii) A large Ar–production rate, $`Q_{Ar}>4`$ SNU (see Figs. 4(b) and 6), is in contradiction to the Homestake result and leads to the increase of the excluded region for negative $`ϵ`$.
2. The vacuum oscillation region:
$$\mathrm{\Delta }m^2=(0.144)\times 10^{10}\text{eV}^2(ϵ=0).$$
(60)
The size of this region depends very weakly on $`ϵ`$ in the interval $`0.3\text{ }<ϵ\text{ }<+0.3`$. The exclusion follows from the interplay between the total rates and the shape of the recoil electron energy spectrum. Notice that the rates only (see Fig. 5(a)) can be reproduced rather well in some parts of the region in Eq. (60). However, in these regions the distortion of the recoil electron spectrum is in contradiction with the Super–Kamiokande results (see Fig. 4(e)). More specifically, for $`\mathrm{\Delta }m^210^{10}`$ eV<sup>2</sup>, a negative slope (or shift of the first moment) in the “reduced” spectrum is expected. The reduced spectrum is defined as the ratio
$$R(E)N_i(E)/N_i(E)^{\mathrm{SSM}},$$
(61)
where $`N_i(E)`$ ($`N_i(E)^{\mathrm{SSM}}`$) is the number of events in a given electron energy, $`E`$, bin $`i`$ with (without) oscillations. In the same way we define any “reduced” observable as the ratio of its value with respect to the expected one in the SSM in the absence of oscillations.
3. There is another region forbidden at 99.99 % CL extending from $`\mathrm{\Delta }m^2=(0.32)\times 10^8\text{eV}^2`$ for $`ϵ=0.3`$ to $`\mathrm{\Delta }m^2=6\times 10^9\text{eV}^2`$ and $`ϵ=0.22`$. As seen in Fig. 3 this region is forbbiden with a lower CL. The exclusion follows mainly from the effect of the rates as seen in Fig. 5(a) being mainly driven by the bad fit to the Gallium rate (see Fig. 7(a)).
We distinguish five regions of the oscillation parameters where maximal mixing is allowed at a confidence level that is lower than 99.9%:
1) Quasi vacuum oscillation region with large $`\mathrm{\Delta }m^2`$ (QVO<sub>L</sub>):
$$\mathrm{\Delta }m^2\left(3\frac{E}{10\mathrm{MeV}}8\right)\times 10^4\text{eV}^2,$$
(62)
where $`E`$ is the average detected energy for a given experiment. The upper bound comes from reactor experiments . Here the flavour conversion is mainly due to averaged vacuum oscillations with only small matter corrections inside the Sun and the Earth.
2) MSW region with high $`\mathrm{\Delta }m^2`$ (HIGH):
$$\mathrm{\Delta }m^2\left(0.13\frac{E}{10\mathrm{MeV}}\right)\times 10^4\text{eV}^2.$$
(63)
This region corresponds to the maximal and near–maximal mixing part of the LMA solution. It is restricted from below by strong Earth regeneration effects (large Day–Night asymmetry and large Ar–production rate). Maximal mixing is acceptable at confidence level larger than $`99.1\%`$. As follows from Fig. 3, the dependence of $`\mathrm{\Delta }\chi ^2`$ on $`\mathrm{\Delta }m^2`$ is rather weak. The global fit becomes substantially better with increase of $`ϵ`$ (shift to positive values): $`ϵ\text{ }>0.25`$ is accepted at 90% CL. For $`ϵ=0.3`$, the 90% CL allowed region expands to $`\mathrm{\Delta }m^2=(210)\times 10^5`$ eV<sup>2</sup>. In contrast, the goodness of the fit decreases when we shift to negative values of $`ϵ`$.
3) MSW region with low $`\mathrm{\Delta }m^2`$ (LOW):
$$\mathrm{\Delta }m^2(0.13)\times 10^7\text{eV}^2.$$
(64)
Here maximal mixing is acceptable at $`99\%`$ CL in the interval
$$\mathrm{\Delta }m^2(0.12)\times 10^7\text{eV}^2(\text{CL}99\%).$$
(65)
For positive $`ϵ`$ the fit improves while for negative $`ϵ`$ it worsens. In particular, $`ϵ=0.2`$ is accepted at $`90\%`$ CL for $`\mathrm{\Delta }m^2(0.81.5)\times 10^7`$ eV<sup>2</sup>. The local minimum occurs at $`\mathrm{\Delta }m^2=1.0\times 10^7\text{eV}^2`$ and $`ϵ=0.21`$. With increase of $`ϵ`$ the accepted region of $`\mathrm{\Delta }m^2`$ shifts to larger values. At $`ϵ=0.3`$ we obtain $`\mathrm{\Delta }m^2=(0.62)\times 10^7\text{eV}^2`$. Conversely negative values of $`ϵ`$ are disfavored.
4) Quasi vacuum oscillation region with small $`\mathrm{\Delta }m^2`$ (QVO<sub>S</sub>):
$$\mathrm{\Delta }m^2(0.11)\times 10^8\text{eV}^2.$$
(66)
In this region the flavour conversion is due to (mainly non-averaged) vacuum oscillations with small matter effects. Maximal mixing is acceptable at $`99\%`$ CL in the interval
$$\mathrm{\Delta }m^2(0.51)\times 10^8\text{eV}^2(\text{CL}99\%).$$
(67)
The fit becomes worse with increase of $`|ϵ|`$, but while for $`ϵ>0`$ the QVO<sub>S</sub> region is essentially excluded, for $`ϵ<0`$ we still have a reasonably good fit.
5) Vacuum oscillation region with relatively large $`\mathrm{\Delta }m^2`$ (VAC<sub>L</sub>):
$$\mathrm{\Delta }m^2(0.41)\times 10^9\text{eV}^2.$$
(68)
Maximal mixing is accepted at a confidence level better than 99% only in a very small interval centered at:
$$\mathrm{\Delta }m^26.6\times 10^{10}\text{eV}^2(\text{CL}99\%).$$
(69)
In this interval, the goodness of the fit depends on $`ϵ`$ very weakly.
Summarizing, maximal (or near–maximal) mixing is allowed at $`99\%`$ or slightly lower CL in several small intervals of $`\mathrm{\Delta }m^2`$ in the QVO<sub>L</sub>, HIGH, LOW, QVO<sub>S</sub> and VAC<sub>L</sub> solution domains. The values $`ϵ=0.05,0.1,0.2`$ are allowed at 99, 95 and 90% CL, respectively. At $`4\sigma `$, practically the whole HIGH, LOW, QVO and VAC<sub>L</sub> ranges are allowed.
Let us point out the role of individual experimental results in constraining maximal mixing (see Fig. 4). The rates in the Gallium and Super–Kamiokande experiments can be well accounted for at maximal (or near–maximal) mixing, although the Super–Kamiokande measurement slightly disfavors a negative $`ϵ`$. The zenith angle distribution measured by Super–Kamiokande gives some preference to the HIGH region and excludes the strong regeneration region. In contrast, the Super–Kamiokande result on the recoil electron energy spectrum gives some preference to the VAC<sub>L</sub>, HIGH and LOW regions and excludes the range $`\mathrm{\Delta }m^2=(0.34)\times 10^{10}\text{eV}^2`$. Both the zenith angle distribution and the shape of spectrum have weak dependence on $`ϵ`$. In contrast, total rates are sensitive to $`ϵ`$, especially in the HIGH and LOW regions.
### B Dependence of the results on features of the analysis
Let us study the dependence of the allowed and excluded regions in the $`\mathrm{\Delta }m^2ϵ`$ plane on features of our analysis.
Total rates versus spectrum and zenith angle distribution: The total rates give the most stable, reliable, and statistically significant information. We have carried out a fit to the three rates only. In Fig. 5(a) we show the dependence of the shift of $`\chi ^2`$ for this analysis with respect to the absolute minimum in the whole plane of oscillation parameters. The absolute minimum for the analysis of the three rates, $`\chi ^2=0.76`$ (for one d.o.f) is achieved in the SMA region. Comparing Fig. 5(a) with Fig. 3, we learn that if we exclude the information from the recoil electron energy spectrum and the Day–Night variation of the event rates from our analysis, then the allowed and forbidden regions are substantially modified. In particular, we note in Fig. 5(a) the following three features:
(i) The goodness of the fit at maximal mixing from the analysis of the three rates only is worse in whole MSW region. At 99% CL only small interval in the LOW region is allowed. The fit improves however with increase of $`ϵ`$. This means that it is the data on the spectrum and the zenith angle distribution which favor maximal mixing.
(ii) Allowed regions appear in the VAC solution range. We learn that the data on the spectrum and the zenith angle distribution exclude (otherwise) allowed VAC regions.
(iii) The regeneration region is still strongly disfavored by the high Ar–production rate.
The Homestake result: Consider the impact of the Ar–production rate on our results. In Fig. 4(b) we show the fit to only this rate. From the figure we see that the Homestake result strongly disfavors maximal mixing for all $`\mathrm{\Delta }m^2`$ above $`10^{10}`$ eV<sup>2</sup>, that is, in all the globally allowed regions. In Fig. 5(b) we show the result of a global fit to the data without the Homestake result. Clearly, the acceptability of maximal mixing improves for all $`\mathrm{\Delta }m^2\text{ }>10^{10}`$ eV<sup>2</sup> with the best fit points being in the LOW region, $`\mathrm{\Delta }m^2=10^7\text{eV}^2`$. For $`ϵ=0.3`$ the best fit is in VAC<sub>L</sub> region. We would like to emphasize the following points about a fit without the Homestake result:
(i) In the HIGH region, maximal mixing is accepted already at 1.7 $`\sigma `$ with very little dependence on $`ϵ`$.
(ii) In the whole LOW region, maximal mixing gives a very good fit.
(iii) In the whole QVO<sub>L</sub> region, maximal mixing gives a very good fit, but positive values of $`ϵ`$ are still disfavored.
(iv) Regions of strong regeneration and VAC (small $`\mathrm{\Delta }m^2`$) solutions are excluded by the Super–Kamiokande data on the spectrum and the Day–Night asymmetry.
Solar neutrino flux uncertainties: Of all the relevant solar neutrino fluxes, the boron neutrino flux suffers from the largest uncertainty, leading to systematic errors in the predicted detection rate that cannot be estimated reliably at present. One way to avoid this problem is to determine the boron neutrino flux experimentally, using the total rate measured in the Super–Kamiokande experiment. Similar results are obtained if the boron neutrino flux is treated as a free parameter in the analysis (in this case the Super–Kamiokande rate, being the most precise and sensitive to the boron neutrino flux, will fix this flux). In Fig. 5(c) we show the results of a global fit with the boron neutrino flux treated as a free parameter. We plot the dependence of $`\mathrm{\Delta }\chi ^2`$, the $`\chi ^2`$ shift with respect to the absolute minimum, on $`\mathrm{\Delta }m^2`$ for three different values of $`ϵ`$. The shape of the curves is very similar to that in Fig. 3. Also the allowed and excluded regions at a given CL practically coincide with those in Fig. 3. (One should take into account that now we have one additional free parameter and therefore a 99% CL corresponds to $`\mathrm{\Delta }\chi ^2=11.36`$). The reason of this similarity lies on the fact that both the spectrum and the Day–Night variation are flux independent. Furthermore in both cases the boron neutrino flux is fixed by the Super–Kamiokande result.
A comparison of Fig. 3 and Fig. 5(c) shows that the results of our analysis are stable with respect to the way that the uncertainty in the boron neutrino flux is treated.
## IV Maximal Mixing and Predictions for Individual Experiments
In this section we consider the predictions of (near–) maximal mixing for various observables. This will clarify the sensitivity of individual experiments to the neutrino oscillation parameters in the relevant range. The extent to which the results of this Section hold in a three generation framework with $`|U_{e3}|0`$ is discussed in Sec. VI.
### A Total rates
In Fig. 7 we plot the values of the expected event rates: Germanium production rate, Argon production rate, and the rate of events at Super–Kamiokande as functions of $`\mathrm{\Delta }m^2`$ for 3 values of $`ϵ`$: $`+0.3,0,0.3`$. The rates are normalized to the no oscillation expectation, $`R_iQ_i/Q_i^{\mathrm{SSM}}`$. For each value of $`ϵ`$ we plot three curves: the central curves give the expected rates using central values of the BP98 fluxes and the upper and lower lines represent the theoretical uncertainty (without the error for the interaction cross sections) from varying the nine parameters in the SSM within $`\pm 1\sigma `$. The horizontal lines give the experimental values within $`\pm 1\sigma `$ experimental errors. The vertical lines in the range $`\mathrm{\Delta }m^2=(310)\times 10^6`$ eV<sup>2</sup> give the expectation from the SMA 99% CL region (again, including the theoretical uncertainties in each point).
Let us first discuss the dependence of the rates on $`\mathrm{\Delta }m^2`$. The rates are proportional to the survival probability. Therefore, the main features of Fig. 7 reflect the dependence of the survival probability on $`\mathrm{\Delta }m^2/E`$ (see Fig. 1 and Eqs. (27), (28) and (29)). For maximal mixing the survival probability in the Sun as a function of $`\mathrm{\Delta }m^2/E`$ is constant: $`P_D=1/2`$ (see Eq. (20)). The probability $`P_{ee}`$ is enhanced by the Earth regeneration effect for $`\mathrm{\Delta }m^2/E`$ in the range $`(10^{15}10^{11})\text{eV}`$ (see Eq. (38)). For $`ϵ0`$, the effects of the adiabatic edge situated at $`\mathrm{\Delta }m^2/E=(10^{12}10^{10})`$ eV (Eq. (42)), and of the non-adiabatic edge situated at $`\mathrm{\Delta }m^2/E=(10^{16}3\times 10^{15})`$ eV (Eq. (45)), become important. For $`\mathrm{\Delta }m^2/E\text{ }<10^{16}`$ eV, an oscillatory behaviour appears due to non-averaged vacuum oscillations between the Sun and the Earth (Eq. (29)). With certain modifications all these features can be seen in Fig. 7. The simplest dependence is for the Super–Kamiokande rate since only one (boron) neutrino flux contributes (Fig. 7(c)). The Ar–production rate has an additional fine structure due to contributions from additional fluxes, and in particular the beryllium neutrino flux. As can be seen from Fig. 7(b), an additional enhancement appears in the regeneration region at $`\mathrm{\Delta }m^2=3\times 10^7\text{eV}^2`$ and the probability as a function of $`\mathrm{\Delta }m^2`$ becomes asymmetric in the regeneration region. In the VAC and QVO<sub>S</sub> regions the boron neutrino “wave” is modulated by the beryllium wave with a smaller amplitude. For the Ge–production rate, all the features of the curves are shifted (with respect to the Super–Kamiokande curves) to smaller values of $`\mathrm{\Delta }m^2`$ by a factor $`30`$. This feature is due the fact that the main contribution to $`Q_{Ge}`$ comes from the $`pp`$-neutrino flux with an average detected energy of 0.3 MeV, about 30 times smaller than the average energy of the boron neutrino flux.
Let us consider now the $`ϵ`$-dependence of the rates. We distinguish here between three different regions of $`\mathrm{\Delta }m^2`$:
1) The QVO<sub>L</sub> region with $`\mathrm{\Delta }m^2\text{ }>3\times 10^4\text{eV}^2`$. Here we have essentially averaged vacuum oscillations, so that the survival probability is given by Eq. (27).
2) The QVO<sub>S</sub> and VAC regions with $`\mathrm{\Delta }m^2\text{ }<10^8\text{eV}^2`$. Here we have essentially non-averaged or partially averaged vacuum oscillations which take the form of Eq. (29). In principle matter effects strongly suppress the oscillations inside the Sun and the Earth. However, the modification of the observables is small, since the size of the Sun (and the Earth) is much smaller than the oscillation length in vacuum.
3) The matter conversion (MSW) region is between these two QVO regions. The oscillation effect is strongly suppressed. As explained in Sec. II, Earth matter effect is important but insensitive to $`ϵ`$. The expression of $`P`$ in this region is given in Eqs. (42) and (45).
From Eqs. (27), (28), and (29) we can find straightforwardly the dependence of the observables on deviations from maximal mixing. For the VO, QVO<sub>S</sub> and QVO<sub>L</sub> regions, the following points are in order:
1. The observables depend very weakly on $`ϵ`$. Corrections to maximal mixing are of $`𝒪(ϵ^2)`$.
2. The dependence is symmetric under the exchange $`ϵϵ`$. Minimum of the survival probability, and consequently, minima of rates, are at $`ϵ=0`$, that is, at maximal mixing.
For the MSW regions, the following points are in order:
1. The survival probability and the rates depend linearly on $`ϵ`$. Corrections to the maximal mixing case are consequently larger.
2. The survival probability and the rates decrease with increase of $`ϵ`$.
There are two transition regions between VO and pure matter conversion. In these regions, the symmetric dependence of the observables transforms into a linear dependence and the sensitivity to deviation from maximal mixing increases. The ambiguity $`ϵϵ`$ disappears.
Thus, in the region of pure matter conversion (inside the Sun) the sensitivity of measurements of rates to a deviation from maximal mixing is maximal. It is in this region that the possibility of maximal mixing can be tested with the highest accuracy. The corresponding $`\mathrm{\Delta }m^2`$ range depends on the neutrino energy. For the highest energies of the solar neutrinos ($`E10`$ MeV) we get the range of maximal sensitivity: $`\mathrm{\Delta }m^2=10^710^4\text{eV}^2`$, while for the lowest energies ($`E0.3`$ MeV) the range is $`\mathrm{\Delta }m^2=3\times 10^93\times 10^6\text{eV}^2`$. Studying effects by the experiments with different energy thresholds we can get high sensitivity to deviations from maximal mixing in whole range of $`\mathrm{\Delta }m^2`$ excluding VO.
In the next two subsections we consider the dependence of specific rates on the oscillation parameters and evaluate the sensitivity of their present measurements to deviations from maximal mixing.
### B Ar–production rate.
Let us now study the implications of the results presented in Fig. 7(b) for the Ar–production rate and in Fig. 7(c) for the Super–Kamiokande rate in the various $`\mathrm{\Delta }m^2`$ regions. As can be seen in Fig. 7(b), the Ar–production rate for maximal mixing in all favorable regions (HIGH, LOW, VAC<sub>L</sub>) lies in the range $`R_{Cl}=0.50\pm 0.08`$ ($`Q_{Ar}=3.9\pm 0.6`$ SNU), which is $`2\sigma `$ above the Homestake result. The predictions for $`ϵ=+0.3`$ ($`0.3`$) are below (above) the values for maximal mixing by just about $`1\sigma `$. As mentioned above, the highest sensitivity of $`Q_{Ar}`$ to $`ϵ`$ is achieved already in the HIGH region, where $`Q_{Ar}`$ depends on $`ϵ`$ linearly. For the Super–Kamiokande rate, we get that for maximal mixing in all favorable regions $`R_{SK}=0.58\pm 0.12`$ which is $`1\sigma `$ above the Super–Kamiokande result. Moreover comparing Figs. 7(b) and 7(c) we see that the Ar–production rate and the Super–Kamiokande rate are strongly correlated as they are both dominated by the contribution from boron flux neutrinos.
Next, in order to reduce the SSM uncertainty we normalize the boron neutrino flux to the Super–Kamiokande rate. More precisely, for each pair of the of oscillation parameters ($`\mathrm{\Delta }m^2,ϵ`$) we find the boron neutrino flux which reproduces the Super–Kamiokande event rate. All other fluxes and their uncertainties are taken according to the BP98. Using this procedure we calculate $`Q_{Ar}`$ (and in the next subsection also $`Q_{Ge}`$). The results of this calculation are shown in Fig. 8.
From Fig. 8 we see that after boron flux normalization procedure the dependence of $`Q_{Ar}`$ on $`ϵ`$ is relatively weak since the ratio between the boron neutrino flux contribution to $`Q_{Ar}`$ and the contribution from charged current interactions to the Super–Kamiokande rate is independent of the survival probability and therefore of $`ϵ`$. The $`ϵ`$ dependence of $`Q_{Ar}`$ comes directly from the suppression of the beryllium and other (CNO and pep) neutrino fluxes, and indirectly from the contribution of neutral current interactions to the Super–Kamiokande rate. Expressing the boron neutrino flux $`f_B`$ via the SK rate:
$$f_B\frac{R_{SK}}{P(1r)+r},$$
(70)
we obtain:
$$Q_{Ar}\frac{Q_{Ar}^BR_{SK}}{(1r+r/P)}+Q_{Ar}^{\mathit{}}P^{}.$$
(71)
Here $`r`$ is the ratio of the $`\nu _\mu e`$ and $`\nu _ee`$ cross-sections, $`Q_{Ar}^B`$ and $`Q_{Ar}^{\mathit{}}`$ are the SSM contributions to the Ar–production rate from the boron neutrino flux and from all other low energy fluxes, respectively, and $`P`$ and $`P^{}`$ are the effective survival probabilities for the boron neutrino and low energy neutrino fluxes, respectively. The bands in Fig. 8 reflect $`1\sigma `$ errors due to the uncertainties in all, but the boron, neutrino fluxes.
As expected, in the QVO regions, QVO<sub>L</sub> and QVO<sub>S</sub>, the $`ϵ`$-dependence is symmetric around $`ϵ=0`$, $`Q_{Ar}=Q_{Ar}(ϵ^2)`$, while in the pure matter conversion regions, HIGH and LOW, $`Q_{Ar}`$ depends on $`ϵ`$ linearly. In the linear regime, the change in the Ar–production rate is $`\mathrm{\Delta }Q_{Ar}(0.70.8)`$ SNU for $`0.3\text{ }<ϵ\text{ }<+0.3`$. This is about $`3\sigma `$ for the present experimental error. In the quadratic regime the change is substantially smaller: $`\mathrm{\Delta }Q_{Ar}(0.10.2)`$ SNU. Clearly, the present sensitivity is not enough to draw definite conclusions. Moreover, even after normalization of the boron neutrino flux to the Super–Kamiokande rate, the predicted rate is higher than the Homestake result for all globally allowed values of $`\mathrm{\Delta }m^2`$. The only statement that one can make is that the Homestake result favors a significant deviation from maximal mixing in the MSW region. Therefore checks of the Homestake result and improved accuracy in measurements of $`Q_{Ar}`$ by a factor of two or higher would have important implications for maximal and near–maximal mixing.
### C Ge–production rate
Let us now study the implications of the results presented in Fig. 7(a) where we show the Ge–production rate as a function of $`\mathrm{\Delta }m^2`$ for different values of the $`ϵ`$ parameter and in and Fig. 9 where we plot $`Q_{Ge}`$ as a function of $`ϵ`$ in the various $`\mathrm{\Delta }m^2`$ regions. While the results plotted in Fig. 7(a) are obtained within the SSM, the predicted rates shown in Fig. 9 have been obtained after normalization of the boron neutrino flux to the Super–Kamiokande measured rate as described previously for Fig. 8. In this case, unlike in the case of Ar–production rate, the results are very slightly modified by the $`{}_{}{}^{8}B`$ flux normalization since the Ge–production rate is dominated by the contribution from the $`pp`$ neutrino flux and the corresponding theoretical uncertainties are smaller.
In the QVO<sub>L</sub> region, the rate is determined by averaged vacuum oscillations and therefore the expected rate is symmetric around maximal mixing. According to Eq. (29), $`P(ϵ=\pm 0.3)=0.545`$ which corresponds to $`Q_{Ge}=70.3\pm 3.5`$ SNU (with $`1\sigma `$ theoretical error). For exactly maximal mixing we have $`P(ϵ=0)=0.5`$ and the rate is minimal, $`Q_{Ge}=64.5\pm 3.5`$ SNU. As seen in Fig. 7(a) the bands for the different $`ϵ`$-values almost overlap and all the predictions are in agreement with the present experimental result. In other words, with present experimental error bars it is impossible to measure deviations from maximal mixing.
In the part of the HIGH region with small $`\mathrm{\Delta }m^2`$ the effects for $`pp`$ neutrinos occur in the transition between linear and quadratic $`ϵ`$ regimes. Consequently, the sensitivity of $`Q_{Ge}`$ to $`ϵ`$ is still below the present experimental accuracy. We find from Fig. 7(a): $`Q_{Ge}(ϵ=0)Q_{Ge}(ϵ=+0.3)6068`$ SNU, while $`Q_{Ge}(ϵ=0.3)70\pm 4`$ SNU, approximately $`1\sigma `$ (experimental) higher. From Fig. 9(a) we see that after the $`{}_{}{}^{8}B`$ normalization the variation of the $`Q_{Ge}`$ in $`ϵ`$ range $`(0.3,+0.3)`$ is $`\mathrm{\Delta }Q_{Ge}(1012)`$ SNU, which is about $`2\sigma `$ (present experimental error).
In the LOW region the sensitivity to $`ϵ`$ is maximal: the $`pp`$ neutrinos undergo pure matter conversion and the rate depends on $`ϵ`$ linearly. We get, for $`\mathrm{\Delta }m^2=10^7\text{eV}^2`$: $`Q_{Ge}(ϵ=0.3)=93\pm 4`$ SNU, $`Q_{Ge}(ϵ=0)=77.5\pm 3.5`$ SNU and $`Q_{Ge}(ϵ=0.3)=61\pm 3`$ SNU. The difference between the predicted values of $`Q_{Ge}(ϵ=0.3)`$ and $`Q_{Ge}(ϵ=+0.3)`$ is more than $`2\sigma `$ (experimental). From Fig. 9(b) we see that after the $`{}_{}{}^{8}B`$ normalization the variation of the $`Q_{Ge}`$ in $`ϵ`$ range $`(0.3,+0.3)`$ is $`\mathrm{\Delta }Q_{Ge}(3335)`$ SNU, which is more than $`5\sigma `$ (present experimental error). In this case one gets certain restrictions on $`ϵ`$, although the confidence level is low. For example, for $`\mathrm{\Delta }m^2=10^7\text{eV}^2`$, the combined SAGE and GALLEX+GNO result gives the $`1\sigma `$ range $`0.12ϵ+0.2`$. Therefore, further decrease of the experimental error bars by a factor of two, from the present 5 SNU to 3 - 4 SNU, could have important implications for the mixing provided that $`\mathrm{\Delta }m^2`$ will be fixed by some other independent measurement. Notice that in the LOW region one expects maximal regeneration effect for the $`pp`$-neutrinos which can be detected as, e.g., seasonal variation of the Ge–production rate .
The situation is similar in the QVO<sub>S</sub> region down to $`\mathrm{\Delta }m^2=5\times 10^9\text{eV}^2`$: the Ge–production rate depends on $`ϵ`$ linearly and $`\mathrm{\Delta }Q_{Ge}30`$ SNU (see Fig. 9(c)).
In the VO region, deviations from the maximal mixing result are determined by $`ϵ^2`$ and the variations (for a given $`\mathrm{\Delta }m^2`$) are small as shown in Fig. 9(d): $`\mathrm{\Delta }Q_{Ge}57`$ SNU. Thus the sensitivity of present data is still low and practically the whole interval $`ϵ(0.3,+0.3)`$ is allowed at $`(23)\sigma `$ level.
In consequence serious implications for maximal mixing require further decrease of the experimental error bars down to 3-4 SNU.
### D The Day–Night asymmetry in electron scattering events
In Fig. 6 we show contours of constant Day–Night asymmetry of the $`\nu e`$ scattering events in the $`\mathrm{\Delta }m^2ϵ`$ plane. The Super–Kamiokande 1117 days result
$$A_{\mathrm{N}\mathrm{D}}=2\frac{ND}{N+D}0.034\pm 0.026$$
(72)
excludes at the $`3\sigma `$ level the region $`A_{\mathrm{N}\mathrm{D}}\text{ }>0.11`$ which corresponds, at maximal mixing, to
$$4\times 10^7\text{ }<\mathrm{\Delta }m^2\text{ }<1.3\times 10^5\text{eV}^2.$$
(73)
The exclusion interval increases slightly with $`ϵ`$. The preferable regions of $`\mathrm{\Delta }m^2`$ for $`ϵ=0`$ are
$$\mathrm{\Delta }m^2=(2.510)\times 10^5\text{eV}^2,\mathrm{\Delta }m^2=(0.62.2)\times 10^7\text{eV}^2.$$
(74)
We emphasize that these results are SSM independent and have no ambiguities related to the analysis of the data.
### E Zenith angle distribution of electron scattering events
In Fig. 10 we show the zenith angle distribution of events in Super–Kamiokande for maximal and near–maximal mixing and for various values of $`\mathrm{\Delta }m^2`$ from the allowed regions. Significant enhancement of the night rate is expected in the HIGH and LOW regions.
In the HIGH region the distribution of events during the night is rather flat and the dependence on $`ϵ`$ is weak, so it will be difficult to use the shape to measure $`ϵ`$. The dependence on $`\mathrm{\Delta }m^2`$ is also weak. Maximal signal is expected in the N3 or/and N5 (core) bins.
In the LOW region the oscillations occur in the matter dominated regime (see Sec. II) when the oscillation length practically coincides with the refraction length, $`l_ml_0`$. For those trajectories crossing the mantle only (N1-N4), the latter can be approximately determined by the average density along the trajectory. Maximal effect (which corresponds to the oscillation phase $`=\pi `$) is realized for the trajectories with $`\mathrm{cos}\theta _Z=0.3`$, i.e. in the N2 bin (see Fig. 10(b)). The phase $`2\pi `$ is collected at $`\mathrm{cos}\theta _Z0.5`$ which corresponds to a minimum of the signal. Notice also that the zenith angle distribution depends on $`ϵ`$ very weakly. For a given trajectory the amplitude of the oscillation is determined by the mixing angle in the Earth matter
$$\mathrm{sin}^22\theta _E\frac{1}{\eta _E^2+12ϵ\eta _E^1}$$
(75)
where $`\eta _E`$ the parameter $`\eta `$ in the Earth defined in Eq. (33). Thus $`\mathrm{sin}^22\theta _E`$ is an increasing function of $`ϵ`$ (see also Eq. (18)). As a consequence the number of events in maxima increases with $`ϵ`$ as seen in Fig. 10(b). Present data do not show any enhancement in the N2 bin.
We conclude that precise measurements of the zenith angle distribution would allow the determination of $`\mathrm{\Delta }m^2`$ and probably resolve the HIGH/LOW ambiguity but are unlikely to play a significant role in the determination of $`ϵ`$.
### F The recoil electron energy spectrum
In Fig. 11 we show the expected recoil electron spectrum for maximal mixing with various values of $`\mathrm{\Delta }m^2`$ in the regions allowed by the global fit. In the HIGH and LOW regions the “reduced” spectra, $`R(E)`$, are flat. Strong distortion is expected in the VAC<sub>L</sub> region. Thus further improvement on the measurement of the recoil electron spectrum can discriminate between the MSW and VAC<sub>L</sub> regions.
For large enough $`\mathrm{\Delta }m^2`$ and $`ϵ0`$ a distortion in the spectrum appears due to the effect of the adiabatic edge (Fig. 12(a)). This feature can be traced from the dependence of the survival probability on $`\mathrm{\Delta }m^2/E`$ (see Fig. 1). For a positive $`ϵ`$ conversion inside the sun leads to negative slope of the reduced spectrum: $`R(E)`$ is larger at low energies. For negative $`ϵ`$, $`R(E)`$ increases with $`E`$ and the slope is positive. In the small $`\mathrm{\Delta }m^2`$ part of the LOW region the distortion of the spectrum is induced by the effect of the non-adiabatic edge (Fig. 1, $`\mathrm{\Delta }m^2/4E=(0.33)\times 10^{15}`$ eV. Here the situation is opposite to that for the HIGH region. The slope is positive for positive $`ϵ`$ and negative for negative $`ϵ`$ as illustrated in Fig. 12(b).
As follows from Figs. 11 and 12, the measurement of the electron energy spectrum provides information mainly on the value of $`\mathrm{\Delta }m^2`$ but it will be difficult to determine $`ϵ`$ by measuring the slope (first moment of the spectrum) in the interval $`0.3ϵ0.3`$ at Super–Kamiokande. At present Super–Kamiokande has presented the measured value $`T=8.14\pm 0.02`$ MeV. The precision of this measurement is to be compared with the maximum theoretically expected variation $`\mathrm{\Delta }T=|T(ϵ=0.3)T(ϵ=0.3)|<0.025`$ MeV which occurs for the two values of $`\mathrm{\Delta }m^2`$ shown in Fig. 12. Thus with the existing experimental precision, in the MSW region, the full range of $`ϵ`$ is allowed at $``$ 1$`\sigma `$.
## V Tests of Maximal Mixing in Future Experiments
In this section we consider the prospects of testing (near–)maximal mixing of the electron neutrinos in future experiments. We study various possibilities to measure the deviation $`ϵ`$ and we estimate the accuracy at which maximal mixing can be established or excluded. The extent to which the results of this Section hold in a three generation framework with $`|U_{e3}|0`$ is discussed in Sec. VI.
### A General requirements
There are two requirements for a precise determination of the mixing:
1. Uncertainties in the original neutrino fluxes should not play a role. For this purpose we will consider SSM independent observables, or at least observables which do not depend on the uncertainties in the boron neutrino flux.
2. At least two independent observables should be measured. As we have seen in Sec. III, the survival probability $`P_{ee}`$ and consequently all observables depend on both $`ϵ`$ and $`\mathrm{\Delta }m^2`$. Even in the case of maximal mixing when the solar survival probability is constant, $`P_D=1/2`$, a dependence of $`P_{ee}`$ on $`\mathrm{\Delta }m^2`$ appears due to the Earth regeneration effect.
Thus, to determine the mixing, one should find two independent observables which (i) are free of flux uncertainties, (ii) can be measured with high accuracy, and (iii) depend on different combinations of $`ϵ`$ and $`\mathrm{\Delta }m^2`$. In what follows we will identify such observables and study the accuracy at which mixing can be measured.
### B GNO and Super–Kamiokande
The main objective of the GNO experiment is to reach an accuracy $`3`$ SNU in the measurement of the Ge–production rate, $`Q_{Ge}`$. Also seasonal variations of $`Q_{Ge}`$ will be studied. Super–Kamiokande will continue to collect data for at least 10 years. With an energy threshold as low as 5 MeV the accuracy in measuring the Day–Night asymmetry will improve to $`0.0100.015`$.
Notice that, in the MSW region, $`A_{\mathrm{N}\mathrm{D}}`$ is mostly sensitive to $`\mathrm{\Delta }m^2`$, whereas $`Q_{Ge}`$ has strong dependence on $`ϵ`$. Therefore the pair of observables ($`A_{\mathrm{N}\mathrm{D}},Q_{Ge}`$) is, in principle, very useful for measurements of the oscillation parameters in the matter conversion region.
In Fig. 13 we show simultaneously contours of constant $`A_{\mathrm{N}\mathrm{D}}`$ at Super–Kamiokande site and $`Q_{Ge}`$ in the $`\mathrm{\Delta }m^2ϵ`$ plane. The iso–contours for $`Q_{Ge}`$ have been obtained for central values of the solar fluxes according to BP98. The theoretical (1$`\sigma `$) uncertainty is about $`\pm 2`$ SNU.
The iso-contours of $`Q_{Ge}`$ and $`A_{\mathrm{N}\mathrm{D}}`$ are nearly perpendicular to each other, which reflects the fact that these observables depend on different combinations of the oscillation parameters. However, the accuracy of the present experimental results is not enough to put statistically significant bounds on the mixing. The present experimental $`1\sigma `$ intervals are $`Q_{Ge}6980`$ SNU and $`A_{\mathrm{N}\mathrm{D}}0.010.06`$. The resulting constraints on the mixing parameters can be read from Fig. 13: in the HIGH region, $`0.2\text{ }>ϵ`$ and $`\mathrm{\Delta }m^2(1.58)\times 10^5\text{eV}^2`$, and in the LOW region $`0.10\text{ }<ϵ\text{ }<+0.15`$ and $`\mathrm{\Delta }m^2(0.52)\times 10^7\text{eV}^2`$. Inclusion of the theoretical uncertainties will expand these regions, especially in the HIGH case where the sensitivity of $`Q_{Ge}`$ to $`ϵ`$ is rather low.
Notice that the present Gallium data somewhat disfavor maximal mixing in the HIGH region. We estimate that the $`1\sigma `$ accuracy of the determination of $`ϵ`$ is
$$\mathrm{\Delta }ϵ0.150.20,(1\sigma ).$$
(76)
Within $`2\sigma `$ experimental errors the allowed regions cover most of HIGH parameter space of Fig. 13(a), and practically the entire LOW parameter space of Fig. 13(b). All values of $`ϵ(0.3,+0.3)`$ become allowed at the $`2\sigma `$ level.
With more data from GNO and higher statistics Super–Kamiokande measurements of the asymmetry one can reach better sensitivity.
### C SNO
The SNO experiment will study various observables in three types of processes:
1. Charged current neutrino-deuteron scattering: the total rate above a certain threshold (we denote the reduced rate of events by \[CC\]), the energy distribution of events (electron energy spectrum), and the time variation of events which can be characterized by the Day-Night asymmetry $`A_{\mathrm{N}\mathrm{D}}^{\mathrm{CC}}`$, the zenith angle distribution, and the seasonal asymmetry.
2. Neutrino-electron scattering: the total rate \[ES\], electron energy spectrum, and time variations.
3. Neutral current neutrino-deuteron scattering: the total rate \[NC\], and time variations.
Since the SNO observables depend on the boron neutrino flux only (we neglect the effect of the hep-neutrino flux), the ratios of rates are flux independent. Also relative time variations and energy spectrum distortion are flux independent. In what follows all the results will be presented for an energy threshold of 5 MeV.
In Fig. 14(a) we show the iso-contours of the double ratio \[ES\]/\[CC\] in the $`\mathrm{\Delta }m^2ϵ`$ plane. In terms of averaged survival probability, $`P(\mathrm{\Delta }m^2,ϵ)`$, the ratio can be written as
$$\frac{\text{[ES]}}{\text{[CC]}}1r+\frac{r}{P}=1r+\frac{2r}{f(1ϵ)}.$$
(77)
The last equality applies in the MSW region. The factor $`f=𝒪(1)`$ and deviates from unity due to Earth matter effects which are small (see Eq. (36)) and $`r`$ is the ratio between the $`\nu _xe`$ and $`\nu _ee`$ (ES) cross sections. The ratio \[ES\]/\[CC\] is larger than 1 and depends rather weakly on the oscillation parameters. From Eq. (77) we find the relation between the accuracy of measurements of \[ES\]/\[CC\], $`\mathrm{\Delta }(\text{[ES]/[CC]})`$ and the corresponding accuracy of determination of $`ϵ`$:
$$\mathrm{\Delta }ϵ\frac{f}{2r}\mathrm{\Delta }\left(\frac{\text{[ES]}}{\text{[CC]}}\right).$$
(78)
That is, the accuracy is lowered by factor $`1/2r3`$. As follows from Fig. 14(b), in the MSW region the ratio \[ES\]/\[CC\] increases with $`ϵ`$ for fixed $`\mathrm{\Delta }m^2`$. It varies within the limits \[ES\]/\[CC\]$`1.15\pm 0.10`$ for $`ϵ=0.0\pm 0.3`$. This variation is comparable with the expected $`1\sigma `$ error which is dominated by the uncertainty in the neutrino-deuteron cross-section ($`6\%`$ (statistical error has been calculated assuming 5000 CC and 500 ES events). Consequently no significant constraints on the oscillation parameters can be obtained, unless the uncertainty in the cross-section is reduced. According to Eq. (78), 10% precision in \[ES\]/\[CC\] corresponds to $`\mathrm{\Delta }ϵ0.3`$.
In Fig. 14(b) we show the iso-contours of the double ratio \[NC\]/\[CC\] in the $`\mathrm{\Delta }m^2ϵ`$ plane. In terms of $`P(\mathrm{\Delta }m^2,ϵ)`$, the ratio can be written as
$$\frac{\text{[NC]}}{\text{[CC]}}\frac{1}{P}=\frac{2}{f(1ϵ)},$$
(79)
where, again, the last equality is valid in the MSW region. From Eq. (79) we get for the accuracy in the determination of $`ϵ`$ in the MSW region:
$$\mathrm{\Delta }ϵ\frac{f}{2}\mathrm{\Delta }\left(\frac{\text{[NC]}}{\text{[CC]}}\right).$$
(80)
Here the prefactor is smaller than 1. Moreover, the double ratio \[NC\]/\[CC\] will be determined with much better accuracy than \[ES\]/\[CC\]. The theoretical uncertainties related to the neutrino-deuteron cross-sections essentially cancel out. The total $`1\sigma `$ error, which includes a statistical error for 5000 CC and 2000 NC events, is about $`3.6\%`$. According to Eq. (80), this corresponds to $`\mathrm{\Delta }ϵ0.04`$ for fixed $`\mathrm{\Delta }m^2`$. As follows from Fig. 14(a), in the MSW region the ratio \[NC\]/\[CC\] increases with $`ϵ`$ for fixed $`\mathrm{\Delta }m^2`$. It varies within the limits $`1.22.7`$ for $`ϵ`$ in the range $`(0.3,+0.3)`$. This variation is much larger than the expected $`1\sigma `$ error, $`\mathrm{\Delta }(\text{[NC]/[CC]})0.08`$ (for \[NC\]/\[CC\] = 2).
In the allowed regions of the oscillation parameters, the ratio \[NC\]/\[CC\] depends strongly on $`ϵ`$. A precise determination of $`\mathrm{\Delta }m^2`$ in these regions can be achieved from measurements of time variations, in particular, the Day- Night asymmetry $`A_{\mathrm{N}\mathrm{D}}^{\mathrm{CC}}`$. In Fig. 15 we show iso–contours of $`A_{\mathrm{N}\mathrm{D}}^{\mathrm{CC}}`$ and \[NC\]/\[CC\] in the $`\mathrm{\Delta }m^2ϵ`$ plane. Notice that the asymmetry of the CC-events is larger that the asymmetry of the ES-events for the same values of the oscillation parameters. The contours have weak dependence on $`ϵ`$. The combined analysis of \[NC\]/\[CC\] (sensitive mostly to $`ϵ`$) and $`A_{\mathrm{N}\mathrm{D}}^{\mathrm{CC}}`$ (sensitive to $`\mathrm{\Delta }m^2`$) can give a precise determination of the oscillation parameters. According to this figure, measurements yielding \[NC\]/\[CC\]$`2\times (1\pm 0.04)`$ and $`A_{\mathrm{N}\mathrm{D}}^{\mathrm{CC}}0.1\times (1\pm 0.3)`$ will determine $`ϵ`$ to an accuracy of order
$$\mathrm{\Delta }ϵ0.050.07(1\sigma ).$$
(81)
Notice that the same pairs of values ($`A_{\mathrm{N}\mathrm{D}}^{\mathrm{CC}}`$, \[NC\]/\[CC\]) appear in the HIGH and LOW regions. The LOW/HIGH ambiguity can be resolved by the KamLAND reactor experiment which will give a positive oscillation signal in the case of the HIGH solution. It can also be resolved by the Borexino experiment which will show strong earth regeneration effect in the LOW region as discussed next.
### D Borexino
The Borexino collaboration will measure the total rate of ES events and search for time variations of the signal.
In Fig. 16 we show iso-contours of the reduced rate $`R_{Be}`$ (suppression factor) and iso-contours of the Day-Night asymmetry in the $`\mathrm{\Delta }m^2ϵ`$ plane. As discussed before for Super-Kamiokande and SNO, the Day-Night asymmetry is sensitive mainly to $`\mathrm{\Delta }m^2`$ while the deviation from maximal mixing can be restricted by the rate for which we can write
$$R_{Be}=r_{Be}+(1r_{Be})P,$$
(82)
where $`r_{Be}0.24`$ is the ratio of the $`\nu _\mu e`$ to $`\nu _ee`$ cross-sections for the beryllium neutrinos.
From Eq. (82) and Fig. 16 we observe the following behaviours:
(i) In the QVO<sub>L</sub> region, the survival probability $`P`$ depends quadratically on $`ϵ`$. We find $`R_{Be}0.62`$ with very weak $`ϵ`$ dependence.
(ii) In the HIGH region, the transition between the quadratic and linear $`ϵ`$ dependence occurs. For $`\mathrm{\Delta }m^210^5`$ eV<sup>2</sup> the rate increases from 0.54 to 0.71 when $`ϵ`$ decreases from $`+0.3`$ to $`0.3`$. The Day–Night asymmetry is very small here.
(iii) In contrast, in the LOW region Borexino has higher sensitivity to the oscillation parameters. For maximal mixing, $`R_{Be}`$ is in the interval 0.65 - 0.70, and the asymmetry can be as large as 30%. The probability depends linearly on $`ϵ`$, so that we get from Eq. (82)
$$\mathrm{\Delta }ϵ\frac{2}{1r_{Be}}\mathrm{\Delta }R_{Be}.$$
(83)
Still, the sensitivity to $`ϵ`$ is low but Borexino will play a crucial role in fixing $`\mathrm{\Delta }m^2`$ and in particular by the measurement of the Day–Night asymmetry will be able to resolve the HIGH/LOW ambiguity which may remain after the measurements at Super–Kamiokande and SNO.
### E LowNu Experiments
A new generation of experiments aiming at a high precision real time measurement of the low energy solar neutrino spectrum is now under study . Some of them such us HELLAZ, HERON and SUPER–MuNu intend to detect the elastic scattering of the electron neutrinos with the electrons of a gas and measure the recoil electron energy and its direction of emission. The proposed experiment LENS plans to detect the electron neutrino via its absorption in a heavy nuclear target with the subsequent emission of an electron and a delayed gamma emission . The expected rates at those experiments for the proposed detector sizes are of the order of $`1`$$`10`$ $`pp`$ neutrinos a day. Consequently, with a running time of two years they can reach a sensitivity of a few percent in the total neutrino rate at low energy, provided that they can achieve sufficient background rejection. This would allow the determination of $`ϵ`$ with a similar precision of a few percent, in particular in the QVO<sub>S</sub>, where SNO and Borexino cannot give much information due to their higher energy threshold.
Sensitivity to the oscillation parameters in the LMA region can also be achieved in experiments detecting low energy $`\overline{\nu }_e`$ fluxes from nuclear reactors (detection threshold is about $`E_{th}=1.8`$ MeV). The Borexino experiment, in addition to detecting solar neutrinos, aims to detect the diffuse fluxes from nuclear reactors in Europe, mainly in France and Germany, at an average distance $`800`$ km. The KamLAND experiment aims at detecting the low energy diffuse $`\overline{\nu }_e`$ fluxes from reactors in Japan from an average distance of $`200`$ km. Both experiments can provide important information in discriminating between the HIGH and LOW solutions. However, they are expected to be very weakly sensitive to $`ϵ`$. Given the short distances traveled by the neutrinos, matter effects are negligible for both experiments. Consequently, the survival probability for these experiments takes the vacuum form of Eqs. (27) and (29) which depend only quadratically on $`ϵ`$. With the expected achievable limits on the survival probability of about $`10`$–20%, a measurement of $`|ϵ|<0.3`$ seems unfeasible.
## VI Effects of the third neutrino
In the general case of three neutrino mixing when $`|U_{e3}|0`$ the definition of deviation from maximal mixing becomes ambiguous. Formally, maximal mixing of the $`\nu _e`$ in the three neutrino context can be defined as the equality $`\theta _{12}=\pi /4`$ where $`\theta _{12}`$ is the rotation angle in the plane of the first and second mass eigenstates. Phenomenologically, deviations from maximal mixing can be defined as a deviation of the averaged probability from 1/2 or a deviation of the depth of oscillations from 1.
Let us consider neutrino mass spectrum which explains also the atmospheric neutrino problem via (mainly) $`\nu _\mu \nu _\tau `$ oscillations. This implies
$$|\mathrm{\Delta }m_{}^2||\mathrm{\Delta }m_{21}^2||\mathrm{\Delta }m_{31}^2||\mathrm{\Delta }m_{32}^2|=|\mathrm{\Delta }m_{\mathrm{atm}}^2|.$$
(84)
In this case (as far as solar neutrinos are concerned) the oscillations driven by $`\mathrm{\Delta }m_{\mathrm{atm}}^2`$ can be averaged. This situation justifies our use of $`\eta _S`$ as defined in Eq. (22) and $`\eta _E`$ as defined in Eq. (33) in the calculation of the survival probabilities. We get the following expression for the survival probability :
$$P_{ee}=(1|U_{e3}|^2)^2P_2+|U_{e3}|^4,$$
(85)
where $`P_2=P(\mathrm{\Delta }m^2,\theta ,V)`$ is the two neutrino mixing survival probability determined by
$$\mathrm{\Delta }m^2=\mathrm{\Delta }m_{}^2,\mathrm{tan}\theta =|U_{e2}/U_{e1}|,V=V_0(1|U_{e3}|^2).$$
(86)
where $`V_0`$ is the matter potential $`V_0=\sqrt{2}G_F\rho Y_e/m_N`$. In principle we can define a parameter $`ϵ_2`$ which describes the deviation from maximal mixing of $`\theta `$ defined in Eq. (86) similarly to the two neutrino parameters $`ϵ`$ and $`\theta _{ex}`$ of Eq. (3). Then the deviation from maximal mixing in the three neutrino context will depend on the specific physical situation. If the mass splitting $`\mathrm{\Delta }m_{}^2`$ and the mixing angle $`\theta `$ induce vacuum oscillations, then $`P_2(1+ϵ_2^2)/2`$ and Eq. (85) gives
$$P_{ee}\frac{1}{2}+\frac{ϵ_2^2}{2}|U_{e3}|^2.$$
(87)
Thus one can define the deviation from maximal mixing in the three neutrino context in this case as $`ϵ_3^v=ϵ_2^22|U_{e3}|^2`$. Clearly, the three neutrino case reduces to the two neutrino case if
$$|U_{e3}|^2\frac{1}{2}ϵ_2^2.$$
(88)
Taking $`|U_{e3}|^2`$ at the level of the CHOOZ bound , $`|U_{e3}|^20.05`$, Eq. (88) gives $`|ϵ_2|0.3`$.
If $`\mathrm{\Delta }m_{}^2`$ and $`\theta `$ lead to the adiabatic conversion in matter, then $`P_2(1ϵ_2)/2`$ and Eq. (85) gives
$$P_{ee}\frac{1}{2}\frac{ϵ_2}{2}|U_{e3}|^2.$$
(89)
Here it is useful to characterize the deviation from maximal mixing through $`ϵ_3^m=ϵ_2+2|U_{e3}|^2`$. Corrections due to the third generation can be neglected provided that
$$|U_{e3}|^2\frac{1}{2}ϵ_2.$$
(90)
For $`|U_{e3}|^20.05`$, Eq. (90) gives $`ϵ_20.1`$.
We have verified that, in the case of adiabatic transitions, the results of our calculations of the expected rates in the two-neutrino mixing scenario can be translated to a good approximation to the case of three-neutrino mixing with the simple replacement of $`ϵ_2`$ with $`ϵ_3^m`$. This applies, for instance, to the contours for the Ar–production rate in Fig. 6 and the predictions in Fig. 7 and Fig. 8 in the range $`10^5\mathrm{\Delta }m^2`$/eV$`{}_{}{}^{2}10^8`$. It also applies to the Ge–production rate in Fig. 13(b) and the corresponding predictions in Fig. 7 and Fig. 9(b) and 9(c). For the predictions of the SNO rates it can be used for Figs. 15 and 14 in the range $`10^4\mathrm{\Delta }m^2`$/eV$`{}_{}{}^{2}10^7`$, and for the predicted rates at Borexino in Fig. 16(b) for $`\mathrm{\Delta }m^2`$/eV$`{}_{}{}^{2}5\times 10^9`$.
In the case of averaged vacuum oscillations (QVO<sub>L</sub>) the results for the three-neutrino scenario can be read from the results presented here with the replacement of $`ϵ_2^2`$ with $`ϵ_3^v`$. For long wavelength oscillations (QVO<sub>S</sub> and VO), the results for the three-neutrino scenario cannot be directly derived from our results. Thus our predictions in these regions only hold for very small values of the mixing angle $`|U_{e3}|`$, well below the present CHOOZ bound.
Concerning the value of $`U_{e3}`$, although certain improvement on the present CHOOZ bound may be expected from long baseline experiments, such as K2K and MINOS , their final sensitivity is still unclear as it depends on their capability of discriminating against the $`\nu _e`$ beam contamination. Ultimate sensitivity can be achieved at experiments performed with neutrino beams from muon-storage rings at the so-called neutrino factories .
## VII Maximal mixing and other experiments
In the previous sections we have concentrated on effects in the solar neutrinos. Mixing of the electron neutrino can be probed in a number of other experiments.
### A Atmospheric neutrinos
Maximal and near–maximal mixing of the electron neutrinos can be probed in the atmospheric neutrino studies. The oscillations in the Earth matter with parameters from the LMA or HIGH regions can give an observable effect in the $`e`$-like events. The electron neutrino flux at the detector can be written as
$$F_e=F_e^0\left[1+P_{e\mu }(r\mathrm{cos}^2\theta _{23}1)\right],$$
(91)
where $`rF_\mu ^0/F_e^0`$ is the ratio of the original electron and muon neutrino fluxes, $`P_{e\mu }=P(\mathrm{\Delta }m_{}^2,ϵ_2)`$ is the two neutrino transition probability, and $`\theta _{23}`$ is the $`\nu _\mu \nu _\tau `$ mixing responsible for the dominant mode of the atmospheric neutrino oscillations. The transition probability can be of order one at $`\mathrm{\Delta }m^2>3\times 10^4`$ eV<sup>2</sup>. It decreases fast with $`\mathrm{\Delta }m^2/E`$ due to matter suppression of the mixing. Thus the biggest effect is expected in the low energy (sub-GeV) events sample. Notice that the probability $`P_{e\mu }`$ enters in Eq. (91) with a “screening factor” $`(r\mathrm{cos}^2\theta _{23}1)`$ which turns out to be small. Indeed, for the sub-GeV sample $`r2`$ and the screening factor is exactly zero for maximal mixing in the atmospheric neutrinos . The factor equals approximately $`\mathrm{cos}2\theta _{23}`$, so that for $`\mathrm{sin}^22\theta _{23}=0.95`$ we get about 0.22.
For $`\theta _{23}<\pi /4`$ the oscillations lead to the excess of the $`e`$-like events. Indeed some excess is hinted by the SK data. The excess can be defined as $`N_e/N_e^01`$, where $`N_e`$ and $`N_e^0`$ are the numbers of events with and without oscillations. This excess can be written in a matter dominant regime of oscillations ($`\eta _E1`$) as
$$\frac{N_e}{N_e^0}1\mathrm{cos}2\theta _{23}\eta _E^2(1ϵ_2\eta _E),$$
(92)
where $`ϵ_2`$ was defined below Eq. (86). The excess depends on $`ϵ_2`$ linearly and it increases with $`\eta _E`$. However it is even more sensitive to deviations of $`\theta _{23}`$ from the maximal mixing value. For $`\mathrm{\Delta }m_{21}^2=10^4`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta _{23}=0.95`$ the excess reaches about 3% and the dependence on $`ϵ`$ is weak. For $`\mathrm{\Delta }m_{21}^2=210^4`$ eV<sup>2</sup> and the same $`\theta _{23}`$, we find the excess about 4.5 %. Significant dependence on $`ϵ`$ appears for $`\mathrm{\Delta }m_{12}^2=3\times 10^4`$ eV<sup>2</sup>. However, it is unlikely given the size of the effect that atmospheric neutrino data will give any significant information on the value of $`ϵ_2`$.
### B Supernova neutrinos
Maximal or near–maximal mixing of the electron neutrinos will significantly modify properties of neutrino bursts from supernovae. The effects depend crucially on features of the whole neutrino mass spectrum and in particular on the value of $`U_{e3}`$ and whether the mass hierarchy is normal ($`\mathrm{\Delta }m_{31}^2>0`$) or inverted ($`\mathrm{\Delta }m_{31}^2<0`$). Let us summarize here the main results (for more details see and references therein).
All oscillation and conversion effects in supernova neutrinos are determined by the total survival probability of the electron neutrinos which in this subsection we will write as $`p`$, and total survival probability of the electron antineutrinos, $`\overline{p}`$. (This property is related to the fact that the original spectra of the muon and the tau neutrinos are identical and that the muon and tau neutrinos cannot be distinguished at the detection point).
The probabilities should include the effects of propagation inside the star, on the way to the Earth and inside the Earth. Using the probabilities $`p`$ and $`\overline{p}`$, one can write the fluxes of the electron neutrinos, $`F_e`$, and electron antineutrinos, $`F_{\overline{e}}`$, at the detector in terms of the original electron (anti)neutrino fluxes, $`F_e^0`$ and $`F_{\overline{e}}^0`$, and the non-electron neutrino flux $`F_x^0F_\mu ^0=F_\tau ^0=F_{\overline{\mu }}^0=F_{\overline{\tau }}^0`$:
$$F_e=pF_e^0+(1p)F_x^0,F_{\overline{e}}=\overline{p}F_{\overline{e}}^0+(1\overline{p})F_x^0.$$
(93)
In general, $`p`$, and $`\overline{p}`$ depend on the neutrino energy.
Let us summarize the results for specific neutrino mass and flavor spectra.
1) If $`|U_{e3}|^2>3\times 10^4`$ the conversion in the resonance related to the largest (atmospheric) splitting ($`\mathrm{\Delta }m_{31}^2`$) will be completely adiabatic and the final effect depends on the type of mass hierarchy. In the case of normal hierarchy ($`\nu _3`$ is the heaviest state) the resonance conversion occurs in the neutrino channel and for the survival probability we get
$$p|U_{e3}|^21.$$
(94)
This probability is practically independent of the properties (mass, flavour) of the first and second mass eigenstates. In particular, there is no sensitivity to $`ϵ`$ and no Earth matter effect is expected for neutrinos.
In contrast, the antineutrino channels will not be affected by the high resonance and $`\overline{p}`$ will be determined by physics of the two light levels.
For parameters in the HIGH and LOW regions, the neutrino propagation in the star is adiabatic, so that the survival probability in the star equals
$$\overline{p}\mathrm{cos}^2\theta _{12}=\frac{1}{2}(1+ϵ).$$
(95)
This probability can be further modified due to oscillations in the matter of the Earth. Thus, we expect the following consequences: (i) disappearance of the $`\nu _e`$ neutronization peak; (ii) hard $`\nu _e`$ spectrum (coinciding with the original $`\nu _\mu `$) spectrum at the cooling stage:
$$F_eF_x^0;$$
(96)
(iii) composite $`\overline{\nu }_e`$ spectrum:
$$F_{\overline{e}}=\frac{1}{2}(F_{\overline{e}}^0+F_x^0)\frac{ϵ}{2}(F_x^0F_{\overline{e}}^0);$$
(97)
(iv) strong Earth matter effect (which leads to different signals at various detectors). For the HIGH mass range, the Earth effect is maximal in the high energy part of the spectrum, $`E>20`$ MeV, whereas for the LOW solution the largest effect is in the low energy part.
According to Eq. (97), the $`ϵ`$-dependent term is proportional to the difference of the original fluxes. Thus due to the uncertainties in the predicticted fluxes it will be difficult to measure $`ϵ`$. In order to reduce the theoretical uncertainty, one could in principle compare numbers of $`\nu _e`$ and $`\overline{\nu }_e`$ events at large $`E`$ which are determined by, respectively, $`F_x^0`$ and $`F_x^0(1ϵ)/2`$ and are proportional to the same flux.
For parameters of the two light states in the VO region, the neutrino propagation in the star is non-adiabatic, so that the survival probability can be writen as
$$\overline{p}(1P_c)\mathrm{cos}^2\theta _{12}+P_c\mathrm{sin}^2\theta _{12}=\frac{1}{2}+ϵ\left(\frac{1}{2}P_c\right),$$
(98)
where the jump probability $`P_c`$ depends on the details of the density profile in the outer regions of the star and cannot be reliably predicted. Practically, the probability should lie between the adiabatic value (95) and the pure vacuum oscillation expression $`p=(1+ϵ^2)/2`$.
In the case of inverted mass hierarchy the sensitivity to $`ϵ`$ appears in the neutrino channel and neutrinos and antineutrinos interchange their roles. Now the resonance is in the antineutrino channel so that
$$\overline{p}|U_{e3}|^21,$$
(99)
and the oscillations in the neutrino channels will be determined by physics of the two light levels. For parameters from the HIGH and LOW regions the propagation in the star is adiabatic and
$$p\mathrm{sin}^2\theta _{12}=\frac{1}{2}(1+ϵ).$$
(100)
Moreover, Earth matter effects are expected for neutrinos. For the VO region we find, similarly to Eq. (98),
$$p\frac{1}{2}ϵ\left(\frac{1}{2}P_c\right).$$
(101)
In this inverted scheme we predict: (i) partial disappearance of the $`\nu _e`$ neutronization peak; (ii) hard spectrum of the electron antineutrinos:
$$F_{\overline{e}}F_x^0;$$
(102)
(iii) composite $`\nu _e`$ spectrum:
$$F_e=\frac{1}{2}(F_e^0+F_x^0)+\frac{ϵ}{2}(F_x^0F_e^0);$$
(103)
(iv) Earth matter effects are expected in the neutrino channel only.
2) If $`|U_{e3}|^2<3\times 10^6`$, the effect of the third neutrino can be neglected: in the resonance channel the transition driven by $`U_{e3}`$ is strongly non-adiabatic, and in the non-resonance channel the mixing is always very small. In this case the problem is reduced to transitions in two level system with parameters determined by $`\mathrm{\Delta }m_{21}^2=\mathrm{\Delta }m_{}^2`$ and $`ϵ`$. As a result both neutrino and antineutrino channels turn out to be sensitive to $`ϵ`$. The effects include those considered above both for normal and inverted mass hierarchy.
For HIGH and LOW regions of parameters the propagation proceeds adiabaticaly, and for the survival probabilities we get the expressions given in (95) and (100). Correspondingly, neutrino and antineutrino spectra will be given by Eq. (97) and (103). Thus we predict that both neutrino and antineutrino spectra will be composite, consisting of nearly equal admixture of the soft and hard components. In the high energy part where effects of the soft components can be neglected we get from Eqs. (97) and (103):
$$\frac{F_e}{F_{\overline{e}}}1+2ϵ.$$
(104)
That is, larger $`\nu _e`$ signal (as compared with $`\overline{\nu _e}`$) is expected for $`ϵ>0`$ and smaller for $`ϵ<0`$. Also Earth matter effect is expected in both neutrino and antineutrino channels.
3) If $`|U_{e3}|^2`$ is in the intermediate region, $`3\times 10^63\times 10^4`$, the adiabaticity in high mass resonance is partially violated and we expect some intermediate situation between those described in 1) and 2). In particular, both $`\nu _e`$ and $`\overline{\nu }_e`$ spectra will be composite, however admixtures of the soft and hard components will be unequal, etc..
To conclude, one expects strong influence of maximal and near–maximal mixing on the properties of the neutrino bursts. However, the uncertainties in the predicted neutrino spectra will make it difficult to obtain high sensitivity to $`ϵ`$.
Notice also that the analysis of the SN1987A data gives the 99% CL bound on $`p>0.65`$. This would correspond to $`\overline{p}>0.3`$ . Some recent calculations show that the difference between $`\overline{\nu }_e`$ and $`\overline{\nu }_\mu `$ original spectra can be rather small, which would somewhat relax the above bound.
### C Neutrinoless double beta decay
The effective Majorana mass of the electron neutrino $`m_{ee}`$ relevant for the neutrinoless double beta decay is sensitive to the distribution of the electron neutrino flavor in the mass eigenstates (see for recent discussion). The contribution to the effective mass $`m_{ee}`$ from the two mass eigenstates responsible for the solar neutrino conversion can be written in terms of the oscillation parameters as
$$m_{ee}=\frac{1}{2}\left|m_1(1+ϵ)+\sqrt{m_1^2+\mathrm{\Delta }m_{}^2}(1ϵ)e^{i\varphi _{12}}\right|,$$
(105)
where $`m_1`$ is the mass of the first eigenstate and $`\varphi _{12}`$ is the relative phase of the first and the second mass eigenvalues. In the case of strong mass hierarchy, $`m_1^2\mathrm{\Delta }m_{}^2`$ we get from Eq. (105)
$$m_{ee}\frac{1}{2}\sqrt{\mathrm{\Delta }m_{}^2}(1ϵ).$$
(106)
According to this equation in the HIGH region the effective mass can be as big as $`(12)\times 10^2`$ eV which can be probed at the next generation of the double beta decay experiments . Notice that the contribution from the third mass eigenstate is strongly restricted by present experimental bound $`m_{ee}<2\times 10^3`$ eV.
Although the dependence of $`m_{ee}`$ on $`ϵ`$ is rather strong, it will be difficult to measure $`ϵ`$ due to the uncertainties in the nuclear matrix elements. Eq. (106) can be considered as a test equation: if the measured values of $`m_{ee}`$, $`ϵ`$ and $`\mathrm{\Delta }m_{}^2`$ indeed satisfy this equation (within experimental and theoretical uncertainties) it will testify for the validity of whole scheme.
In the case of strong mass degeneracy, $`m_1^2\mathrm{\Delta }m_{}^2`$ we get
$$m_{ee}\frac{1}{2}m_1\left|(1+ϵ)+(1ϵ)e^{i\varphi _{12}}\right|,$$
(107)
where both non-oscillation parameters $`m_1`$ and $`\varphi _{12}`$ are unknown. For $`\varphi _{12}=\pi `$ and $`(0)`$ the mass equals $`m_{ee}=m_1ϵ`$ $`(m_1)`$, so that $`ϵ`$ determines the lower bound on $`m_{ee}`$. If $`m_{ee}`$ and $`ϵ`$ are measured, the above equalities will determine the upper bound and the lower bounds on the absolute scale of the neutrinos mass: $`m_{ee}<m_1<m_{ee}/ϵ`$.
In the case of inverted mass hierarchy the two states responsible for the solar neutrino conversion are degenerate: $`m_1m_2\sqrt{\mathrm{\Delta }m_{\mathrm{atm}}^2}`$ and the effective majorana mass can be written as
$$m_{ee}\frac{1}{2}\sqrt{\mathrm{\Delta }m_{\mathrm{atm}}^2}\left|(1+ϵ)+(1ϵ)e^{i\varphi _{12}}\right|.$$
(108)
In this case the measurement of $`ϵ`$ will allow us to determine the phase $`\varphi _{12}`$. According to Eq. (108), $`\sqrt{\mathrm{\Delta }m_{\mathrm{atm}}^2}ϵ<m_{ee}<\sqrt{\mathrm{\Delta }m_{\mathrm{atm}}^2}`$ which can be used as a test inequality for a given scheme.
Thus measurements of $`ϵ`$ in the oscillation experiments will allow to determine or restrict the effective mass $`m_{ee}`$ in the context of certain schemes of neutrino masses and mixing.
## VIII Conclusions
In this work we have explored the phenomenological consequences of (near–)maximal mixing of electron neutrinos with other standard neutrinos. The possibility of such maximal or near–maximal lepton mixing constitutes an intriguing challenge for fundamental theories of flavour. Our aim was twofold. First we have formulated the present status of maximal mixing of $`\nu _e`$ in the light of existing experimental data from solar neutrino experiments. Second we have explored the best ways to measure deviations from such maximal mixing at future experiments.
We show in Sec. II that both probabilities and observables depend on $`ϵ`$ quadratically in the regions of $`\mathrm{\Delta }m^2`$ where the effects are due to vacuum oscillations, and they depend on $`ϵ`$ linearly when matter effects dominate. Consequently, for $`|ϵ|1`$ the highest sensitivity to deviation from maximal mixing can be achieved in the $`\mathrm{\Delta }m^2`$ ranges of the MSW effect.
The results of a global fit to the existing solar neutrino data are presented in Sec. III and summarized in Figs. 2 $``$5. From this analysis we find that values of the mixing parameter $`|ϵ||12\mathrm{sin}^2\theta |<0.3`$ are allowed at 99% or lower CL for $`\mathrm{\Delta }m^2\text{ }>1.5\times 10^5`$ eV<sup>2</sup> (which contains the HIGH and QVO<sub>L</sub> regions) and for $`4\times 10^{10}`$ eV$`{}_{}{}^{2}\text{ }<\mathrm{\Delta }m^2\text{ }<2\times 10^7`$ eV<sup>2</sup> (which contains the defined LOW, QVO<sub>S</sub> and upper VAC<sub>L</sub> regions).
The role of the individual existing experiments on the determination of these regions is discussed in Sec. IV. We conclude that the present sensitivity to the mixing angle arises from the measurements of total event rates. The present data from Homestake experiment in Ar–production rate gives the strongest constraint on maximal or near–maximal mixing as it favours a significant deviation from $`ϵ=0`$. This conclusion is independent of the existing theoretical uncertainty on the boron flux as discussed in Sec. IV B. On the other hand the measurement of the zenith angle dependence and the recoil electron energy spectrum are important in the determination of the allowed mass ranges but are very weakly sensitive to deviation from maximal mixing. With the present existing sensitivity all values of $`|ϵ|<0.3`$ are allowed within 4$`\sigma `$.
In Sec. V we have discussed the ways to improve our knowledge on deviations from maximal mixing at future experiments. We concentrate on observables which are SSM (or at least boron flux) independent. In order to determine both the mass and the mixing we study pairs of observables. First we have looked at the maximal sensitivity which may be achievable on the presently running experiments GNO and Super–Kamiokande. In principle the measurement of the Ge–production rate at GNO and the Day–Night asymmetry at Super–Kamiokande and SNO can give crossed information on the oscillation parameters in the matter conversion region. In practice, however, the expected sensitivity is not enough to substantially improve the present knowledge on $`\mathrm{\Delta }m^2`$ and $`ϵ`$. The role of SNO and Borexino experiments is discussed in Secs. V C and V D. We show that with the expected theoretical and statistical uncertainty the most sensitive observable to the mixing angle is the rate \[NC\]/\[CC\] measurable at SNO. For instance, a measurement yielding \[NC\]/\[CC\] $`2\times (1\pm 0.04)`$ and $`A_{\mathrm{N}\mathrm{D}}^{\mathrm{CC}}0.1\times (1\pm 0.3)`$ will determine $`ϵ`$ to an accuracy of order $`\mathrm{\Delta }ϵ0.07`$. There exist however an ambiguity on the allowed mass range between HIGH and LOW regions. We show that the LOW/HIGH ambiguity can be resolved by the measurement of the Day–Night asymmetry at Borexino experiment which is sensitive to strong Earth regeneration effect in the LOW region or by the detection of oscillations in long baseline reactor experiments such as KamLand. However no substantial improvement on the knowledge of $`ϵ`$ is expected neither from Borexino nor from the new generation of low energy experiments either with solar or reactor neutrinos.
###### Acknowledgements.
We thank J. N. Bahcall, P. I. Krastev and E. Lisi for valuable discussions. We are particularly indebted to Y. Suzuki for providing us with Super–Kamiokande data on the night and day spectra. YN is partially supported by the Department of Energy under contract No. DE–FG02–90ER40542, by the Ambrose Monell Foundation, by AMIAS (Association of Members of the Institute for Advanced Study), by the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities, and by the Minerva Foundation (Munich). AYS acknowledge partial support from NSF grant No. PHY95-13835 to the Institute for Advanced Study. This work was also supported by the spanish DGICYT under grants PB98-0693 and PB97-1261, by the Generalitat Valenciana under grant GV99-3-1-01 and by the TMR network grant ERBFMRXCT960090 of the European Union.
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# Geometry of the smooth dual of 𝐺𝐿(𝑛)
## Abstract
Let $`𝒜(n)`$ be the smooth dual of the $`p`$-adic group $`G=GL(n)`$. We create on $`𝒜(n)`$ the structure of a complex algebraic variety. There is a morphism of $`𝒜(n)`$ onto the Bernstein variety $`\mathrm{\Omega }G`$ which is injective on each component of $`𝒜(n)`$. The tempered dual of $`G`$ is a deformation retract of $`𝒜(n)`$. The periodic cyclic homology of the Hecke algebra of $`G`$ is isomorphic to the periodised de Rham cohomology supported on finitely many components of $`𝒜(n)`$.
Introduction
Let $`F`$ be a finite extension of $`_p`$ and let $`G=GL(n)=GL(n,F)`$. A representation of $`G`$ on a complex vector space $`V`$ is smooth if the stabilizer of each vector in $`V`$ is an open subgroup of $`G`$. The set of equivalence classes of irreducible smooth representations of $`G`$ is the smooth dual $`𝒜(n)`$ of $`G`$. The Bernstein variety $`\mathrm{\Omega }G`$ is the set of $`G`$-conjugacy classes of pairs $`(M,\sigma )`$ where $`M`$ is a Levi subgroup of $`G`$ and $`\sigma `$ is an irreducible supercuspidal representation of $`M`$. Each irreducible smooth representation of $`G`$ is a subquotient of an induced representation $`i_{GM}\sigma `$. The pair $`(M,\sigma )`$ is unique up to conjugacy. This creates a finite-to-one map, the infinitesimal character, from $`𝒜(n)`$ onto $`\mathrm{\Omega }G`$.
The set $`\mathrm{\Omega }G`$ has the structure of a complex algebraic variety. This structure comes ultimately from twisting a supercuspidal representation $`\sigma `$ of $`M`$ by unramified quasicharacters of $`M`$. The set $`\mathrm{\Psi }(M)`$ of all unramified quasicharacters of $`M`$ has the structure of a complex torus.
The local Langlands correspondence is a natural bijection from $`𝒢(n)`$ to $`𝒜(n)`$ where $`𝒢(n)`$ is the set of equivalence classes of admissible $`n`$-dimensional complex representations of the Weil-Deligne group $`W_F^{}`$.
Each representation of the Weil group $`W_F`$ may itself be twisted by an unramified quasicharacter of $`W_F`$. In this way, by adapting Bernstein’s construction, we create the structure of complex algebraic variety on $`𝒢(n)`$. We then transport this structure to $`𝒜(n)`$ by the local Langlands correspondence. So $`𝒜(n)`$ acquires the structure of complex algebraic variety with infinitely many components. We obtain the following results.
###### Theorem 1
The smooth dual $`𝒜(n)`$ has a natural structure of complex algebraic variety. There is a canonical morphism of $`𝒜(n)`$ onto the Bernstein variety $`\mathrm{\Omega }G`$ which is injective on each component of $`𝒜(n)`$.
The Langlands parameters effect a stratification of the smooth dual $`𝒜(n)`$.
###### Theorem 2
The tempered dual of $`G`$ is a deformation retract of $`𝒜(n)`$ .
We view the deformation retraction as the geometric counterpart of the assembly map in the Baum-Connes conjecture for $`GL(n)`$, see .
Let $`(G)`$ be the Hecke algebra of compactly supported, uniformly locally constant complex-valued functions on $`G`$.
###### Theorem 3
The periodic cyclic homology of $`(G)`$ is isomorphic to the periodised de Rham cohomology supported on finitely many components of the smooth dual of $`G`$.
From the point of view of noncommutative geometry, the variety $`𝒜(n)`$ is an affine scheme underlying the non-unital noncommutative ring $`(G)`$.
We find it interesting that the complex algebraic variety structure on $`𝒜(n)`$ has become visible only after employing the local Langlands correspondence, in other words only after paying attention to the number theory of the local field $`F`$.
1. The Complex Algebraic Variety $`𝒢(n)`$
The Weil group $`W_F`$ fits into a short exact sequence $`0I_FW_F\stackrel{d}{}0,`$ where $`I_F`$ is the inertia group of $`F`$. The local Langlands correspondence asserts that for all $`n1`$ there exists a natural bijection
$$\pi _F:𝒢^0(n)𝒜^0(n)$$
where $`𝒢^0(n)`$ is the set of equivalence classes of irreducible continuous $`n`$-dimensional complex representations of the Weil group $`W_F`$, and $`𝒜^0(n)`$ is the set of equivalence classes of irreducible supercuspidal representations of $`GL(n)`$.
We identify the elements of the set $`𝒢^0(1)`$, the quasicharacters of $`W_F`$, with quasicharacters of $`F^\times `$ via the reciprocity isomorphism $`r_F:W_F^{ab}F^\times `$. The local Langlands correspondence is compatible with twisting by quasicharacters, see \[6, 1.2\].
A complex representation of the Weil-Deligne group $`W_F^{}`$ is a pair $`(\rho ,N)`$ consisting of a continuous representation $`\rho :W_FGL(V),dim_{}V=n`$, together with a nilpotent endomorphism $`NEndV`$ such that $`\rho (w)N\rho (w)^1=wN`$. A representation $`\rho ^{}=(\rho ,N)`$ is admissible if $`\rho `$ is semisimple. For any $`n1`$ the representation $`sp(n)`$ is defined by $`V=^n=e_0+\mathrm{}+e_{n1},\rho (w)e_i=w^ie_i`$ and $`Ne_i=e_{i+1}(0i<n1),Ne_{n1}=0`$.
Let $`𝒢(n)`$ be the set of equivalence classes of admissible $`n`$-dimensional representations of the Weil-Deligne group $`W_F^{}`$, and let $`𝒜(n)`$ be the set of equivalence classes of irreducible smooth representations of $`GL(n)`$. To any indecomposable representation $`\rho sp(r)`$ of $`W_F^{}`$ associate the essentially square-integrable representation $`Q(\mathrm{\Delta })`$ with $`\mathrm{\Delta }=\{\pi _F(\rho ),||\pi _F(\rho ),\mathrm{},||^{r1}\pi _F(\rho )\}`$. To any admissible representation $`\rho ^{}=(\rho _1sp(r_1))\mathrm{}(\rho _msp(r_m))`$ of $`W_F^{}`$ associate the Langlands quotient $`Q(\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_m)`$ where $`\mathrm{\Delta }_j`$ is the segment
$$\mathrm{\Delta }_j=\{\pi _F(\rho _j),\mathrm{},||^{r_j1}\pi _F(\rho _j)\}$$
This creates as in \[7, p. 381\] a natural bijection
$$\pi _F:𝒢(n)𝒜(n)$$
A quasicharacter $`\psi :W_F^\times `$ is unramified if $`\psi `$ is trivial on the inertia group $`I_F`$. In that case we have $`\psi (w)=z^{d(w)}`$ with $`z^\times `$. The group of unramified quasicharacters of $`W_F`$ is denoted $`\mathrm{\Psi }(W_F)`$. We have $`\mathrm{\Psi }(W_F)^\times `$.
Let now $`\rho ^{}=\rho _1sp(r_1)\mathrm{}\rho _msp(r_m)`$ be an admissible representation of $`W_F^{}`$. The set
$$\{\psi _1\rho _1sp(r_1)\mathrm{}\psi _m\rho _msp(r_m):\psi _1,\mathrm{},\psi _m\mathrm{\Psi }(W_F)\}$$
will be called the orbit of $`\rho ^{}`$ under the action of $`(\mathrm{\Psi }(W_F))^m=\mathrm{\Psi }(W_F)\times \mathrm{}\times \mathrm{\Psi }(W_F)`$ ( $`m`$ factors ). Then $`𝒢(n)`$ admits a partition into orbits. Since $`(\mathrm{\Psi }(W_F))^m=(^\times )^m`$, a complex torus, each orbit has the structure of a complex algebraic variety. In this way, the set of admissible $`n`$-dimensional representations of the Weil-Deligne group acquires the structure of a complex algebraic variety. Each component in this variety is the quotient of a complex torus by a product of symmetric groups.
Let $`\mathrm{\Omega }G`$ be the Bernstein variety of $`G`$. Each point in $`\mathrm{\Omega }G`$ is a conjugacy class of cuspidal pairs $`(M,\sigma )`$. A quasicharacter $`\psi :M^\times `$ is unramified if $`\psi `$ is trivial on $`M^{}`$. The group of unramified quasicharacters of $`M`$ is denoted $`\mathrm{\Psi }(M)`$. We have $`\mathrm{\Psi }(M)(^\times )^{\mathrm{}}`$ where $`\mathrm{}`$ is the parabolic rank of $`M`$. The group $`\mathrm{\Psi }(M)`$ now creates orbits: the orbit of $`(M,\sigma )`$ is $`\{(M,\psi \sigma ):\psi \mathrm{\Psi }(M)\}`$. Denote this orbit by $`D`$, and set $`\mathrm{\Omega }=D/W(M,D)`$, where $`W(M)`$ is the Weyl group of $`M`$ and $`W(M,D)`$ is the subgroup of $`W(M)`$ which leaves $`D`$ globally invariant. The orbit $`D`$ has the structure of a complex torus, and so $`\mathrm{\Omega }`$ is a complex algebraic variety. We view $`\mathrm{\Omega }`$ as a component in the algebraic variety $`\mathrm{\Omega }G`$.
We recall the extended quotient. Let the finite group $`\mathrm{\Gamma }`$ act on the space $`X`$. Let $`\widehat{X}=\{(x,\gamma ):\gamma x=x\}`$, let $`\mathrm{\Gamma }`$ act on $`\widehat{X}`$ by $`\gamma _1(x,\gamma )=(\gamma _1x,\gamma _1\gamma \gamma _1^1)`$. Then $`\widehat{X}/\mathrm{\Gamma }`$ is the extended quotient of $`X`$ by $`\mathrm{\Gamma }`$. There is a canonical projection $`\widehat{X}/\mathrm{\Gamma }X/\mathrm{\Gamma }`$.
The Bernstein variety $`\mathrm{\Omega }G`$ is the disjoint union of ordinary quotients. We now replace the ordinary quotient by the extended quotient to create a new variety $`\mathrm{\Omega }^+G`$. So we have
$$\mathrm{\Omega }G=D/W(M,D)\text{and}\mathrm{\Omega }^+G=\widehat{D}/W(M,D)$$
Let $`\rho _1,\mathrm{},\rho _m`$ be irreducible complex representations of $`W_F`$ with $`dim_{}\rho _j=a_j`$ and let
$$\mathrm{\Delta }_j=\{\pi _F(\rho _j),\mathrm{},||^{r_j1}\pi _F(\rho _j)\}\text{and}R(\mathrm{\Delta }_j)=\pi _F(\rho _j)\mathrm{}||^{r_j1}\pi _F(\rho _j)$$
###### Theorem 1
Let $`\pi _F`$ be the local Langlands correspondence and let $`inf.ch.`$ be the infinitesimal character of Bernstein \[3, (III.4) p.75\]. Then we have a commutative diagram
$$\begin{array}{ccc}𝒢(n)& \stackrel{\pi _F}{}& 𝒜(n)\\ \alpha & & inf.ch& & \\ \mathrm{\Omega }^+G& \underset{\beta }{}& \mathrm{\Omega }G\end{array}$$
in which the maps $`\alpha `$ and $`\beta `$ are as follows
$$\begin{array}{ccc}\rho _1sp(r_1)\mathrm{}\rho _msp(r_m)& \stackrel{\pi _F}{}& Q(\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_m)\\ \alpha & & inf.ch.& & \\ (GL(a_j)^{r_j},\pi _F(\rho _j)^{r_j})& \underset{\beta }{}& (GL(a_j)^{r_j},R(\mathrm{\Delta }_j))\end{array}$$
The map $`\beta :\mathrm{\Omega }^+G\mathrm{\Omega }G`$ is not the canonical projection, but a twisted projection. The map $`\beta `$ is a morphism of complex algebraic varieties, injective on each component of $`\mathrm{\Omega }^+G`$.
Let $`𝒢_\mathrm{\Omega }(n)=(\beta \alpha )^1\mathrm{\Omega }`$. In the Bernstein decomposition $`𝒢(n)=𝒢_\mathrm{\Omega }(n)`$ we have
$$𝒢_\mathrm{\Omega }(n)=\mathrm{\Omega }^+$$
where we write $`\mathrm{\Omega }^+=\widehat{D}/W(M,D)`$ when $`\mathrm{\Omega }=D/W(M,D)`$.
Proof. We have to delve into some combinatorics. Let $`dim_{}\rho =a`$, let $`(M,\sigma )=(GL(a)^r,\pi _F(\rho )^r)`$, let $`D`$ be the orbit of $`(M,\sigma )`$. Then $`W(M,D)=S_r`$ the symmetric group on $`r`$ letters, and $`\mathrm{\Omega }=^{\times r}/S_r=Sym^r^\times `$, the $`r`$-fold symmetric product of $`^\times `$.
Write $`\mathrm{\Gamma }=S_r`$. Let $`t_1,\mathrm{},t_k`$ be a partition of $`r`$ in which $`t_1<\mathrm{}<t_k`$ and $`t_j`$ is repeated $`n_j`$ times, $`1jk`$. Then we have $`n_1t_1+\mathrm{}+n_kt_k=r`$. Let $`\gamma `$ be the permutation of $`r`$ letters which corresponds to this partition, so that $`\gamma `$ is the product of cycles $`(1,\mathrm{},t_1)\mathrm{}(1,\mathrm{},t_k)`$.
Let
$$\rho _\gamma ^{}=\rho sp(t_1)\mathrm{}\rho sp(t_k)$$
The first summand is repeated $`n_1`$ times, …, the last summand is repeated $`n_k`$ times. The orbit of $`\rho _\gamma ^{}`$ has the structure $`Sym^{n_1}^\times \times \mathrm{}\times Sym^{n_k}^\times `$. Note that $`\beta (\alpha (\rho _\gamma ^{}))\mathrm{\Omega }`$. Consider the orbit of $`\rho _\gamma ^{}`$. Each $`\rho ^{\prime \prime }`$ in the orbit of $`\rho _\gamma ^{}`$ is such that $`\alpha (\rho ^{\prime \prime })`$ is fixed by $`\gamma `$, so that $`\alpha (\rho ^{\prime \prime })D^\gamma `$.
Now the centralizer $`Z(\gamma )=(/t_1S_{n_1})\times \mathrm{}\times (/t_kS_{n_k})`$. But each cyclic group $`/t_1,\mathrm{},/t_k`$ acts trivially on $`D^\gamma `$. Hence $`D^\gamma /Z(\gamma )=D^\gamma /(S_{n_1}\times \mathrm{}\times S_{n_k})Sym^{n_1}^\times \times \mathrm{}\times Sym^{n_k}^\times `$.
Let $`𝒪(\rho _\gamma ^{})`$ be the orbit of $`\rho _\gamma ^{}`$. When we restrict $`\alpha `$ to this orbit we get
$$𝒪(\rho _\gamma ^{})D^\gamma /Z(\gamma )$$
Choose one $`\gamma `$ in each $`\mathrm{\Gamma }`$-conjugacy class. We get
$$𝒢_\mathrm{\Omega }(n)=𝒪(\rho _\gamma ^{})=D^\gamma /Z(\gamma )=\mathrm{\Omega }^+$$
If we start with $`\rho _1sp(r_1)\mathrm{}\rho _msp(r_m)`$ where $`\rho _1,\mathrm{},\rho _m`$ are not all equivalent (after unramified twist), then we partition this sum in an obvious way and apply the above argument piecewise. To this end, we note the following decompositions. Let $`X=X_1\times X_2`$, $`\mathrm{\Gamma }=\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$. Then we have
$$\widehat{X}/\mathrm{\Gamma }=\widehat{X_1}/\mathrm{\Gamma }_1\times \widehat{X_2}/\mathrm{\Gamma }_2\text{and}𝒜_{\mathrm{\Omega }_1\times \mathrm{\Omega }_2}(n)=𝒜_{\mathrm{\Omega }_1}(n)\times 𝒜_{\mathrm{\Omega }_2}(n)$$
where $`𝒜_\mathrm{\Omega }(n)=(inf.ch.)^1\mathrm{\Omega }`$ and $`\mathrm{\Omega }_1,\mathrm{\Omega }_2`$ are disjoint.
Example. Let $`T`$ be the diagonal subgroup of $`G=GL(2)`$ and let $`\mathrm{\Omega }`$ be the component in $`\mathrm{\Omega }G`$ containing the cuspidal pair $`(T,1)`$. Then $`\sigma 𝒜(2)`$ is arithmetically unramified if $`inf.ch.\sigma \mathrm{\Omega }`$. If $`\pi _F(\rho ^{})=\sigma `$ then $`\rho ^{}`$ is a $`2`$-dimensional representation of $`W_F^{}`$ and there are two possibilities:
$`\rho ^{}`$ is reducible, $`\rho ^{}=\psi _1\psi _2`$ with $`\psi _1,\psi _2`$ unramified quasicharacters of $`W_F`$. So $`\psi _j(w)=z_j^{d(w)},z_j^\times ,j=1,2`$. We have $`\pi _F(\rho ^{})=Q(\psi _1,\psi _2)`$ where $`\psi _1`$ does not precede $`\psi _2`$. In particular we obtain the $`1`$-dimensional representations of $`G`$ as follows:
$$\pi _F(||^{1/2}\psi ||^{1/2}\psi )=Q(||^{1/2}\psi ,||^{1/2}\psi )=\psi det$$
$`\rho ^{}`$ is indecomposable, $`\rho ^{}=\psi sp(2)`$. Then $`\pi _F(\rho ^{})=Q(\mathrm{\Delta }_1)`$ with $`\mathrm{\Delta }_1=\{\psi ,||\psi \}`$. In particular we have $`\pi _F(||^{1/2}\psi sp(2))=Q(\mathrm{\Delta }_2)`$ with $`\mathrm{\Delta }_2=\{||^{1/2}\psi ,||^{1/2}\psi \}`$ so $`\pi _F(||^{1/2}\psi sp(2))=\psi St(2)`$ where $`St(2)`$ is the Steinberg representation of $`G`$.
The orbit of $`(T,1)`$ is $`D=(^\times )^2`$, and $`W(T,D)=/2`$. Then $`\mathrm{\Omega }(^\times )^2//2Sym^2^\times `$. The extended quotient is $`\mathrm{\Omega }^+=Sym^2^\times ^\times `$. The twisted projection $`\beta `$ sends $`\{z_1,z_2\}`$ to $`\{z_1,z_2\}`$ and $`z`$ to $`\{z,q^1z\}`$. Note the twist by $`q^1`$ where $`q`$ is the cardinality of the residue field of $`F`$.
2. A geometric model for the Hecke algebra $`(G)`$
Let $`\rho _1,\mathrm{},\rho _m`$ be irreducible representations of $`W_F`$ such that $`det\rho _j`$ is unitary, $`1jm`$, and let $`\rho ^{}=\rho _1sp(r_1)\mathrm{}\rho _msp(r_m)`$. Let $`\psi _j\mathrm{\Psi }(W_F)`$ and give each $`\rho _j`$ an unramified twist as follows:
$$\rho ^{\prime \prime }=\psi _1||^{(1r_1)/2}\rho _1sp(r_1)\mathrm{}\psi _m||^{(1r_m)/2}\rho _msp(r_m)$$
Then $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ are in the same component of $`𝒢(n)`$.
###### Theorem 2
The map defined by
$$\underset{j}{}(\psi _j||^{(1r_j)/2}\rho _jsp(r_j))\underset{j}{}(\psi _j|\psi _j|^1||^{(1r_j)/2}\rho _jsp(r_j))$$
determines a deformation retraction of $`𝒜(n)`$ onto the tempered dual $`𝒜^t(n)`$.
Proof. This depends on the Langlands classification of tempered representations of $`GL(n)`$, as in \[7, Prop. 2.2.1\] and \[7, p. 384\].
###### Theorem 3
The periodic cyclic homology of $`(G)`$ is isomorphic to the periodised de Rham cohomology supported on finitely many components of the smooth dual of $`G`$.
Proof. We have the Bernstein decomposition for the Hecke algebra $`(G)`$: $`(G)(\mathrm{\Omega }).`$ We now use \[2, Theorem 7.8\] and note that de Rham cohomology is invariant under deformation retraction to get
$$HP_0((\mathrm{\Omega }))H^{ev}(\mathrm{\Omega }^+;)\text{and}HP_1((\mathrm{\Omega }))H^{odd}(\mathrm{\Omega }^+;)$$
where the variety $`\mathrm{\Omega }^+`$ is given its classical topology. It now follows from Theorem 2 that
$$HP_0((\mathrm{\Omega }))H^{ev}(𝒢_\mathrm{\Omega }(n);)\text{and}HP_1((\mathrm{\Omega }))H^{odd}(𝒢_\mathrm{\Omega }(n);)$$
Since $`HH_j((\mathrm{\Omega }))=0`$ for all $`j>n`$, periodic cyclic homology commutes with direct limits (cf. \[4, Theorem 2\]) and the Theorem follows.
J.B.: School of Mathematical Sciences, University of Exeter, North Park Road, Exeter, EX4 4QE, U.K., brodzki@maths.ex.ac.uk.
R.J.P.: Department of Mathematics, University of Manchester, Manchester, M13 9PL, U.K., roger@ma.man.ac.uk.
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# The Cosmic Microwave Background for a Nearly Flat Compact Hyperbolic Universe
## 1 Introduction
The two crucial properties of the universe at large scales are its curvature and its topology. Both properties are encoded in the cosmic microwave background (CMB), see e. g. \[Hu and White, 1996, Cornish et al., 1998, Kamionkowski and Buchalter, 2000\], which is measured with ever increasing resolution. The detection of the first acoustic peak in the CMB angular power spectrum $`C_l`$ \[Knox and Page, 2000\] provides evidence for a flat, or nearly flat, cold dark matter universe with a non-vanishing cosmological constant and/or an extra exotic energy component. The recent Boomerang \[de Bernardis et al., 2000\] and MAXIMA-1 \[Hanany et al., 2000\] measurements yield evidence even for the second acoustic peak which is, however, less pronounced than expected by standard CMB models. These CMB scenarios are based on isentropic initial perturbations in a universe composed of radiation, baryonic matter according to the Big Bang nucleosynthesis, cold dark matter, and a non-vanishing cosmological constant. Possible explanations for the low second peak and the surprisingly large scale of the first peak are discussed in \[White et al., 2000, Tegmark and Zaldarriaga, 2000, Cornish, 2000, Weinberg, 2000b\]. The constraints obtained from the MAXIMA-1 experiment are $`\mathrm{\Omega }_{\text{tot}}=0.90\pm 0.15`$, $`\mathrm{\Omega }_{\text{bar}}h_0^2=0.025\pm 0.010`$, $`\mathrm{\Omega }_{\text{cdm}}h_0^2=0.13\pm 0.10`$, and a spectral index $`n=0.99\pm 0.09`$ at 95% confidence level \[Balbi et al., 2000\]. Here $`\mathrm{\Omega }_{\text{cdm}}`$ and $`\mathrm{\Omega }_{\text{bar}}`$ denote the ratio of the cold dark matter (cdm) and baryonic (bar) energy densities, respectively, to the critical energy density.
The standard models describe the structure of the acoustic peaks, but fail to match with the low quadrupole moment $`C_2`$ of the COBE experiment \[Tegmark and Hamilton, 1997\]. This can be interpreted as a hint for a non-trivial topology of our universe. The standard models suppose a trivial topology implying a universe with infinite volume for negative and zero curvature. Models with a non-trivial topology lead to a finite volume and to a suppression in the angular power spectrum $`C_l`$ for low multipoles. This motivates the study of non-trivial topologies \[Lachièze-Rey and Luminet, 1995\] which leads to multiple images of a single source called topological lensing \[Uzan et al., 2000\], and the just mentioned suppression in the large scale CMB anisotropy on which we concentrate in this paper. For flat models the topological length scale is constrained to be significantly larger than half the diameter of the observable universe \[Levin et al., 1998\] which renders these models unattractive. (See, however, \[Roukema, 2000\].) Thus in the following we discuss models with negative curvature.
The computation of the CMB anisotropy for compact hyperbolic universes can be carried out in two different ways. On the one hand one can compute the fluctuations by using the so–called method of images which requires the group elements which define the fundamental cell of the considered non-trivial topology \[Bond et al., 1998, Bond et al., 2000a, Bond et al., 2000b\]. On the other hand one can use a method which requires the eigenmodes of the fundamental cell with respect to the Laplace-Beltrami operator of the considered space. With the latter method the CMB anisotropy is computed for hyperbolic universes with a vanishing cosmological constant $`\mathrm{\Lambda }`$ for a compact orbifold \[Aurich, 1999\] and several compact manifolds \[Inoue et al., 1999, Cornish and Spergel, 2000\].
Our main aim in this paper is to incorporate a non-vanishing cosmological constant $`\mathrm{\Lambda }`$ and, in addition, an extra smooth dark energy component $`\epsilon _\text{x}`$ \[Turner and White, 1997\] in the anisotropy calculations as suggested by current observations, where we concentrate on the lower part of the angular power spectrum $`C_l`$, which is affected by the non-trivial topology. Recent investigations, in particular of the luminosity-redshift surveys of Type Ia supernovae \[Riess et al., 1998, Perlmutter et al., 1999\], strongly indicate that current observations require apart from a matter density component $`\mathrm{\Omega }_{\text{mat}}=\mathrm{\Omega }_{\text{cdm}}+\mathrm{\Omega }_{\text{bar}}0.3\pm 0.1`$ an additional unclustered, dark energy component of 60% of the total energy density of the universe with negative pressure \[Wang et al., 2000, Hu et al., 2000\] corresponding to an accelerated expansion. In the following we will assume that the new component is a mixture of vacuum energy, $`\mathrm{\Omega }_{\text{vac}}`$, or cosmological constant $`\mathrm{\Lambda }`$, and a smooth dark energy component $`\mathrm{\Omega }_\text{x}`$. Whereas a cosmological constant $`\mathrm{\Lambda }>0`$ corresponds to a constant homogeneous energy component $`\epsilon _{\text{vac}}>0`$ with negative pressure and equation of state $`w_{\text{vac}}=1`$, the extra dark energy component $`X`$ considered by us consists of a dynamical, time-dependent energy density $`\epsilon _\text{x}>0`$ with negative pressure and equation of state $`w_\text{x}=\frac{2}{3}`$. (Here $`w`$ denotes the ratio of pressure to energy density.) The $`X`$-component is similar to a particular version of quintessence which is generated by a slowly evolving scalar field with an exponential or inverse power law potential and $`1<w_{\text{quint}}0`$ \[Ratra and Peebles, 1988, Peebles and Ratra, 1988, Wetterich, 1988a, Wetterich, 1988b, Caldwell et al., 1998\]. The quintessence models and the more recent “tracker field” models \[Zlatev et al., 1999\] have been introduced to solve the two cosmological constant problems \[Weinberg, 2000a\]: i) why is $`\epsilon _{\text{vac}}`$ so small, ii) why is $`\epsilon _{\text{vac}}`$ not only small, but also of the same magnitude as the present mass of the universe? (See also \[Armendariz-Picon et al., 2000a, Armendariz-Picon et al., 2000b\] for the concept of “k-essence” and \[Weinberg, 2000a\] for an anthropic explanation.) While at present there is no direct evidence for a scalar field, let alone for a particular form of the potential, observations are consistent with an $`X`$-component if $`w_\text{x}=0.65\pm 0.07`$ (assuming flat models) \[Wang et al., 2000\]. However, since the concept of a scalar quintessence field is a purely classical phenomenological one, we expect the quantum mechanics of such a theory to be plagued with the usual problem of nonrenormalizability \[Ratra and Peebles, 1988, Peebles and Ratra, 1988\]. We therefore adopt in this paper the point of view of an effective quintessence model, where we do not start from a given potential for the scalar field, but rather give the redshift behavior $`\epsilon _\text{x}a^1`$ and thus the equation of state $`w_\text{x}=\frac{2}{3}`$, which is consistent with the above cited experimental bounds ($`a`$ is the cosmic scale factor, see the next section). Furthermore, we will assume that the $`X`$-component is spatially constant which can be understood as follows \[Ratra and Peebles, 1988, Peebles and Ratra, 1988\]. In linear perturbation theory spatial gradients in the scalar field act like particles with very low mass that cannot bind to a nonrelativistic gravitational potential well that is much smaller than the Hubble length $`H^1`$. Thus dynamical studies of groups and clusters of galaxies with size $`H^1`$ cannot detect concentrations in $`\epsilon _\text{x}`$, and thus $`\epsilon _\text{x}`$ can be assumed to be spatially constant. This is confirmed by recent determinations of $`\mathrm{\Omega }_{\text{mat}}`$ obtained from relative velocities of galaxies yielding results in the range $`0.2\mathrm{\Omega }_{\text{mat}}0.4`$ \[Juszkiewicz et al., 2000, Willick, 2000\]. Since these measurements are sensitive to a spatially inhomogeneous dark energy component \[Lahav et al., 1991\], we have to assume the dark energy component, required by CMB measurements, to be smooth. Therefore, we assume below an $`X`$-component being spatially homogeneous like the vacuum energy and $`\mathrm{\Omega }_{\text{mat}}=0.3`$. (For further arguments concerning the smoothness assumption, see \[Turner and White, 1997, White, 1998\].)
The fundamental cell which we consider in this paper is the same pentahedron as considered in \[Aurich, 1999\] with the same Dirichlet eigenmode spectrum. (For more details, see also \[Aurich and Marklof, 1996\].) One is interested in fundamental cells with volumes as small as possible for several reasons, e. g. the creation probability of the universe increases dramatically with decreasing volume \[Atkatz and Pagels, 1982\] and, furthermore, to obtain an appreciable effect in the CMB anisotropy. Unfortunately, the smallest hyperbolic manifold $``$ is unknown. The smallest one known has a volume $`\text{vol}()=0.94272\mathrm{}R^3`$, which is far above the lower bound $`\text{vol}()>0.16668\mathrm{}R^3`$ \[Gabai et al., 1996\]. Here $`R`$ denotes the curvature radius of the universal covering space.
Concerning the impact on the CMB anisotropy the main difference with respect to infinite-volume models arises from the discrete eigenvalue spectrum $`\{E_n\}`$ which compact manifolds possess in contrast to the continuous spectrum of the infinite volume models. The volume of the manifold has an important influence on the CMB anisotropy because it determines the magnitude of the lowest eigenmode and the number of eigenmodes below a given value $`E`$. For a manifold $``$, the number $`𝒩(E)`$ of eigenvalues below $`E`$ is asymptotically given by Weyl’s law $`(R=1)`$
$$𝒩(E)\frac{\text{vol}()}{6\pi ^2}k^3\text{with}k:=\sqrt{E1}.$$
The first lowest eigenmodes determine the largest scales of the CMB anisotropies. The smaller the volume the stronger is the suppression of the first multipoles in the angular power spectrum $`C_l`$.
The considered pentahedron has a volume $`\text{vol}()0.7173068R^3`$. Since only one of two symmetry classes is taken into account, the following computations correspond to a hyperbolic manifold with $`\text{vol}()0.3586534R^3`$. However, a volume comparison is complicated by the fact that Weyl’s law has additional terms for orbifolds in comparison to manifolds, in particular the surface term \[Aurich and Marklof, 1996\] being absent in the case of manifolds. The additional surface term leads to a suppression of $`𝒩(E)`$ in comparison with manifolds as shown in figure 1 for $`E<3026`$. The figure demonstrates that the considered orbifold mimics a manifold with effective volume $`0.25R^3`$. The statistical properties of the eigenmodes are expected to be of the same random nature as observed in quantum chaos \[Aurich and Steiner, 1993, Inoue, 1999\].
## 2 The background model
The standard cosmological model based on the Friedmann-Lemaître-Robertson-Walker metric $`(c=1)`$
$$ds^2=a^2(\eta )\left\{d\eta ^2\gamma _{ij}dx^idx^j\right\}$$
is governed for negative curvature ($`K=1`$) by the Friedmann equation
$$a_{}^{}{}_{}{}^{2}a^2=\frac{8\pi G}{3}T_0^0a^4,$$
where $`a(\eta )`$ is the cosmic scale factor and $`\eta `$ the conformal time. The prime denotes differentiation with respect to $`\eta `$. The energy-momentum tensor for an ideal fluid is given by
$$T_\nu ^\mu =(\epsilon +p)u^\mu u_\nu p\delta _\nu ^\mu ,$$
where $`u^\mu `$ is the four-velocity of the fluid, and $`\epsilon =\epsilon (\eta )`$ denotes the energy density and $`p=p(\eta )`$ the pressure. In the following we consider multi-component models containing a matter-energy density $`\epsilon _{\text{mat}}`$, a radiation density $`\epsilon _{\text{rad}}`$ as well as a non-vanishing cosmological constant $`\mathrm{\Lambda }=8\pi G\epsilon _{\text{vac}}`$ and a spatially constant $`X`$-component $`\epsilon _\text{x}`$ with an equation of state $`p_\text{x}=\frac{2}{3}\epsilon _\text{x}`$. Then the $`00`$component of the energy-momentum tensor is given in comoving coordinates by
$$T_0^0=\underset{\begin{array}{c}k=0\\ k2\end{array}}{\overset{4}{}}\epsilon _{k,0}\left(\frac{a_0}{a}\right)^k,$$
expressed in terms of the current radiation density $`\epsilon _{4,0}=\epsilon _{\text{rad},0}`$, the current matter density $`\epsilon _{3,0}=\epsilon _{\text{mat},0}`$, the current $`X`$-component energy density $`\epsilon _{1,0}=\epsilon _{\text{x},0}`$ and a vacuum energy density $`\epsilon _{0,0}=\epsilon _{\text{vac}}`$. Here $`a_0:=a(\eta _0)`$ is the scale factor of the present epoch. The present conformal time $`\eta _0`$ is implicitly given by
$$a(\eta _0)=\frac{1}{H_0\sqrt{1\mathrm{\Omega }_{\text{tot}}}},\mathrm{\Omega }_{\text{tot}}=\mathrm{\Omega }_{\text{rad}}+\mathrm{\Omega }_{\text{mat}}+\mathrm{\Omega }_\text{x}+\mathrm{\Omega }_{\text{vac}},$$
where $`H_0=h_0100\text{ km s}^1\text{Mpc}^1`$ denotes Hubble’s constant and $`\mathrm{\Omega }_k:=\epsilon _{k,0}/\epsilon _{\text{crit}}`$ with $`\epsilon _{\text{crit}}=3H_0^2/(8\pi G)`$. With
$$\mathrm{\Omega }_2:=\mathrm{\Omega }_{\text{curv}}:=\frac{K}{(a_0H_0)^2}=\frac{1}{(a_0H_0)^2}=1\mathrm{\Omega }_{\text{tot}}>0$$
the Friedmann equation reads
$$a^{}(\eta )=H_0\sqrt{\underset{k=0}{\overset{4}{}}\mathrm{\Omega }_ka_0^ka^{4k}}.$$
(1)
This gives the infinitely far future $`\eta _{\mathrm{}}`$ as
$$\eta _{\mathrm{}}=\sqrt{1\mathrm{\Omega }_{\text{tot}}}_0^{\mathrm{}}\frac{dx}{\sqrt{_{k=0}^4\mathrm{\Omega }_kx^k}},$$
(2)
which yields $`\eta _{\mathrm{}}<\mathrm{}`$ for a large class of models, see below. Notice that the various components redshift like $`a^k`$ with an associated equation of state $`w_k:=p_k/\epsilon _k=(k3)/3`$. The current value of the deceleration parameter is given by $`q_0=\mathrm{\Omega }_{\text{rad}}+\frac{1}{2}\mathrm{\Omega }_{\text{mat}}\frac{1}{2}\mathrm{\Omega }_\text{x}\mathrm{\Omega }_{\text{vac}}`$.
Let us define the following quantities
$$A:=\frac{1}{2}\mathrm{\Omega }_{\text{mat}}H_0^2a_0^3=\frac{2a_{\text{eq}}}{\widehat{\eta }^2},$$
$$B:=\frac{1}{4}\mathrm{\Omega }_\text{x}H_0^2a_0,$$
$$C:=\frac{1}{12}A^2\widehat{\eta }^2\mathrm{\Lambda }$$
with
$$\widehat{\eta }:=\frac{2\sqrt{\mathrm{\Omega }_{\text{rad}}}}{H_0a_0\mathrm{\Omega }_{\text{mat}}}\left(1+\sqrt{2}\right)\eta _{\text{eq}},$$
where the subscript “eq” marks the epoch of matter-radiation equality, and $`a_{\text{eq}}:=a(\eta _{\text{eq}})=a_0(\mathrm{\Omega }_{\text{rad}}/\mathrm{\Omega }_{\text{mat}})`$. With the initial conditions $`a(0)=0`$ and $`a^{}(0)>0`$, equation (1) has the unique solution
$$a(\eta )=\frac{A}{2}\frac{𝒫(\eta )\frac{1}{12}\widehat{\eta }𝒫^{}(\eta )+AB\widehat{\eta }^2}{(𝒫(\eta )\frac{1}{12})^2C}.$$
(3)
Here $`𝒫(\eta )`$ denotes the Weierstrass $`𝒫`$-function which can numerically be evaluated very efficiently by
$$𝒫(\eta )=𝒫(\eta ;g_2,g_3)=\frac{1}{\eta ^2}+\underset{k=2}{\overset{\mathrm{}}{}}c_k\eta ^{2k2}$$
(4)
with
$$c_2:=\frac{g_2}{20},c_3:=\frac{g_3}{28}$$
and \[Abramowitz and Stegun, 1965\]
$$c_k=\frac{3}{(2k+1)(k3)}\underset{m=2}{\overset{k2}{}}c_mc_{km}\text{ for }k4.$$
The so-called invariants $`g_2`$ and $`g_3`$ are determined by the cosmological parameters
$`g_2`$ $`=`$ $`{\displaystyle \frac{1}{12}}+4C2AB`$
$`g_3`$ $`=`$ $`{\displaystyle \frac{1}{216}}+{\displaystyle \frac{8CA^2\mathrm{\Lambda }}{12}}+{\displaystyle \frac{AB}{6}}A^2B^2\widehat{\eta }^2.`$
For cosmologically plausible parameter choices the series (4) needs only to be evaluated by taking into account the first twenty terms and thus the explicit solution (3) is much more efficient than the usual integration of the Friedmann equation.
Expanding (3) in a series at $`\eta =0`$ gives the scale factor at early times
$`a(\eta )`$ $`=`$ $`A\{\widehat{\eta }\eta +{\displaystyle \frac{\eta ^2}{2}}+\widehat{\eta }{\displaystyle \frac{\eta ^3}{3!}}+(1+12AB\widehat{\eta }^2){\displaystyle \frac{\eta ^4}{4!}}`$
$`+\widehat{\eta }(1+36AB+48C){\displaystyle \frac{\eta ^5}{5!}}+O\left(\eta ^6\right)\}.`$
This expansion shows that the $`X`$-component term $`B`$ influences the scale factor one power in $`\eta `$ lower than the cosmological constant term $`C`$.
In the case $`\mathrm{\Lambda }>0`$ and/or $`B>0`$ the conformal time $`\eta `$ is restricted to $`0\eta <\eta _{\mathrm{}}<\mathrm{}`$, where $`\eta _{\mathrm{}}`$, defined in (2), is obtained from the implicit relation
$$𝒫(\eta _{\mathrm{}};g_2,g_3)=\frac{1}{12}+\sqrt{C}.$$
In the case $`\mathrm{\Lambda }>0`$ the scale factor $`a(\eta )`$ has a simple pole at $`\eta =\eta _{\mathrm{}}`$ with residue $`\sqrt{3/\mathrm{\Lambda }}`$, i. e.
$$a(\eta )=\frac{\sqrt{3/\mathrm{\Lambda }}}{\eta _{\mathrm{}}\eta }\frac{3}{\mathrm{\Lambda }}B+O(\eta \eta _{\mathrm{}}),$$
which leads to an exponential expansion of the universe with scale factor $`R(t)=a(\eta (t))=O(\mathrm{exp}(\sqrt{\mathrm{\Lambda }/3}t))`$ for cosmic time $`t\mathrm{}`$.
For $`\mathrm{\Lambda }=0`$ and $`\mathrm{\Omega }_\text{x}>0`$ one has at $`\eta =\eta _{\mathrm{}}`$ a double pole, i. e.
$$a(\eta )=\frac{1/B}{(\eta \eta _{\mathrm{}})^2}\frac{1}{12B}+O((\eta \eta _{\mathrm{}})^2),$$
leading to $`R(t)=Bt^2+\mathrm{}`$ for $`t\mathrm{}`$, which follows from the exact formula
$$t=\frac{1}{12B}\eta (t)+\frac{1}{B}\zeta (\eta _{\mathrm{}}\eta (t))\frac{1}{B}\zeta (\eta _{\mathrm{}}),$$
where $`\zeta (\eta ):=\zeta (\eta ;g_2,g_3)`$ denotes the Weierstrass zeta function. In this special case formula (3) reduces to ($`𝒫(\eta _{\mathrm{}})=\frac{1}{12}`$)
$$a(\eta )=\frac{1}{B}𝒫(\eta \eta _{\mathrm{}})\frac{1}{12B}\text{ with }0\eta <\eta _{\mathrm{}}.$$
Finally, if both the cosmological constant and the $`X`$-component vanish, $`B=C=0`$, the invariants simplify to $`g_2=\frac{1}{12}`$ and $`g_3=\frac{1}{216}`$, and the $`𝒫`$-function can be expressed in terms of an elementary function \[Abramowitz and Stegun, 1965\]
$$𝒫(\eta ;\frac{1}{12},\frac{1}{216})=\frac{1}{12}+\frac{1}{4\mathrm{sinh}^2\frac{\eta }{2}},$$
which, with (3), leads to $`a(\eta )=A(\widehat{\eta }\mathrm{sinh}\eta +\mathrm{cosh}\eta 1),`$ which is the well-known expression for the scale factor of a two-component model consisting of radiation and matter only. This leads to $`R(t)=t+A\mathrm{ln}t+O(1)`$ for $`t\mathrm{}`$.
As an illustration we show in figure 2 the cosmic scale factor $`a(\eta )`$ for three different, nearly flat models ($`\mathrm{\Omega }_{\text{tot}}=0.9`$). The full curve corresponds to a two-component model consisting of radiation and matter ($`\mathrm{\Omega }_{\text{mat}}=0.9`$) only. For this model $`\eta _{\mathrm{}}=\mathrm{}`$ holds. The dashed curve represents a three-component model consisting of radiation, matter and an $`X`$-component ($`\mathrm{\Omega }_\text{x}=0.6`$). The approach to the double pole at $`\eta _{\mathrm{}}=1.739`$ is clearly visible. The dotted curve shows $`a(\eta )`$ for a three-component model consisting of radiation, matter and vacuum energy ($`\mathrm{\Omega }_{\text{vac}}=0.6`$). One observes a steep rise to the pole at $`\eta _{\mathrm{}}=1.385`$. In addition we have indicated the present scale factor $`a_0`$ by a dot.
In table 1 the scale factor $`a_0`$ and the present age of the universe $`t_0`$ as well as several cosmologically important times, i. e. $`\eta _{\text{eq}}`$, $`\eta _{\text{SLS}}`$, $`\eta _0`$ and $`\eta _{\mathrm{}}`$, are given for several combinations of $`\mathrm{\Omega }_k`$. The present age of the universe $`t_0`$ is very close to the limit given by globular cluster ages $`13.5\pm 2.0`$ Gyr \[Jimenez, 1999, Primack, 2000\]. White dwarf cooling rates lead to an age of our galaxy of $`9.3\pm 2.0`$ Gyr \[Winget et al., 1987\] or $`8.0\pm 1.5`$ Gyr \[Leggett et al., 1998\]. For a smaller Hubble constant, e. g. $`h_0=0.6`$, the age $`t_0`$ obtains larger values with 13.18 Gyr $`t_0`$ 15.24 Gyr.
## 3 The CMB anisotropy
In the following we consider only scalar perturbations and their influence on the CMB. Furthermore, we assume that the vacuum energy $`\epsilon _{\text{vac}}`$ and the other dark energy component $`\epsilon _\text{x}`$ are spatially constant. The metric with scalar perturbations is written in the conformal-Newtonian gauge in terms of scalar functions $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ as
$$ds^2=a^2(\eta )\left\{(1+2\mathrm{\Phi })d\eta ^2(12\mathrm{\Psi })\gamma _{ij}dx^idx^j\right\},$$
where $`\mathrm{\Phi }=\mathrm{\Psi }`$ for a diagonal $`T_{\mu \nu }`$. Assuming negligible entropy perturbations $`\delta S=0`$, the evolution of the metric perturbation $`\mathrm{\Phi }`$ gives in first-order perturbation theory in the conformal-Newtonian gauge \[Mukhanov et al., 1992\]
$$\mathrm{\Phi }^{\prime \prime }+3\widehat{H}(1+c_s^2)\mathrm{\Phi }^{}c_s^2\mathrm{\Delta }\mathrm{\Phi }$$
$$+\{2\widehat{H}^{}+(1+3c_s^2)(\widehat{H}^2+1)\}\mathrm{\Phi }=\mathrm{\hspace{0.33em}0},$$
where $`\widehat{H}:=a^{}/a`$ and $`\mathrm{\Delta }`$ denotes the Laplace-Beltrami operator. The quantity $`c_s^2`$ can be interpreted as the sound velocity and is given by
$$c_s^2=\frac{1}{3+\frac{9}{4}\epsilon _{\text{mat}}/\epsilon _{\text{rad}}}.$$
Here it is assumed that the vacuum energy $`\epsilon _{\text{vac}}`$ and the other dark energy component $`\epsilon _\text{x}`$ are spatially constant.
Specifying $`\mathrm{\Phi }`$ at $`\eta =0`$ such that it corresponds to a scale-invariant (Harrison-Zel’dovich) spectrum, allows the computation of the time-evolution of the metric perturbation $`\mathrm{\Phi }`$. This in turn gives the input to the Sachs-Wolfe formula \[Sachs and Wolfe, 1967\] which reads for isentropic initial conditions
$`{\displaystyle \frac{\delta T}{T}}`$ $`=`$ $`2\mathrm{\Phi }(\eta _{\text{SLS}},\stackrel{}{x}(\eta _{\text{SLS}})){\displaystyle \frac{3}{2}}\mathrm{\Phi }(0,\stackrel{}{x}(0))`$ (5)
$`+2{\displaystyle _{\stackrel{}{x}(\eta _{\text{SLS}})}^{\stackrel{}{x}(\eta _0)}}𝑑\eta {\displaystyle \frac{\mathrm{\Phi }(\eta ,\stackrel{}{x}(\eta ))}{\eta }},`$
from which one obtains the desired temperature fluctuations $`\delta T`$ of the CMB. The conformal time at recombination, which defines the surface of last scattering, is denoted by $`\eta _{\text{SLS}}`$. For $`\eta _{\text{SLS}}\eta _{\text{eq}}`$ the first two terms on the right-hand side are approximately
$$2\mathrm{\Phi }(\eta _{\text{SLS}},\stackrel{}{x}(\eta _{\text{SLS}}))\frac{3}{2}\mathrm{\Phi }(0,\stackrel{}{x}(0))\frac{1}{3}\mathrm{\Phi }(\eta _{\text{SLS}},\stackrel{}{x}(\eta _{\text{SLS}})).$$
(6)
This is the so-called ordinary or naive Sachs-Wolfe term (NSW), whereas the other term in (5) is called integrated Sachs-Wolfe term (ISW).
The metric perturbation $`\mathrm{\Phi }`$ is expanded with respect to the eigenmodes
$$\mathrm{\Delta }\psi _n(\stackrel{}{x})=E_n\psi _n(\stackrel{}{x}),k_n:=\sqrt{E_n1},$$
of the considered compact orbifold, i. e.
$$\mathrm{\Phi }(\eta ,\stackrel{}{x})=\underset{n=1}{\overset{\mathrm{}}{}}f_n(\eta )\psi _n(\stackrel{}{x}),$$
which yields for $`f_n(\eta )`$ the differential equation
$`f_n^{\prime \prime }(\eta )`$ $`+`$ $`3\widehat{H}(1+c_s^2)f_n^{}(\eta )`$ (7)
$`+\{c_s^2E_n+2\widehat{H}^{}+(1+3c_s^2)(\widehat{H}^2+1)\}f_n(\eta )=\mathrm{\hspace{0.33em}0},`$
where $`\widehat{H}`$ and $`c_s^2`$ are determined by the background model. The initial conditions are ($`\alpha >0`$ is a normalization constant)
$$f_n(0)=\frac{\alpha }{\sqrt{k_n(k_n^2+1)}}\text{ and }f_n^{}(0)=\frac{f_n(0)}{8\widehat{\eta }},$$
(8)
which carry over to a Harrison-Zel’dovitch spectrum having a spectral index $`n=1`$ and selecting only the non-decaying modes. Using the eigenmodes the perturbation is defined obeying the periodicity condition imposed by the fundamental cell.
The time dependence of $`f_n(\eta )`$, determined by the background model (3), is obtained by numerical integration of (7) and is shown in figure 3 for three different models. The first model, shown in figure 3a), is completely dominated by matter, $`\mathrm{\Omega }_{\text{tot}}=\mathrm{\Omega }_{\text{mat}}=0.9`$, whereas the other two models belong to $`\mathrm{\Omega }_{\text{mat}}=0.3`$ and 0.6 for $`\mathrm{\Omega }_\text{x}`$ and $`\mathrm{\Omega }_{\text{vac}}`$, shown in b) and c), respectively. The latter two cases have a finite $`\eta _{\mathrm{}}`$ at which the perturbation vanishes, i. e. $`(\eta \eta _{\mathrm{}})`$
$$f_n(\eta )(\eta _{\mathrm{}}\eta )^{\frac{7\sqrt{17}}{2}}$$
for the case $`\mathrm{\Omega }_\text{x}>0`$ and $`\mathrm{\Omega }_{\text{vac}}=0`$, and
$$f_n(\eta )\eta _{\mathrm{}}\eta $$
for $`\mathrm{\Omega }_\text{x}=0`$ and $`\mathrm{\Omega }_{\text{vac}}>0`$. Furthermore, perturbation modes with wavelength $`\lambda =\frac{2\pi }{k}>\eta _{\mathrm{}}`$ will never enter the horizon in models with $`\mathrm{\Omega }_\text{x}>0`$ and/or $`\mathrm{\Omega }_{\text{vac}}>0`$. The first decline of $`f_n(\eta )/f_n(0)`$ from 1 to $`\frac{9}{10}`$ for small values of $`\eta `$ is due to the transition from the radiation- to the matter-dominated epoch (see, e. g. \[Mukhanov et al., 1992\]), which leads to the approximation (6).
With the background model (3), the time-evolution (7), and the Sachs-Wolfe formula (5), the angular power spectrum of the CMB anisotropy
$$C_l=\frac{1}{2l+1}\underset{m=l}{\overset{l}{}}|a_{lm}|^2$$
can be computed, where $`a_{lm}`$ are the expansion coefficients of the CMB anisotropy $`\delta T`$ with respect to the spherical harmonics $`Y_l^m(\theta ,\varphi )`$.
The angular power spectra $`\delta T_l:=\sqrt{l(l+1)C_l/2\pi }`$ are computed for several models. The considered compact orbifold as well as the position of the observer is the same as in \[Aurich, 1999\]. In figure 4 the angular power spectrum $`\delta T_l`$ is shown for the case $`\mathrm{\Omega }_{\text{bar}}=0.05`$, $`\mathrm{\Omega }_{\text{cdm}}=0.25`$, $`\mathrm{\Omega }_\text{x}=0.0`$, $`\mathrm{\Omega }_{\text{vac}}=0.6`$ and $`h_0=0.7`$. The curve is obtained from CMBFAST \[Zaldarriaga and Seljak, 1999\], i. e. is obtained for an infinitely extended hyperbolic universe, whereas the dots are obtained by the procedure outlined above, i. e. for a compact hyperbolic universe. Because the latter computation takes only the first 749 eigenmodes with $`k<55`$ into account, one observes at $`l40`$ a decline towards zero. This is solely due to the truncation in the $`k`$-summation because the eigenmodes are only computed up to this $`k`$ value. The considered modes are all above horizon at recombination, and thus the processes leading to the acoustic peak can be ignored. In the following we are only concerned with the low multipoles $`C_l`$ which are affected by the non-trivial topology. These low multipoles are not affected by truncating the $`k`$-summation. The figure shows clearly the suppression in power for $`l10`$.
Since $`\delta T_l`$ is computed for a fixed observer, i. e. it represents a “one-sky realization”, the spectrum is not smooth like the CMBFAST spectrum. It rather shows fluctuations which an experiment would observe which necessarily measures a one-sky realization.
The lower part of the angular power spectra $`\delta T_l`$, computed for several models with a vanishing and non-vanishing $`X`$-component, are shown in figures 5 and 6, respectively. The spectra in figures 5a) and 5b) for $`\mathrm{\Omega }_{\text{tot}}=0.5`$ and $`\mathrm{\Omega }_{\text{tot}}=0.8`$, respectively, show a plateau being normalized to $`30\mu \text{K}`$. At higher values of $`l`$ the angular power spectra $`\delta T_l`$ rise again (see figure 4) which is not shown here because our calculations do not take into account the necessary processes leading to the acoustic peak, since the modes considered here are well above the horizon at recombination. For low values of $`l`$ one observes a nearly linear increase of $`\delta T_l`$ which is caused by the finite size of the fundamental cell which in turn causes a cut-off in the $`k`$-spectrum. (The straight lines are drawn solely to guide the eyes.) One observes that the “bend” point, where the behavior turns from a linear increase to a plateau, decreases towards smaller values of $`l`$ for increasing vacuum energy. If the amount of vacuum energy is replaced by the same energy contribution of an $`X`$-component one obtains quantitatively analogous angular power spectra because the behavior of $`f_n(\eta )`$ shown in figure 3 is similar for vacuum energy and the $`X`$-component for $`\eta \eta _0`$. The two models shown in figures 5c) and 5d) possess an even larger $`\mathrm{\Omega }_{\text{tot}}`$, i. e. $`\mathrm{\Omega }_{\text{tot}}=0.9`$ and $`\mathrm{\Omega }_{\text{tot}}=0.95`$, respectively. Here the suppression is much less pronounced than in the cases with $`\mathrm{\Omega }_{\text{tot}}0.85`$. In the case $`\mathrm{\Omega }_{\text{tot}}=0.9`$ the quadrupole moment is larger than the other low multipoles which is due to the large integrated Sachs-Wolfe contribution (see below). In the other case $`\mathrm{\Omega }_{\text{tot}}=0.95`$ one observes very large fluctuations for low values of $`l`$.
For four models with an $`X`$-component the angular power spectra $`\delta T_l`$ are shown in figure 6. In figure 6a) and 6b) two models with $`\mathrm{\Omega }_{\text{tot}}=0.9`$ are shown, where in the first case the energy density is equally distributed between $`\mathrm{\Omega }_{\text{mat}}`$, $`\mathrm{\Omega }_\text{x}`$ and $`\mathrm{\Omega }_{\text{vac}}`$, and in the second case the $`X`$-component dominates. One observes similar angular power spectra which is again explained by the similar behavior of $`f_n(\eta )`$. In both cases the multipoles with $`l10`$ are suppressed. In figure 6c) and 6d) two models with a vanishing vacuum energy are shown for $`\mathrm{\Omega }_{\text{tot}}=0.9`$ and $`\mathrm{\Omega }_{\text{tot}}=0.95`$, respectively. In the latter case the suppression of low multipoles is blurred by very large fluctuations, which occur as in figure 5d).
The angular power spectrum does not go to zero at the smallest values of $`l`$. This is due to the competition of the two contributions to $`\delta T_l`$, i. e. the NSW and the ISW term in (5). As shown in figure 7 the NSW term gives a contribution which indeed vanishes for small values of $`l`$ because of the cut-off in the eigenmode spectrum. But the ISW term adds an almost constant contribution which even increases towards small values of $`l`$. This interplay is responsible for the fact that $`\delta T_l`$ does not fall below $`15\mu K`$, where a plateau of $`30\mu K`$ is assumed. The relative contribution of the two terms is largely determined by the chosen initial conditions at $`\eta =0`$.
The inflationary models naturally suggest isentropic initial conditions and these are imposed in the above calculations. However, imposing isocurvature initial conditions leads to a much smaller ISW contribution relative to the NSW term. This is shown in figure 7, where in figure 7a) isentropic initial conditions according to most inflationary models and in figure 7b) isocurvature initial conditions are chosen. If the observed increase in $`\delta T_l`$ would be approximately linear towards zero for nearly flat models, this would imply a small ISW contribution and this would point to isocurvature initial conditions.
To summarize the results, the anomalously low quadrupole moment obtained from the COBE measurements can be taken as a first sign for a universe with a finite volume. The presented calculations demonstrate that low multipoles occur for the considered compact fundamental domain even for nearly flat, but hyperbolic, models with $`\mathrm{\Omega }_{\text{tot}}0.9`$. For even larger values of $`\mathrm{\Omega }_{\text{tot}}0.95`$ very large fluctuations occur which may also be an indication for a finite volume. Furthermore, the kind of increase of $`\delta T_l`$ gives a clue to the initial conditions. Future experiments which survey the complete CMB sky like MAP and PLANCK, will have the required signal to noise ratio to reveal a possible finite universe.
### ACKNOWLEDGMENTS
We would like to thank the Rechenzentrum of the University of Karlsruhe for the access to their computers.
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# Entanglement versus Disentanglement: Quantum Cryptography
## Abstract
In quantum information, the role of entanglement and disentanglement is itself a subject of research and debate. Earlier works on quantum cryptography have almost established that entanglement has no special advantage in quantum cryptography. In this paper we reveal that entanglement is better ingredient than disentanglement for our alternative quantum cryptography.
In quantum information, there are some tasks which can be accomplished only by entanglement -such as dense coding and teleportation . But, there are some other tasks which can be realized both by entanglement and disentanglement. Quantum computation algorithm and quantum cryptography \[3-5\]are two major applications of both entanglement and disentanglement.
Between entanglement and disentanglement which one is better ingredient in quantum computation and quantum cryptography ? This question is yet not settled in quantum computation, however, entanglement enjoys some favoritism from the researchers in this field. In quantum cryptography this question is believed to have been settled. In his entanglement-based quantum key distribution protocol Ekert pointed out that quantum encryption can be executed after the transmission of the quantum state. This seemed to be advantageous to ensure the security of the key. Bennett, Brassard and Mermin comparatively studied Ekert’s entanglement-based and Bennett-Brassard’s disentanglement-based quantum key distribution (QKD) protocols. They concluded that so far security is concerned, the operational advantage of Ekert’s protocol is apparent. Since then, many quantum cryptographic protocols have been proposed and both type of cryptosystems have been extensively studied. But neither of the two type of systems can stake claim of its superiority.
In their comparative study , Bennett et al made an another important observation. They found that entanglement and disentanglement based cryptosystem are indistinguishable. That is, which type is being used cannot be distinguished by others. If sender uses entangled state but tells dishonestly to the receiver that he used disentangled state for the encryption, then receiver could not also verify the veracity of sender’s statement. In that sense, two cryptosystems are indistinguishable. It is recently understood conventional quantum bit commitment protocol ( a cryptographic application) completely fails because of this indistinguishability of two systems. Therefore, Bennett et al’s work has become helpful to examine other cryptographic tasks. Their work was based on conventional cryptography. Recently alternative disentanglement and entanglement based cryptographic protocols have been proposed . Many conclusions drawn from conventional quantum cryptography do not hold good in alternative quantum cryptography. So a fresh comparative study is necessary.
Alternative disentanglement-based cryptosystem uses mixed quantum state to encode a bit value but alternative entanglement based system uses many pure entangled states for the same purpose. Despite this dissimilarity, they have many similarities. Both can operate entirely on quantum channel and can provide quantum authentication. In both the systems, key can carry meaningful information. Secure bit commitment encoding and secure quantum coin tossing are possible for both the systems.
Yet the two systems are not well understood. We have seen that classical channel cannot be used in disentanglement-based system when each individual bit is separately made secure, but we do not know whether same is true for our entanglement-based system. We also do not know whether conventional cryptography or its prototype can be recovered from these alternative systems or not. Here we shall see that on these two questions two cryptosystems differ.
First, we shall present a modified (alternative) entanglement-based QKD protocol in which classical channel can be used when each bit is separately made secure and a prototype of conventional QKD protocol can be recovered from this protocol. This kind of modification is not possible for our disentanglement based system. This will imply that our entanglement-based system can be made much much faster than our disentanglement-based system.
Suppose a source emits pairs of spin $`1/2`$ particle in their singlet state. Two users, Alice and Bob, get one particle from each pair. Alice and Bob secretly share the information of two sequences of measurements. Suppose two sequences of direction of spin-measurements are:
$`S_0^n=\{x,x,y,y,x,y,y,y,x,x,y,x,\mathrm{}\mathrm{}\mathrm{}\mathrm{}..\}`$,
$`S_1^n=\{y,x,x,x,y,y,x,y,y,y,x,x,\mathrm{}\mathrm{}\mathrm{}\mathrm{}..\}`$,
where 0 and 1 in the subscripts stand for bit values and ”x” and ”y” are two orthogonal directions of measurements.
Let us assume they jointly decide the bit values. The bit values can be decided when both of them use the same sequence of measurements. To produce a key, both use $`S_0`$ and $`S_1`$ at random on their own sequences of EPR particles. When both use $`S_0`$ or $`S_1`$, the corresponding results will be perfectly correlated. But if one use $`S_0`$ and other $`S_1`$ or vice versa, the results will not be perfectly correlated. So $`50\%`$ bit value choices are discarded. The remaining $`50\%`$ bits form the key. We shall first assume they reveal results through classical public channel.
Their measurements yield the two sequences of data sets:
$`\{R_1^A,R_2^A,R_3^A,R_4^A,R_5^A,R_6^A,R_7^A,R_8^A,\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}.\}`$.
$`\{R_1^B,R_2^B,R_3^B,R_4^B,R_5^B,R_6^B,R_7^B,R_8^B,\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}.\}`$.
The first one is Alice’s sequence and second one is Bob’s. Half of their data sets contain perfectly correlated data.
Eavesdropper’s problem is to know the secret code of measurements. For simplicity, let us think they want to produce a single bit and only Bob’s particles are exposed to Eve. Eve can directly or indirectly measure using her own sequence of measurements. She gets a set of data from her measurements and taps Alice’s set of data when Alice reveals the results. But these two sets of data will neither reveal any bit information nor complete information of Alice’s choice of measurements. Now it is Eve’s turn to reveal the results. Alice’s and Eve’s results cannot be perfectly correlated. But this can be interpreted by Alice as a case of non-identical choice of bit values. Note that in $`S_0`$ and $`S_1`$ there is a common subsequence $`S_c`$. So if Alice does not get perfect correlation, she can check the data corresponding to $`S_c`$. Irrespective of choice of sequences of measurements, the data corresponding to $`S_c`$ will be always perfectly correlated. This second test will expose eavesdropping. Still it is not the last nail to eavesdropping.
The data are not secure because public channel is not authenticated channel. Eve can impersonate. After Alice’s disclosure of data, Eve, impersonating Bob, can reveal ”fake data” correlating with Alice’s data. Same thing she can do with Bob’s data impersonating Alice. Note that this attack works only for ”fake correlation”. That is, this attack will work when the users choose the same bit value. But they also choose different bit values in $`50\%`$ cases. In those cases, data are not perfectly correlated, only the data corresponding to $`S_c`$ are perfectly correlated. So initially if they do not get perfect correlation between their data sets, they will get perfect correlation in the subsets. As $`S_c`$ is hidden in $`S_0`$ and $`S_1`$, Eve could not generate ”fake correlation” in the data corresponding to the subset $`S_c`$. Therefore Eve can only impersonate to select the bits not to reject the bits. Eve can leave the task of rejecting the bits for the legitimate users. It seems that system fails.
There is a rescue. The ”fake correlation” attack works as both of them reveal all the data of the same events. The ”fake correlation” cannot be produced if they do not reveal any data. But the data has to be revealed if the system is to run. If they do not reveal all the results of the same events, yet the system can work but ”fake correlation” attack cannot. For clarity, suppose they divide the results of each set into two subsets. Alice’s subsets are $`r_1^A`$ and $`r_2^A`$ and Bob’s corresponding subsets are: $`r_1^B`$ and $`r_2^B`$. They reveal the data of non identical sets -that is, either $`r_1^A`$ and $`r_2^B`$ or $`r_2^A`$ and $`r_1^B`$. Because, the data of two correlated subsets are not revealed, ”fake correlation” attack will not work.
To create many bits, the strategy, discussed above, has to be repeatedly used to ensure bit by bit security. If any bit is found corrupted, the next bit will not be produced. If eavesdropping is detected they must reject $`S_0`$ and $`S_1`$ and may try with another two preselected sequences of measurements.
Is it possible to recover existing quantum cryptography from alternative quantum cryptography and vice versa ? Let us see the basic difference of the conventional and alternative systems. In conventional quantum cryptography, a pure state or a pure entangled state represent a classical bit/bits. On the other hand in alternative quantum cryptography, many states represent a classical bit. The bits of the conventional cryptosystem do not carry meaningful information but it carries meaningful information in alternative cryptosystem. Therefore recovery of alternative system from conventional system is not possible. But if we can produce pure state -bits (which does not carry any meaningful information) from alternative system then at least recovery of prototype of conventional system, if not the same system, will be possible. We have two options - recovery of conventional entanglement-based system from alternative entanglement based system and conventional disentanglement-based system from alternative disentanglement-based system. Next we shall see that the former can be easily realized.
We have seen that when both of them use $`S_0`$ or $`S_1`$, the data are perfectly correlated. These two sets of data can make a key provided they are not revealed. Suppose Alice divides the results into three subsets $`r_1^A`$, $`r_2^A`$ and $`r_3^A`$. Taking the results of same instances( which events will be chosen to construct the three subsets are not secret) Bob prepares his three subsets $`r_1^B`$, $`r_2^B`$ and $`r_3^B`$. Alice reveals $`r_1^A`$ and Bob $`r_2^B`$ or Alice reveals $`r_2^A`$ and Bob $`r_2^B`$. Both go through the correlation test. If correlation is found they know their chosen bit value and side by side they know the undisclosed subsets $`r_3^A`$ and $`r_3^B`$ contain perfectly correlated data. To construct $`r_3`$, it is better to use the data corresponding to $`S_c`$ so that they always get perfectly correlated data even when they use non-identical sequences of measurements. Continuing the process two different kind of keys (fast and slow keys) can be produced. The former does not exist even in the mind of the users and the later can exist in the mind of the users. If they want to produce only the former type they can share only a single sequence of measurement instead of the two.
The recovery of conventional entanglement-based system from alternative entanglement-based system is not possible. The reason is, a sequence of single photon polarized states produces sequence of results. These results can represent bit value. But if these results are revealed they cannot be secure.
Why do these systems differ on the above two issues ? From a sequence of two-particle entangled state Alice and Bob can get correlated random bits. And some of the correlated data are used for authentication via public channel and rest of the correlated data itself can make a fast key (which is a recovery of a prototype of conventional QKD protocol). Measurements on a sequence of disentangled states never produce correlated data. Therefore, entanglement is a necessary condition to have the above mentioned two utilities. As we get fast key from entanglement based system, entanglement is better secure number generator than disentanglement in our case.
Throughout our discussion we relied on bit by bit security. If we do not want bit by bit security but want security of many bits (meaningful) at a time, then such security needs to be proved. In that scenario there may have a possibility of using authenticated classical channel (authentication by some additional shared classical bits) in our disentanglement-based system. In that eventuality, we can say the use of classical channel and bit by bit security are mutually exclusive in our disentanglement-based system.
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# Controlling of exciton condensate by external fields and phonon laser
## Abstract
The novel method of observation and controlling of Bose-Einstein condensation in the system of spatially and momentum-space indirect excitons in coupled quantum wells using in-plane magnetic and normal electric fields is proposed. Fields are used for exciton dispersion engeneering. In the presence of in-plane magnetic field ground state of spatially indirect exciton becomes also indirect in the momentum space. Manipulation of electric field magnitude is used for tuning to resonance of transition of ”dark” long-life spatially and momentum-space indirect exciton in condensate into spatially and momentum-space direct exciton in radiative zone, with phonon creation, and with subsequent recombination of the direct exciton in radiative zone. Recombination rate of indirect excitons in condensate according to proposed scheme sharply rises at resonance. As a result besides enhanced spectral narrow photoluminescence connected with recombination of direct exciton in radiative zone with specific angular dependence, almost monochromatic and unidirectional beam of phonons appears. An opportunity of obtaining of ”phonon laser” radiation on the basis of this effect in considered.
PACS: 71.35Ji,71.35Lk,63.20Ls
The system perspective for observation of different phases of excitons is the system of spatially indirect excitons (SIE) in coupled quantum wells (CQW) (see Ref. and references therein). The point is that for the system of SIE in CQW - with spatially separated electron and hole (e and h) - recombination is suppressed by exponentialy small overlaping of e and h wave functions. Thus, the case can be easily achieved when the relaxation times of electrons, holes or excitons is sufficiently smaller than their lifetime. Therefore, there is an opportunity for reliable observation of different equilibrium exciton phases. In superfluid phase of the system of spatially separated e and h interesting phenomena can be observed, such as nondissipative electric currents in each quantum well, Josephson effects etc. . Very interesting recent experiments (see also essential previous works ) bear witness to an existence of coherent effects in SIE system at low temteratures.
In this paper we propose new method of the observation and controlling of coherent phase of SIE by external fields. We consider direct-gap $`{}_{}{}^{1}S`$-excitons in $`GaAs/AlGaAs`$ CQW. In-plane magnetic field shifts SIE dispersion minimum, where Bose-Einstein condensate (BEC) of SIE forms, out of the radiative zone. Energy of SIE in BEC is linear dependent on magnitude of normal electric field. Thus changing the electric field value enables one to tune the system of indirect excitons in condensate to resonance of process of transformation into spatially and momentum space direct excitons (SDE) in radiative zone, by emitting phonons, and subsequent recombination of the SDE. It occurs that photoluminescence (PL) of indirect excitons in condensate according to proposed scheme sharply rises and has specific angular dependence at resonance. Besides, almost monochromatic and unidirectional beam of phonons appears. Possibility of phonon laser creation using this effect is discussed.
There are four different bound states for a pair of e and h in CQW in the absence of external fields (Fig.1 a). Two of them, with e and h in the same quantum well, correspond to SDE placed in one of two quantum wells. The other two bound states correspond to SIE with e and h in different quantum wells. As a result of spatial separation of e and h SIE has normal to CQW electric dipole momentum $`\pm eD`$, where $`D`$ is interwell distance (signs correspond to two possible locations of e and h in two wells). In the case of identical quantum wells SDE and SIE have two-fold degeneracy. We suppose also that the uncertainty of exciton wave vector due to CQW roughness etc. is much smaller than all momenta involved in the problem. In the absense of external fields SDE is the lowest excitonic state.
Electric and magnetic fields change the dispersion of SIE (Fig.1 b). Normal electric field splits SIE sublevels as $`\mathrm{\Delta }\omega =eDE/\mathrm{}`$ due to electric dipole - electric field interaction. It can be shown that moving SIE has in-plane magnetic dipole momentum $`\pm \frac{eDp}{cM}`$ normal to its velocity, where $`p`$ is in-plane momentum and $`M`$ is a mass of the SIE. Thus in parallel magnetic field SIE sublevel dispersion curves move oppositely apart from the center of Brillouin zone by the quantity of $`K_m=\frac{eDH}{c\mathrm{}}`$ due to velocity-selective magnetic dipole - magnetic field interaction . This was experimentally observed in the work . SDE level, as in electric field, remains unchanged and degenerated.
Thus, by in-plane magnetic and normal electric fields one can change SIE sublevel position in dispersion space. In particular, applying normal electric field makes it possible for one of the SIE sublevels to become the ground state of the excitonic system (instead of SDE at E=0).
The process of recombination of excitons with emitting a photon is possible only inside the radiative zone - a small area in the center of Brillouin zone with $`k<k_rE_g/(c\mathrm{})`$, where $`E_g`$ is the energy gap in semiconductor, $`c`$ is the speed of light in the media and $`k_r`$ is the boundary wave vector of the radiative zone (in $`GaAs`$ $`k_r310^5cm^1`$). This is due to photon inability to carry away an in-plane momentum greater than $`k_r`$. Out of the radiative zone the process of photon recombination of the exciton is forbidden, and the process of exciton recombination with production of an acoustic phonon and a photon becomes the most probable exciton recombination process . This process is of the second order of magnitude of the exciton interaction with photon and acoustic phonon fields. Greatest parts of energy and 2D momentum of exciton are carried away by photon and phonon, correspondingly.
From preceding discussion, the following scheme of the observation of BEC of SIE can be developed. In the presence of such normal electric and in-plane magnetic fields that SIE is the ground excitonic state with minimum of its dispersion curve situated out of the radiative zone, SIE ground level is being pumped by laser radiation into the radiative zone (Fig. 2a). After laser switch off the system is evolving using two channels - one part of pumped excitons in radiative zone recombines with production of photons, and the other part of excitons relaxate into the minimum of dispersion of SIE level out of the radiative zone, tranferring their extra energy and momentum to acoustic phonons (Fig. 2b). Thereafter, the system of SIE reaches a quasi-equilibrium state with the temperature of reservoir (lattice), and in case of sufficiently low temperature of lattice the system of SIE forms BEC. We will consider $`T=0`$ for simplicity; at temperatures much less than Kosterlitz-Thouless temperature optical properties are changed only slightly with respect to those at $`T=0`$. The state of the system at this moment of evolution is the state where almost all excitons are at SIE level in its dispersion minimum with wave vector $`\underset{m}{\overset{}{\mathrm{K}}}`$ (for low density of excitons) and all the other levels are not populated. In high magnetic field we have $`K_mk_r`$ (e.g. $`H=10T`$ and $`D=10nm`$ corresponds to $`K_m=210^6cm^1`$).
Bose-condensed SIE situated in momentum space at $`K=K_m`$ position are ”dark”, that is the process of emitting of photons is weak. The rate of recombination SIE in BEC through intermediate production of virtual exciton in radiative zone (with production of photon and phonon with wave vectors in the region of momentum space $`d^3q,d^3k`$) is
$`dW_{phn,pht}=({\displaystyle \underset{s}{}}{\displaystyle \frac{(M_{s,pht}M_{s,phn})^2}{(\omega _{pht}(\stackrel{}{𝐤})\omega _s(\stackrel{}{\mathrm{k}}))^2+\eta _s^2}}+{\displaystyle \underset{s_1,s_2,s_1s_2}{}}({\displaystyle \frac{M_{s_1,pht}M_{s_1,phn}M_{s_2,pht}M_{s_2,phn}}{2(\omega _{pht}(\stackrel{}{𝐤})\omega _{s_1}(\stackrel{}{\mathrm{k}})+i\eta _{s_1})(\omega _{pht}(\stackrel{}{𝐤})\omega _{s_2}(\stackrel{}{\mathrm{k}})i\eta _{s_2})}}+c.c.))`$ (1)
$`\delta (\stackrel{}{\mathrm{k}}+\stackrel{}{\mathrm{q}}\underset{m}{\overset{}{\mathrm{K}}})\delta (\omega _{pht}(\stackrel{}{𝐤})+\omega _{phn}(\stackrel{}{𝐪})\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}}))(2\pi )^3\rho N_\stackrel{}{𝐤},N_\stackrel{}{𝐪}d^3qd^3k`$ (2)
In Eq.(2) $`\rho `$ is 2D density of SIE in CQW; $`N_\stackrel{}{𝐤},N_\stackrel{}{𝐪}`$ are numbers of photons and phonons in the states with wave vectors $`\stackrel{}{𝐤}`$ and $`\stackrel{}{𝐪}`$, correspondingly; notation $`\stackrel{}{\mathrm{x}}`$ implies 2D in-plane vector $`x`$ or in-plane component of 3D vector $`\stackrel{}{𝐱}`$; $`M_{s,phn}`$ is the matrix element of transformation of SIE into $`s`$-state exciton in radiative zone with acoustic phonon creation, and $`M_{s,pht}`$ is the matrix element of photon recombination of $`s`$-exciton in radiative zone; $`\eta _s`$ is the width of $`s`$th exciton level. Using $`GaAs`$ parameters, we estimate for SDE $`\eta _{SDE}10^{10}sec^1`$. Recombination rate consists of resonant terms, corresponding to the processes with intermediate creation of virtual s-exciton in radiative zone, and nonresonant terms.
Main contribution to the rate of SIE recombination process is originated from recombination transitions, in which virtual s-excitons are mostly close to their mass shell. SIE and SDE are such s-excitonic states in our case. At this moment of evolution (Fig.2b) the process goes mainly through virtual SIE in radiative zone. For the estimation we take phonon energy to be equal to $`\omega _{phn}(\stackrel{}{𝐪})c_vK_m`$, where $`c_v`$ is a speed of sound. In this case $`\omega _{pht}(\stackrel{}{𝐤})\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})c_vK_m`$ and we get:
$`W_{phn,pht}{\displaystyle \frac{(M_{SDE,pht}M_{SDE,phn})^2}{(c_vK_m(\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})\omega _{SDE}(0)))^2+\eta _{SDE}^2}}+{\displaystyle \frac{(M_{SIE,pht}M_{SIE,phn})^2}{(c_vK_m(\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})\omega _{SIE}(0)))^2+\eta _{SIE}^2}}`$ (3)
In this formula we also neglected non-resonant terms. Parameter $`(\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})\omega _{SIE}(0))`$ is not dependent on electric field magnitude, while $`(\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})\omega _{SDE}(0))`$ is a linear function of electric field magnitude, and equals zero when the electric field has a particular (resonant) value. For CQW system used in Ref. at magnetic field $`10T`$ these parameters are correspondingly equal to $`0.810^{13}sec^1`$ and $`3.510^{13}sec^1`$.
At this stage of the system evolution one can reduce the magnitude of electric field down to the moment, when the conservation laws for energy and 2D momentum are satisfied in the processes of SIE in BEC transformation into real SDE in radiative zone with production of acoustic phonon and SDE in radiative zone photon recombination (Fig.2c) (in other words, when resonance condition is satisfied). Using Eq.(3), one can get the estimation for the ratio of PL rates at initial condition and after tuning to the resonance in the system used in Ref.:
$`{\displaystyle \frac{W_{resonance}}{W_{initial}}}\alpha ^2{\displaystyle \frac{(\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})\omega _{SIE}(0))^2}{\eta _{SDE}^2}}10^{5÷6}\alpha ^2,\alpha ={\displaystyle \frac{M_{SDE,pht}}{M_{SIE,pht}}}{\displaystyle \frac{M_{SDE,phn}}{M_{SIE,phn}}}`$ (4)
Martix elements of phonon-exciton interaction vertex are determined by Bardeen-Shockley deformational potential Ref.. Matrix element of photon recombination of an exciton in radiative zone is proportional to the overlaping integral of e and h in the exciton . Using Eq.(4) one can show that PL intensity can be increased, by changing the electric field, at least by two orders. Therefore, after tunning the system to the resonance the recombination rate will be greatly increased, and this will result in sharp PL intensity growth.
Now we will analize angular dependence of resonant PL. Let magnetic field be directed along x-axis. In this case $`\underset{m}{\overset{}{\mathrm{K}}}`$ vector is parallel to y-axis. Studying resonant radiation, in Eq.(2) we take only resonant term responsible for SIE recombination by intermediate SDE creation. Integrating (2) over $`d^3q`$, we obtain the rate of photon emission in k-space area $`d^3k`$:
$`dW_{pht}^{res}`$ $``$ $`{\displaystyle \frac{1}{\left[(\omega _{pht}(\stackrel{}{𝐤})\omega _{SDE}(\stackrel{}{\mathrm{k}}))^2+\eta _{SDE}^2\right]}}{\displaystyle \frac{(\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})\omega _{pht}(\stackrel{}{𝐤}))d^3k}{\sqrt{(\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})\omega _{pht}(\stackrel{}{𝐤}))^2c_v^2(\underset{m}{\overset{}{\mathrm{K}}}\stackrel{}{\mathrm{k}})^2}}}`$ (5)
The maximum energy of the photon is $`\mathrm{}c\frac{\omega _{SIE}(\underset{m}{\overset{}{\mathrm{K}}})c_vK_m}{c+c_v}`$. Photons with any energy less than its maximum value can be created in the process, but the rate of production of photons with small energy is suppressed by Lorentz factor.
The process becomes resonant in the case when the maximum of the Lorentz factor and the singularity of the second factor coincide for at least some values of $`\stackrel{}{\mathrm{k}}`$. This condition can be made clear graphically (Fig.3a). Resonant condition is satisfied if in the space ($`\stackrel{}{\mathrm{p}},\omega `$) a dispersion cone of 2D phonon, drawn down from BEC position, intersects SDE dispersion surface in the radiative zone. We stress, that 2D dispersion of 3D phonons gives only the minimum energy of a phonon with a given component of 3D wave vector on CQW plane (the same is true for photons). The minimum energy corresponds to zero normal to CQW component of phonon wave vector. For this reason, resonant process with SIE recombination with intermediate creation of real SDE takes place even in case of greater energy of exciton in BEC than resonance one (i.e. in smaller elecric fields), but with less intensity.
The resonance condition corresponds to a distinct range of electric field magnitude. We represent exciton in BEC energy dependence on electric field magnitude as $`\omega _{SIE}(K_m)=c_v(K_m+1/2k_rT(E))`$, where $`T(E)`$ is a linear function of electric field. Since $`K_mk_r`$, we can admit that in the case of resonance, 2D phonon dispersion (cone) intersects the SDE dispersion surface in a circle (Fig.2 b):
$`(k_xk_0)^2+k_y^2=k_rk_0T(E)+k_0^2,k_0={\displaystyle \frac{Mc_v}{\mathrm{}}}`$ (6)
For $`GaAs`$ we have $`k_0=10^5cm^1`$. The condition of the resonance can be represented as $`2+\frac{k_r}{k_0}T(E)\frac{k_0}{k_r}`$.
We will below admit resonance approximation, i.e. substitution for the Lorentz factor by delta-function $`\delta (\omega _{pht}(\stackrel{}{𝐤})\omega _{SDE}(\stackrel{}{\mathrm{k}}))`$. This approximation is valid when $`\eta _{SDE}c_vk_r`$, so that it is applicable in our case since $`c_vk_r=1.510^{11}sec^1`$. In result one can get:
$`dW_{pht}^{res}(\theta ,\varphi ){\displaystyle \frac{d\mathrm{\Omega }}{\sqrt{T(E)\beta sin^2(\varphi )sin^2(\theta )cos(\theta )}}},\beta =k_r/k_03`$ (7)
Here $`d\mathrm{\Omega }`$ is the spatial angle, $`\theta `$ and $`\varphi `$ are azimuth and polar anlges of spherical coordinates in photon wave vector $`\stackrel{}{𝐤}`$ space. Azimuth axis ($`\theta =0`$) is directed parallel to $`\underset{m}{\overset{}{\mathrm{K}}}`$. Eq. (7) gives a resonant radiation angular dependence. Resonant PL is absent when $`T(E)<k_0/k_r`$. At resonance ($`2+k_r/k_0T(E)k_0/k_r`$) resonant photons are emitted into spatial angle, which is formed by intersection of sphere with radius $`k_r`$ and a cylinder with base (6) and generatrix parallel to z-axis (Fig.3 c). In the resonant approximation on boundary of this spatial angle the rate of resonant PL has an integrating singularity. It can be shown that real rate of PL right on this boudary is approximately as much as $`\frac{c_vk_r}{\eta _{SDE}}15`$ times greater than in the other area of spatial angle of the resonant radiation. If $`T(E)>2+k_r/k_0`$ resonant PL intensity has no singularity on the boundary of spatial angle of resonant PL, and with increasing $`T(E)`$ (reducing electric field) tends to become more and more isotropic.
When resonance condition is satisfied, energy level scheme of our system is similar to that of three-level impulse one-pass laser with inverse population of upper level and with rapid lower transition (Fig.4 a). Rapidity of the lower transition with respect to the the rate of upper transition is supplied by tunneling character of transformation of SIE into SDE with creation of acoustic phonon. The difference is that in our case phonons are emitted instead of photons. Wave vectors of resonant phonons form a prolate ellipsoid of revolution with base (6) and the ratio of radii $`\sqrt{K_m/k_0}`$ (Fig.4 b). The best quality of unidirectivity and monochromatism of the resonant phonons is evidently achieved in electric field that corresponds to the relation $`T(E)=k_0/k_r`$, when the ellipsiod reduces to a point, so that for discussion of phonon ”laser” radiation we will consider this case.
At this resonance condition the coherence of resonant phonons at early stages of the process of phonon emission is determined by total (in a sense ”homogeneous” and ”inhomogeneous”) width of phonons. Width of wave vector distribution of phonons can be represented as $`\mathrm{\Delta }q=\sqrt{k_0(\mathrm{\Delta }q_{hom}+\mathrm{\Delta }q_{inh})}`$, where $`\mathrm{\Delta }q_{hom}`$ is due to widths of SDE level and phonon state, and $`\mathrm{\Delta }q_{inh}`$ is connected with SIE momenta spread that corresponds to depleting of BEC due to interactions of excitons; $`\mathrm{\Delta }q_{inh}=\sqrt{U(0)\rho M}/\mathrm{}`$, where $`U(0)`$ is zero Fourier component of interaction between SIE; $`\mathrm{\Delta }q_{hom}=(\eta _{SDE}+\eta _{phn})c_v^1`$, where $`\eta _{phn}`$ is the phonon attenuation rate.
There are two sufficient conditions for phonon radiation to be the ”laser” one: generation condition and condition of macroscopic population of each quantum state. In our case the first condition is the prevailing rate $`\eta _{spon}`$ of spontaneous phonon emission in the process over the rate of phonon attenuation in the sample $`\eta _{phn}`$. Attenuation of the phonon is determined by the media properties, temperature, the size of the sample and scattering of phonon on the planes of CQW. Using data from Ref. and Eq.(4) one can get, that $`\eta _{spon}`$ can be not simply greater than the phonon attenuation rate, but even can reach the phonon energy $`K_mc_v10^{12}sec^1`$.
The second condition is macroscopic population of each phonon quantum state (when all SIE in BEC are recombined and thus have contributed to phonon radiation). Phonon radiation occupies $`V(\mathrm{\Delta }q)^3/(2\pi )^3`$ quantum states, where $`V`$ is a volume of the sample, while the number of SIE is $`S\rho `$, $`S`$ is the area of CQW. The condition is $`L_z(2\pi )^3\rho /(k_0(\sqrt{U(0)\rho M}/\mathrm{}+\mathrm{\Delta }q_{hom}))^{3/2}`$, where $`L_z`$ is the width of the sample in $`z`$-direction. The quantity $`L_z`$ appeared due to that phonons and SIE differ in their dimensions (3D phonons vs. 2D excitons).
In conclusion we emphasize that discussions above point out possible existence of interesting effect - ”phonon laser” (details will be published elsewhere). Using all advantages of the analogy with photon laser, one can propose to adjust a phonon resonator for the purpose of increasing phonon coherence. The resonator can experimentally be realized as phonon mirrors (the media with greater speed of sound) on (y-direction) edges of the sample. In this case one-pass phonon ”laser” becomes multi-pass one. The restriction for the width of the sample (second condition) can be made weaker by considering heterostructure consisting of many CQW (the distance between CQW must satisfy the above condition for $`L_z`$).
Note that, coherent phonons modulate the dielectric function of the media. This gives an opportunity of detection of the effect of coherent phonon generation by study the modulation of optical properties of the media, which can be observed by time-resolved femtosecond spectroscopy (e.g. by pump-probe method). Another possibility to observe coherent properties of phonon radiation is analysis of its statistics, e.g. by Hunbury-Brown and Twiss method.
The work has been supported by Russian Foundation of Basic Research, INTAS and Programm ”Solid State Nanostructures”.
Captures to Figures.
Fig.1 Dispersion laws of direct and indirect excitons in coupled quantum wells a) Without extrenal electromagnetic fields SDE level is lower than SIE level by the difference of Coloumb interactions in SDE and SIE b) In parallel magnetic field SIE sublevel dispersions move oppositely apart by wave vector $`K_m=\frac{eHD}{c\mathrm{}}`$; in normal electric field one of SIE sublevels lowers as $`\mathrm{\Delta }\omega =eED/\mathrm{}`$ and the other rises as $`\mathrm{\Delta }\omega `$
Fig.2 a) Laser pumping of SIE level in the presence of parallel magnetic and normal electric fields. $`K_m=\frac{eHD}{c\mathrm{}}`$ is the displacement of dispersion minimum in parallel magnetic field, $`k_r`$ is boundary wave vector of radiative zone. b) Relaxation of indirect excitons. One part of excitons recombines with production of photons, and the other part relaxates to the minimum of SIE dispersion $`K_m`$. c) Changing of BEC position when electric field magnitude is being reduced.
Fig.3 a) Position of direct and indirect exciton dispersions when condition of the resonance is satisfied. b) intersection of dispersion surface of indirect excitons with 2D phonon dispersion. c) segments of spatial angle correspond to the boundary of resonant photoluminescence angular dependence.
Fig.4 a) Three-level laser analog of phonon laser. b) Phonon wave vector distribution. Wave vectors of phonons form prolate ellipsond of revolution with ratio of its radii $`\sqrt{\frac{K_m}{k_0}}`$.
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# References
Polygonization and plastic bending of whiskers induced by pulsed power electron beam
A.A. Petrova
Tomsk Polytech University, Tomsk, Russia
## Abstract
Dynamics of bending of NaCl whisker is studied with use of high speed camera VFU-1 synchronized with high power electron accelerator GIN-600. The dynamics of bending turns to be non-monotonous and quite complicated. The figures of etching of whisker after bombardment by electron beam are analysed, thin polygonization in samples is found. The time of bending in experiment and time of forming of dislocation walls computed on the base of theory of plastic bending are compared.
It was shown in the paper that intense plastic bending of whiskers induced by nanosecond pulse of bombardment by high density electron beam takes place.
Once more process of relaxation of mechanical tensions generated by bombardment is found. This process represents itself as an intense plastic deformation connected with arising and sliding of dislocations. This process quite quickly compensates thermoelastic tensions arising in a sample during bombardment and decreases resulting tension to be lower than threshold of movement for cracks thus preventing fragile destruction. However, dynamics and mechanism of bending were not studied in details.
In this paper estimation of time of plastic bending and its medium speed obtained from experiment are given and analysis of these data is carried out.
The time of plastic bending was measured in experiment using methods described in details in .
Parameters of experimental system
| Maximal energy of accelerated electrons | 0.35 $`MeV`$ |
| --- | --- |
| Current of beam of electrons | 5 $`kA`$ |
| Density of current | 4 $`kA/cm^2`$ |
| Diameter of homogeneous part of beam | 8 $`mm`$ |
| Time of pulse | 2– 30 $`ns`$ |
| Speed of filming | 25000– 25000000 stills per second |
| Maximal time distinguishing | 2 $`\mu s`$ |
| Total number of stills | 49 |
| Diameter of one still | 10 $`mm`$ |
| Vacuum in experimental camera | 1 $`Pa`$ |
The scheme of experiment is given at Fig. 1. It was carried out in the following consequence.
A sample of NaCl whisker (6) fixed on special needle holder was posed in vacuum experimental box (2). The flash lamp (7) for lighting the sample during filming is established inside of experimental box. The electron accelerator (1) is connected via vacuum tract with experimental camera (2). The optical tract connects experimental box with high speed photographic system VFU-1 which is able to make still-by-still filming of an object. Synchronization scheme (4) orders moments of switching of accelerator, flash lamp and high speed photographic camera. The sample brighten by flash lamp is pierced by high density electron beam during some nanoseconds after which it begins to bend slowly, and high speed photographic camera carried out still-by-still filming of the process.
The results of filming are the primary source of information about plastic bending of whiskers under nanosecond pulse of bombardment by high density electron beam. The results of experiment are given by Fig. 2, in which one typical time dependence of deviation h of free end of whisker from initial position during bending process. One can see that full time of bending is $`1.5ms`$. There is an induction period 0.1–0.7 ms during which one bending is small and is not observed practically. The calculation of thermoelastic tensions arising in thin plates and bars as a result of inhomogeneous bombardment by high density electron beam is carried out in .
The spectrum of electron beam contains essential number of electrons with energy no more than 100 keV. Distribution of dose of bombardment per thickness of whisker is inhomogeneous: as the front part of whisker is heated more, as rear one is heated less. The thermoelastic tensions arisen stretch rear part of whisker and compress its front part. Under these tensions whisker should elastically bend up, i.e. with prominent part turned to beam, and make elastic bending oscillations with respect to curved quasistatic profile. Both these effects are of course present but they cannot be seen for photographic camera. They can be fixed by laser interferometer . Absence of visible deformation during first 0.1 – 0.7 $`ms`$ shows also that mechanical forces acting in crystal are not so large, i.e. thermoelastic tensions induced by inhomogeneous bombardment are relatively quickly compensated by tensions of other origin. Those ones must manifest as thermoelastic tensions are relaxing because of smoothing of temperatures during heat conduction process.
The characteristic time of smoothing temperature due to heat conductivity is
$$\tau =a^2/\chi (1)$$
(1)
where $`a`$ is a characteristic linear scale of temperature inhomogeneous domain, $`\chi `$ is a temperature conductivity coefficient. In our experiment $`a`$=(0.5 - 1) $`b`$, where $`b`$ is width of whisker, $`\chi =\lambda /(c\rho )`$, $`\lambda `$ is a heat conductivity coefficient, $`c`$ is a capacity, $`\rho `$ is a density. For NaCl we have $`\lambda =7.4W/(m\times K)`$ at 300 K; $`\rho =2.16\times 10^3kg/m^3`$, $`c=0.86\times 10^3J/(kg\times K)`$ hence $`\chi =4\times 10^6m^2/s`$. Substituting $`a=(0.51.0)b=(48.897.6)10^6m`$ and the value of $`\chi `$ to expression (1) we obtain $`\tau =(0.62.4)ms`$.
Hence we conclude that the time of relaxation of quasistatic thermoelastic tensions generated by inhomogeneous bombardment (and consequently by inhomogeneous heating) of whisker is about 1 ms. It is easy to find nature of forces compensating thermoelastic ones. When the value of thermoelastic tensions stretching upper part of whisker becomes more than value of ”flowing tooth” ($`0.7\times 10^8Pa`$ for NaCl ), the intense sliding and generation of dislocations of the same sign begins. These dislocations slide from upper stretched surface lower to medium unstrained part of a sample and stay on the threshold of lower compressed part. This is the dislocation mechanism of bending for plates and bars .
As a result of this process upper part of whisker stretched by thermoelastic tensions is filled by big number of extraplanes limited below by edge dislocations of one sign. These extraplanes generate tensions preventing from bending of crystal up, with prominent part turned to beam. As temperature is smoothed and thermoelastic tensions are relaxed the situation is changed to opposite one. The upper part of the sample filled with superfluous extraplanes becomes compressed, and lower one – stretched. The samples is bent to opposite side, with concave side turned to beam.
The experiment shows that bending takes place in process of elastic-plastic non-harmonic decaying oscillations of the sample. Firstly h increases, reaches maximum, then decreases and again grows up to some stationary value. This is the typical dynamics of bending (Fig.2).
Two types of pictures are observed. On first ones after bombardment and etching of the sample one can see branching on lateral sides of whisker turned by its angle to electron beam (Fig.3). On second ones after bombardment and etching one can see branching turned by its angle to opposite side (Fig.4). First type corresponds to whiskers containing occasional defects, cracks or bubbles, and second one – to whiskers in those edge dislocations in main planes of sliding were introduced via stretching on microdeformation machine. Quantitative pictures (Figs. 3 and 4) are not equivalent but qualitatively lead to the same result. In both cases the sample bent to the beam.
To understand why figures of etching are different let us turn to mechanisms of plastic bending and describe behaviour of dislocations during bending. It is known that during plastic bending dislocations are stayed to be grouped in sliding lines (see Fig. 5) laying in near sliding plates with concentrating in most energetically advantageous directions for corresponding material. Such a direction for NaCl is $`<110>`$ (Fig. 6). At Fig. 6 movement of dislocation loop in plane from bombarded surface to opposite side of crystal in a sample is given. The lateral (parallel to the vector $`b`$) sides of loop are spiral dislocations, and the side of loop which is perpendicular to vector $`b`$ is an edge dislocation which reaches surface of crystal faster. If edge dislocation as a component of dislocation loop does not reach surface of sample one can observe a pyramid- like pit on lateral surface (100) under etching in point of its crossing with lateral surface. These pits form strips which one can see at Figs. 3 and 4. Only after additional heat acting process of decaying of rows of dislocations in sliding planes and forming of vertical rows of dislocations (polygonization) can take place (Fig. 7). Let us imagine that these edge dislocations slide down sliding plates and stand approximately one under another. The picture of etching is changed but structure of bending of the crystal is not changed since additional planes which edges are edge dislocations are introduced from previous direction, e.g. from side opposite to bombardment.
Interaction of edge dislocations of the same sign in parallel sliding planes in favorable conditions leads to forming of vertical rows of edge dislocations. The vertical row of dislocations situated on distance $`l`$ from each other forms a border of blocks symmetrically turned with respect to each other around axis parallel to dislocation for an angle $`\theta =b/l`$.
These processes are most visible during plastic bending of crystal. Arising of rows of dislocations takes place successively in plastic deformation process. In our case in NaCl crystals thin polygonization is observed. It is possible only under essential plastic deformation when second sliding system begins to act. This sliding plane is perpendicular to plane (011) on Fig. 6. Not only pairs of dislocations but also groups of them in orthogonal sliding planes can be stable. Such a groups in the case of dislocations of the same sign are so called dislocation walls, i.e. rows of dislocations perpendicular to sliding planes. When deformation is inhomogeneous (as in our case) and essential surplus of dislocations of one sign presents favorable conditions for forming walls (i.e. rows of dislocations) are established . Reconstruction of border-like dislocation structures from nets of chaotically situated dislocations, especially after heating, and in our case electron beam heats the sample creating temperature gradient about 100 degrees between surface on which electron beam falls and that one opposite, is connected with decreasing of elastic energy stored in crystal when dislocations form walls. Forming of subborders in process of plastic deformation itself takes place as a result of local reconstruction of dislocations in sliding strips which also leads to decreasing of elastic energy . If one turns the curved sample by prominent side to beam and continues pulse bombardment the sample will become straight after several pulse, then it is curved to opposite side. This phenomenon is agreed with polygonization mechanism and represents itself as an usual process in ion crystals . In curved samples polygonization etching pits lay along straight lines corresponding to borders of blocks.
The value of mutual desorientation of blocks is well agreed with number of dislocations in these borders. The distance between edge dislocations of the same sign in NaCl wall $`h=100b`$ which is found from pictures of etched samples. Then angle of desorientation is $`\theta =b/100b`$ where $`b`$ is a Burgers vector $`b=4\times 10^8cm`$ since vertical row of dislocations situated on distance l from each other forms border of blocks symmetrically turned with respect to each other around axis parallel to dislocation line, and the desorientation angle $`\theta =b/h`$. In our case $`\theta =0.6^0`$.
Number of dislocation walls calculated on pictures of etched samples is about 20, so total angle of bending of whisker coincides with theoretical angle $`12^0`$. It is necessary to find speeds of sliding of edge dislocations in NaCl crystal under action of tensions and to estimate time necessary for forming dislocation walls. Real tensions acting in NaCl whisker under bombardment by fast electrons are $`\sigma =0.7\times 10^8Pa`$. In the dependence $`v(\sigma )`$ was found from which one concludes that speed of edge dislocations at such tensions is about $`10^2m/s`$. If medium thickness of whisker about 60 $`\mu m`$ dislocation needs $`6\times 10^7s`$ to pass distance from place of arising to opposite side of crystal. This time is enough for construction of dislocation walls since time of bending is about 1 ms.
Conclusions.
1. Dynamics of elastic-plastic bending of NaCl whisker induced by nanosecond pulse of bombardment by intense electron beam is studied with use of high-speed still-by-still filming. The measured time of irreversible plastic bending of whisker lies in an interval 1.5 - 2.5 ms and practically coincides with characteristic time of relaxation of quasistatic thermoelastic tensions arising due to inhomogeneous heating of whisker by beam.
2. The observed bending is not elastic in all characteristics: time, space and thermodynamic ones. The main contribution to bending is given by elastic-plastic deformation connected with movement and multiplication of dislocations.
3. Different pictures of etching of whiskers bent to beam are found. At detailed consideration it turned to be that dislocation walls are seen on the pictures. The estimations of sliding of dislocations are made. The angle of bending of samples is connected with the angle of desorientation of blocks forming during polygonization. These estimations are agreed with experimental data.
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# Untitled Document
METASTABILITY AND SMALL EIGENVALUES
IN MARKOV CHAINS
Anton Bovier <sup>1</sup><sup>1</sup>1 Weierstrass-Institut für Angewandte Analysis und Stochastik, Mohrenstrasse 39, D-10117 Berlin, Germany. e-mail: bovierwias-berlin.de, Michael Eckhoff<sup>2</sup><sup>2</sup>2Institut für Mathematik, Universität Potsdam, Am Neuen Palais 10, D-14469 Potsdam, Germany. e-mail: eckhoffrz.uni-potsdam.de,
Véronique Gayrard<sup>3</sup><sup>3</sup>3DMA, EPFL, CH-1021 Lausanne, Switzerland, and Centre de Physique Théorique, CNRS, Luminy, Case 907, F-13288 Marseille, Cedex 9, France. email: Veronique.Gayrardepfl.ch, Markus Klein<sup>4</sup><sup>4</sup>4Institut für Mathematik, Universität Potsdam, Am Neuen Palais 10, D-14469 Potsdam, Germany. e-mail: mkleinfelix.math.uni-potsdam.de
Abstract: In this letter we announce rigorous results that elucidate the relation between metastable states and low-lying eigenvalues in Markov chains in a much more general setting and with considerable greater precision as was so far available. This includes a sharp uncertainty principle relating all low-lying eigenvalues to mean times of metastable transitions, a relation between the support of eigenfunctions and the attractor of a metastable state, and sharp estimates on the convergence of probability distribution of the metastable transition times to the exponential distribution.
Keywords: Markov chains, metastability, eigenvalue problems, exponential distribution
I. Introduction.
The phenomenon of metastability has been a fascinating topic of non-equilibrium statistical mechanics for a long time. Currently, it has found renewed interest in the investigation of glassy systems and aging phenomena which appear to play a central role in many physical and non-physical systems. An approach to link metastability to spectral characteristics, in particular low-lying eigenvalues and the corresponding eigenfunctions has been proposed by Gaveau and Schulman \[GS\]. Such an approach is appealing not only because it allows to characterize metastability in terms that are intrinsically dynamic and make no a priori reference to geometric concepts such as “free energy landscapes”, but also since it allows numerical investigations of metastable states via numerical spectral analysis (see Schütte et al. \[S,SFHD\] for applications to conformational dynamics of biomolecules).
Relating metastability to spectral characteristics of the Markov generator or transition matrix is in fact a rather old topic. First mathematical results go back at least as far as Wentzell \[W\] (see also \[M\] for more recent results) and Freidlin and Wentzell (see \[FW\]). Freidlin and Wentzell relate the eigenvalues of the transition matrix for Markov processes with exponentially small transition probabilities to exit times from “cycles”; Wentzell has a similar result for the spectral gap in the case of certain diffusion processes. All these relations are on the level of logarithmic equivalence, i.e. of the form $`lim_{ϵ0}ϵ\mathrm{ln}(\lambda _i^ϵT_i^ϵ)=0`$ where $`ϵ`$ is the small parameter, and $`\lambda _i^ϵ,T_i^ϵ`$ are the eigenvalues, resp. exit times. This rather crude level of precision persists also in the more recent literature and prevents, in particular, applications to systems with unbounded numbers of metastable states which are particularly relevant for glassy systems.
In this letter we announce results that – for a large class of Markov chains – improve this situation considerably: in particular we allow for the number of metastable states to grow (with e.g. the ‘volume’), and we give precise control of error terms for ‘finite volume’ systems. Moreover, we provide representations for all quantities concerned that are computable in terms of certain ‘escape probabilities’ that are in turn well controllable through variational representations \[BEGK1\]. A more detailed exposition of our results, as well as the proofs, will be given in two forthcoming papers \[BEGK2,EK\]. 2. Stable set and metastable states. We will consider in the sequel Markov chains $`X_t`$ with state space $`\mathrm{\Gamma }_N`$, discrete time<sup>5</sup><sup>5</sup>5All results apply, however, also to continuous time. $`t`$, and transition matrices $`P_N`$. We will assume that for any fixed $`N`$ they are ergodic, and have a unique invariant distribution $`_N`$. We are interested in the situation when these chains exhibit “metastable” behaviour; loosely speaking, this means that the state space $`G_N`$ can be decomposed into subsets $`S_{N,i}`$ such that the typical times the process takes to go from one such set to another are much larger than the time it takes to “look like” being in equilibrium with respect to the conditional distribution $`_N(|S_{N,i})`$. Some reflection shows that this statement has considerable difficulties and cannot be interpreted literally, and that a precise definition of metastability is a rather tricky business (see e.g. the recent discussion in \[BK\]). We will give a precise definition that is, however, inspired by this vague consideration. The main point here is that one should make precise the two time scales we alluded to. We will take the following attitude: to look ergodic within $`S_{N,i}`$, the process should have at least enough time to reach the “most attractive” state within $`S_{N,i}`$, while at least the times to go from two such states in different metastable regions should be long compared to that time. Note that this concept is rather flexible and allows, in general, to define metastable states corresponding to different time scales.
The following definition of “stable sets” follows this ideology; however, we prefer to use certain probabilities rather than actual times as criteria, mainly because these are more readily computable. Linking them in a precise manner to times will be part of our results. We will write $`\tau _I^x`$, for $`x\mathrm{\Gamma }_N`$, $`I\mathrm{\Gamma }_N`$, for the first non-zero time the process started at $`x`$ arrives at $`I`$.
Definition 2.1: A set $`_N\mathrm{\Gamma }_N`$ will be called a set of stable points if it satisfies the following assumptions: There exist finite positive constants $`a_N`$, $`b_N`$, $`c_N`$, and $`r_N`$ such that they satisfy for some sequence $`\epsilon _N`$ s.t. $`|\mathrm{\Gamma }_N|\epsilon _N0`$, $`a_N^1\epsilon _Nb_N`$.
(i) For all $`z\mathrm{\Gamma }_N`$,
$$\left[\tau __N^z<\tau _z^z\right]b_N$$
$`(2.1)`$
(ii) For any $`xy_N`$,
$$\left[\tau _y^x<\tau _x^x\right]a_n^1$$
$`(2.2)`$
We associate with each $`x_N`$ its local valley
$$A(x)\{z\mathrm{\Gamma }_N:\left[\tau _x^z=\tau __N^z\right]=\underset{y_N}{sup}\left[\tau _y^z=\tau __N^z\right]\}$$
$`(2.3)`$
Then
$$r_N\frac{_N(x)}{_N(A(x))}R_xc_N^1$$
$`(2.4)`$
We will also write $`T_{x,I}[\tau _I^x\tau _x^x]^1`$. An important characteristic of the sets $`I_N`$ is $`T_Isup_{x_N}T_{x,I}`$. A simplifying assumption, that will be seen to ensure sufficient “spacing” of the low lying eigenvalues is that of “genericity”, defined as follows:
Definition 2.2: We say that our Markov chain is generic on the level of the set $`_N`$, if there exists a sequence $`ϵ_N0`$, s.t.
(i) For all pairs $`x,y_N`$, and any set $`I_N\backslash \{x,y\}`$ either $`T_{x,I}ϵ_NT_{y,I}`$ or $`T_{y,I}ϵ_NT_{x,I}`$.
Each of the elements of $`_N`$ in the generic case will then correspond indeed to a metastable state. Our first task is to identify precisely the notion of the exit time from a metastable state. To do so, we define for any $`x_N`$ the set $`_N(x)\{y_N:_N(y)>_N(x)\}`$; these are the points that are even more stable than $`x`$. The metastable exit time, $`t_x`$ from $`x`$ is then defined as the time of the first arrival from $`x`$ in $`_N(x)`$. With this notion we can formulate our main result:
Theorem 2.3: Assume that $`_N`$ is a stable set and that the genericity assumptions are satisfied with $`\epsilon _N`$ such that $`r_Nc_Nϵ_N0`$. Set $`t_x\tau _{_N(x)}^x`$. Then
(i) For any $`x_N`$,
$$𝔼\tau _x=R_x^1T_{x,_N(x)}(1+o(1))$$
$`(2.5)`$
(ii) For any $`x_N`$, there exists an eigenvalue $`\lambda _x`$ of $`1P_N`$ that satisfies
$$\lambda _x=\frac{1}{𝔼t_x}\left(1+o(1)\right)$$
$`(2.6)`$
Moreover, there exists a constant $`c>0`$ such that for all $`N`$
$$\sigma (1P_N)\backslash _{x_N}\lambda _x(cb_N,1]$$
$`(2.7)`$
(iii) For any $`x_N`$, for all $`t>0`$,
$$[t_x>t𝔼t_x]=e^{t(1+o(1))}(1+o(1))$$
$`(2.8)`$
(iv) If $`\psi _x`$ denotes the eigenvector of $`1P_N`$ corresponding to the eigenvalue $`\lambda _x`$, then
$$\psi _x(y)=\{\begin{array}{cc}[\tau _x^y<\tau _{_N(x)}^y](1+o(1)),\hfill & \text{if}[\tau _x^y<\tau _{_N(x)}^y]ϵ_N\hfill \\ O(ϵ_N),\hfill & \text{otherwise}\hfill \end{array}$$
$`(2.9)`$
Remark: Explicit bounds on the error terms are given \[BEGK2,EK\].
Let us make some additional comments on this theorem. First of all, they identification of what constitutes a metastable exit is crucial, and, in particular the fact that these processes include the transition through the ‘‘saddle point’’, guaranteed in our case by the insistence that the process has actually arrived in $`_N(x)`$. Without taking this into account, the precise uncertainty principle (ii) could not hold. It is interesting to note that on the level of this theorem, the difficulties associated with the control of the passage through a saddle are not visible, and that we have the exact formula $`\mathrm{}`$T.1 for the mean exit time. Of course the difficulty is hidden in the quantities $`T_{x,y}`$ whose computation is far from trivial. However, we have shown in \[BEGK1\] that at least in the reversible case, using a variational representation, very precise control of such quantities can be gained in certain setting. Somewhat less precise results can also be obtained in some non-reversible situations \[EK\]. Concerning our estimate on the eigenfunction, it is easy to see that \[BEGK2\] $`[\tau _x^y<\tau _{_N(x)}^y]`$ is either very close to one or very close to zero, except on a set of points whose invariant measure is extremely small. Therefore, the corresponding right-eigenfunctions $`\psi _x^r(z)=_N(z)\psi _x(z)`$, are essentially proportional to the measure $`_N`$ conditioned on the local valley of corresponding to $`x`$ (all up to errors of order $`ϵ_N`$), i.e. they do indeed represent metastable measures, as suggested in \[GS\].
3. Some ideas of the proofs. The first major ingredient for the proof is a representation formula for the Green’s function of the transition matrix $`P_N`$ in terms of certain probabilities. It implies in particular that for any $`I\mathrm{\Gamma }_N`$,
$$𝔼t_I^x=\underset{y\mathrm{\Gamma }_N\backslash I\backslash x}{}\frac{_N(y)}{_N(x)}[\tau _x^y<\tau _I^y][\tau _I^x<\tau _x^x]+\frac{1}{[\tau _I^x<\tau _x^x]}$$
$`(3.1)`$
This formula was first derived for the reversible case in \[BEGK1\]. An apparently independent derivation that also covers the non-reversible case was given recently in \[GM\]. This formula allows to prove in a rather simple way, $`\mathrm{}`$T.1. However, the realization that this formula actually arises from a representation of the Greens function makes it even more useful.
Our analysis of the spectrum of $`1P_N`$ passes through the analysis of the Laplace transforms, $`G_{y,I}^x(u)𝔼e^{ut_y^x}1\mathrm{I}_{\tau _y^x<\tau _I^x}`$ of transition times of process that is ‘killed’ upon arrival in a set $`I\mathrm{\Gamma }_N`$. We write $`P_N^J`$ for the transition matrix of such a process, and we write $`\lambda _J`$ for the smallest eigenvalue of $`(1P^J)`$. It then turns out that all eigenvalues of $`(1P_N)`$ below $`\lambda _J`$ can be characterised as follows: Set $`u(\lambda )\mathrm{ln}(1\lambda )`$. Define the $`|J|\times |J|`$ matrix $`𝒢_J(u)`$ whose elements are
$$\delta _{m^{},m}G_{m,J}^m^{}(u),m^{},mJ$$
$`(3.2)`$
Then $`\lambda `$ is an eigenvalue of $`(1P_N)`$ below $`\lambda _J`$, if and only if
$$detG_J(u(\lambda ))=0$$
$`(3.3)`$
This equation is rather easy to understand if $`|J|=1`$. In this case, $`\mathrm{}`$3.3 becomes simply $`G_m^m(u(\lambda )=1`$. By a simple renewal argument, one sees that $`G_x^m(u)=\frac{G_{x,m}^m(u)}{1G_{m,x}^m(u)}`$. Therefore, $`u(\lambda )`$ defined by $`\mathrm{}`$3.3 is the first value at which $`sup_{xG_N}G_x^m(u)=+\mathrm{}`$. The general formula $`\mathrm{}`$3.2 is somewhat less intuitive. Basically, one makes an ansatz for the eigenfunctions of $`(1P_N)`$ in terms of the Laplace transforms of the form
$$\mathrm{\Psi }(x)=\underset{mJ}{}\varphi _mG_{m,J}^x(u)$$
$`(3.4)`$
One then finds that condition $`\mathrm{}`$3.3 is sufficient for the ansatz to yield eigenfunctions wit $`u=u(\lambda )`$. Moreover, one can show that if $`\lambda `$ is an eigenvalue, then the eigenfunctions can be represented in this form and $`\mathrm{}`$3.3 must be satisfied.
To complete the proof one needs good control over the Laplace transforms; this is partly provided again by the representation of the Green’s function, complemented by lower bounds on eigenvalues $`\lambda _J`$ obtained from a Donsker-Varadhan \[DV\] argument. The actual proofs are rather involved and must be left to the longer publications \[BEGK2,EK\]. Let us finally mention that the good control over the spectrum of $`(1P_N)`$ allows in turn a very good control of the analytic properties of the Laplace transforms which allow in turn the sharp estimates on the probability distribution of metastable transition times stated under (iv).
References
\[BEGK1\] A. Bovier, M.Eckhoff, V. Gayrard and M. Klein, “Metastability in Stochastic Dynamics of Disordered Mean-Field Models“, WIAS-preprint 452, to appear in Prob. Theor. Rel, Fields, (2000).
\[BEGK2\] A. Bovier, M.Eckhoff, V. Gayrard and M. Klein, “Metastability and low-lying spectra in reversible Markov chains“, in perparation (2000).
\[BK\] G. Biroli and J. Kurchan, “Metastable states in glassy systems”, http://www.xxx.lanl.gov/cond-mat/0005499 (2000).
\[DV\] M.D. Donsker and S.R.S. Varadhan, “On the principal eigenvalue of second-order elliptic differential operators”, Comm. Pure Appl. Math. 29, 595-621 (1976).
\[EK\] M. Eckhoff, M. Klein, “Metastability and low lying spectra in non-reversible Markov chains”, in preparation (2000).
\[FW\] M.I. Freidlin and A.D. Wentzell, “Random perturbations of dynamical systems”, Springer, Berlin-Heidelberg-New York, 1984.
\[GM\] B. Gaveau and M. Moreau, “Metastable relaxation times and absorbtion probabilities for multidimensional stochastic systems”, J. Phys. A: Math. Gen. 33, 4837-4850 (2000).
\[GS\] B. Gaveau and L.S. Schulman, “Theory of nonequilibrium first-order phase transitions for stochastic dynamics”, J. Math. Phys. 39, 1517-1533 (1998).
\[M\] P. Mathieu, “Spectra, exit times and long times asymptotics in the zero white noise limit”, Stoch. Stoch. Rep. 55, 1-20 (1995).
\[S\] Ch. Schütte, “Conformational dynamics: modelling, theory, algorithm, and application to biomolecules”, preprint SC 99-18, ZIB-Berlin (1999).
\[SFHD\] Ch. Schütte, A. Fischer, W. Huisinga, P. Deuflhard, “A direct approach to conformational dynamics based on hybrid Monte Carlo”, J. Comput. Phys. 151, 146-168 (1999).
\[Sc\] E. Scoppola, “Renormalization and graph methods for Markov chains”. Advances in dynamical systems and quantum physics (Capri, 1993), 260-281, World Sci. Publishing, River Edge, NJ, 1995.
\[W\] A.D. Wentzell, “On the asymptotic behaviour of the greatest eigenvalue of a second order elliptic differential operator with a small parameter in the higher derivatives”, Soviet Math. Docl. 13, 13-17 (1972).
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# Analytic Expression for Exact Ground State Energy Based on an Operator Method for a Class of Anharmonic Potentials
## I Introduction
The operator method has been widely used as an elegant analytical tool in quantum mechanics and quantum statistics for studying exactly solvable models. Since the publication of a remarkable paper by Delbecq and Quesne , the operator method has been extended to many other possible physical applications and realizations . Indeed, the generalization of the notion of “shift operators” or “ladder operators” as a spectrum generating algebra can be effectively studied both mathematically and physically in terms of “nonlinear algebra” introduced by Delbecq and Quesne. In fact, Chen et al has recently reproduced the results of Bethe Ansatz for the XXX model through the shift operator method. A general observation in these works appears to be a possibility of treating the well-known technique of the algebraic Bethe Ansatz for quantum many-body problem as a special case of the shift operator method. Thus, all previous works evaluated using the algebraic Bethe Ansatz should hopefully be possible using a nonlinear algebraic method via shift operators. A further implication of this notion is the emergence of an underlying Yangian algebra. Furthermore, the factorization method employed in supersymmetric quantum mechanics (SQM) can also be deemed as a special case of the nonlinear Delbecq-Quesne algebra.
In many realizations of physical models, the operator method can provide us with an alternative and a clearer picture regarding the analytical determination of energy eigenstates and eigenvalues. Specifically, it has been found to be a useful tool in the solution of bound state problems such as the string postulate in Bethe Ansatz or the raising and lowering of energy levels in SQM. As in Bethe Ansatz and SQM, knowledge of a certain reference state, like the highest weight state in Bethe Ansatz or the ground state ansatz for the superpartner potential in SQM, in the operator method permits the analysis and determination of the full spectrum of a physical system. Moreover, as mentioned in ref. , the knowledge of the ground state wave function also determines the potential and hence the Hamiltonian of a system up to a constant value. However, in principle given the Hamiltonian of a physical system, surely it should be possible to determine all the energy levels. In this paper, we show that provided certain consistency relations are satisfied, the entire spectrum of a given Hamiltonian can be determined by shift operator method even if there is no prior knowledge of a reference state.
To illustrate the power of the shift operator method, we apply the technique to one-dimensional oscillator with anharmonic potentials. The study of anharmonic potential has always been an exciting and interesting field due to its broad applications in quantum field theory , nuclear models, atomic and molecular physics , condensed matter physics, statistical physics and chemical physics . Indeed, numerous numerical methods including renormalized strong coupling expansion, perturbation expansion, supersymmetric quantum mechanics, WKB, iteration based on the generalized Bloch equation, state-dependent diagonalization, Hill determinant method, phase-integral approach, iterative Bogoliubov transformations, eigenvalue moment method, perturbative-variation, and algebraic method have been proposed to investigate these anharmonic potentials . However, due to its inherent intractability, no analytic solution for energy spectrum and eigenstates has been obtained so far except in some special cases. Indeed, for the first time, we have partially resolved some of these technical issues and obtained analytic solutions for a wide class of anharmonic potentials under certain conditions.
It is instructive to see how SQM can be regarded as a special case of the nonlinear algebraic operator method. Before we proceed further to show this point, we first recall some relevant definitions in operator method and nonlinear algebra . Let $`\widehat{H}`$ be an observable satisfying the eigenfunction equation $`\widehat{H}|\psi >=E|\psi >`$. The operators $`L^\pm `$ satisfying
$$[\widehat{H},L^\pm ]=L^\pm f^\pm (\widehat{H}),$$
(1)
are called “shift operators” of $`\widehat{H}`$. In Eq.(1), $`f^\pm (\widehat{H})`$ are real functions of $`\widehat{H}`$ so that the values of its action on eigenstates of $`H`$ are reals. Thus, these values can be interpreted as energy gaps if $`\widehat{H}`$ is the Hamiltonian of the system. Nevertheless, it is not necessary for the mutually adjoint condition
$$(L^+)^{}=L^{}$$
to hold and there is generally no constraint on the commutation relation between $`L^+`$ and $`L^{}`$, i.e. $`[L^+,L^{}]`$. However, if the mutually adjoint condition is satisfied, we obtained the “nonlinear algebra” defined by the relations
$`[\widehat{H},L^+]`$ $`=`$ $`f(\widehat{H})L^+,`$ (2)
$`[\widehat{H},L^{}]`$ $`=`$ $`L^{}f(\widehat{H}),`$ (3)
$`[L^{},L^+]`$ $`=`$ $`g(\widehat{H}).`$ (4)
A simple rearrangement of the above relations gives
$`(\widehat{H}f(\widehat{H}))L^+`$ $`=`$ $`L^+\widehat{H},`$ (5)
$`L^{}(\widehat{H}f(\widehat{H}))`$ $`=`$ $`\widehat{H}L^{}.`$ (6)
It is interesting to compare Eq.(6) with the analogous relations in SQM
$`\widehat{H}_1Q^+`$ $`=`$ $`Q^+\widehat{H}_2`$ (7)
$`Q^{}\widehat{H}_1`$ $`=`$ $`\widehat{H}_2Q^{}`$ (8)
$`Q^{}`$ $`=`$ $`(Q^+)^{}`$ (9)
and identify
$`Q^+L^+`$ $`Q^{}L^{}`$ (10)
$`\widehat{H}_1\widehat{H}f(\widehat{H})`$ $`\widehat{H}_2\widehat{H}.`$ (11)
Hence, one can always regard SQM as a specific type of nonlinear algebraic method.
The purpose of this paper is to describe the operator method and apply it to study the one-dimensional oscillator with anharmonic potentials. In section II, we show how a shift operator method can be formulated using a simple instructive example of the harmonic oscillator. The general procedures to solve anharmonic potentials are presented in section III. In section IV, we solve the relevant equations for the one-dimensional oscillator with cubic-quartic potential and show that we can get the exact analytical expression for the ground state energy provided certain consistency relations are satisfied. We then compare our analytical result with some numerical results in section V. Our analytical expressions can be used to verify the accuracy of existing numerical results. We conclude with some brief remarks in section VI.
## II Applying Operator Method to the Harmonic Oscillator
In this section, we apply the operator method to the harmonic oscillator. The Hamiltonian of the harmonic oscillator is given by
$$\widehat{H}=\frac{1}{2}(\widehat{p}^2+\widehat{x}^2).$$
(12)
Together with the well-known commutation relation between $`\widehat{x}`$ and $`\widehat{p}`$, i.e. $`[\widehat{x},\widehat{p}]=i`$ ($`\mathrm{}=1`$), we get
$`[\widehat{H},\widehat{x}]`$ $`=`$ $`i\widehat{p}`$ (13)
$`[\widehat{H},\widehat{p}]`$ $`=`$ $`i\widehat{x}.`$ (14)
Using the notation
$$[\widehat{H},(\widehat{x},\widehat{p})]:=([\widehat{H},\widehat{x}],[\widehat{H},\widehat{p}]),$$
Eq.(13) can also be rewritten as
$`[\widehat{H},(\widehat{x},\widehat{p})]`$ $`=`$ $`(\widehat{x},\widehat{p})\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),`$ (17)
$``$ $`(\widehat{x},\widehat{p})M`$ (18)
where, the matrix, $`M`$, is called the “coefficient matrix”. To diagonalize the coefficient matrix, we look for a transformation, $`U`$, such that $`M=UDU^1`$ where $`D`$ is a diagonal matrix. In general, provided all the entries of $`U`$ commute with Hamiltonian $`\widehat{H}`$, the entries of the matrix $`U`$ are functions of $`c`$ numbers, the Hamiltonian or any conserved quantity $`\widehat{I}`$ of the system. Thus, we have
$`[\widehat{H},(\widehat{x},\widehat{p})U(c,\widehat{H},\widehat{I})]`$ $`=`$ $`(\widehat{x},\widehat{p})U(c,\widehat{H},\widehat{I})U^1(c,\widehat{H},\widehat{I})\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)U(c,\widehat{H},\widehat{I})`$ (21)
$`=`$ $`(\widehat{x},\widehat{p})U(c,\widehat{H},\widehat{I})\left(\begin{array}{cc}\lambda _1(c,\widehat{H},\widehat{I})& 0\\ 0& \lambda _2(c,\widehat{H},\widehat{I})\end{array}\right)`$ (24)
It is easy to see that
$$U=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ i& i\end{array}\right),$$
so that Eq.(17) can be rewritten as
$`[\widehat{H},(\widehat{x},\widehat{p}){\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}1& 1\\ i& i\end{array}\right)]`$ (27)
$`=`$ $`(\widehat{x},\widehat{p}){\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}1& 1\\ i& i\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).`$ (32)
Identifying
$$(\widehat{a}^+,\widehat{a})=(\widehat{x},\widehat{p})\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ i& i\end{array}\right),$$
we obtain
$$\begin{array}{ccc}[\widehat{H},\widehat{a}^+]& =& \widehat{a}^+\\ [\widehat{H},\widehat{a}]& =& \widehat{a}.\end{array}.$$
The “creation” and the “annihilation” operators are thus obtained naturally. Following the literature, we can also call them “shift” or “ladder” operators. Although, the application of the shift operator to harmonic oscillator is very simple, it can still serve an illustrative example.
For the harmonic oscillator, a closed algebra is obtained. This closure in the algebra enables the whole spectrum of the system to be determined using only one pair of shift operators, in which one of the operator raises while the other lowers the energy levels. In general, the existence of a closed algebra is not assured, and so, we usually have to deal with an unclosed algebra with infinitely many shift operators. Thus, we have to use infinitely many shift operators pairs to generate the whole spectrum. Moreover, corresponding to each energy level, we have a pair of distinct shift operators in which one raises while the other lowers the energy level.
It is interesting to note that, the eigenvalues of the coefficient matrix correspond to the amount of shifted energy associated with the various shift operators. And, except for some special Hamiltonians, the coefficient matrix, its eigenvalues and the transformation matrix, are all functions of the Hamiltonian. In nonlinear algebraic method, the energy shifts are no longer uniform and the gap between any two adjacent energy levels is related to their relative positions in the spectrum. Thus, we find that shift operator method can provide us with a clearer physical picture for the mapping of the energy level in a physical system.
Note that we can only extract the energy gaps rather than the energy levels. Moreover, we have no information regarding to the energy of the ground state. However, in some cases, it is possible to determine the exact energy of the ground state. This last result constitutes the gist of one of the most important aspect of our paper.
## III Operator Method to Anharmonic Oscillator: the General Procedure
In this section, we confine ourselves to the one-dimensional harmonic oscillator in which the anharmonic potential contains cubic and quartic terms. Nevertheless, the procedure developed here is very general and can be applied to other types of anharmonic potentials. The Hamiltonian is given by
$$\widehat{H}=\frac{d^2}{dx^2}+\alpha \widehat{x}^2+\beta \widehat{x}^3+\gamma \widehat{x}^4,$$
(33)
and it is easy to show that
$`[\widehat{H},\widehat{x}^n]`$ $`=`$ $`2n\widehat{x}^{n1}{\displaystyle \frac{d}{dx}}n(n1)\widehat{x}^{n2}`$ (34)
$`[\widehat{H},\widehat{x}^n{\displaystyle \frac{d}{dx}}]`$ $`=`$ $`n(n1)\widehat{x}^{n2}{\displaystyle \frac{d}{dx}}+2n\widehat{x}^{n1}\widehat{H}`$ (36)
$`2\alpha (n+1)\widehat{x}^{n+1}2\beta (n+{\displaystyle \frac{3}{2}})\widehat{x}^{n+2}2\gamma (n+2)\widehat{x}^{n+3}.`$
Rewriting Eqs.(34,36) into matrix form, we have
$`[\widehat{H},(\widehat{x},\widehat{x}^2,\widehat{x}^3,\mathrm{},\widehat{x}^n,\mathrm{})]`$ $`=`$ $`(\widehat{x},\widehat{x}^2,\widehat{x}^3,\mathrm{},\widehat{x}^n,\mathrm{})M_1`$ (38)
$`+({\displaystyle \frac{d}{dx}},x{\displaystyle \frac{d}{dx}},x^2{\displaystyle \frac{d}{dx}},\mathrm{},x^n{\displaystyle \frac{d}{dx}},\mathrm{})N_1+L_1`$
$`[\widehat{H},({\displaystyle \frac{d}{dx}},x{\displaystyle \frac{d}{dx}},x^2{\displaystyle \frac{d}{dx}},\mathrm{},x^n{\displaystyle \frac{d}{dx}},\mathrm{})]`$ $`=`$ $`(\widehat{x},\widehat{x}^2,\widehat{x}^3,\mathrm{},\widehat{x}^n,\mathrm{})M_2`$ (40)
$`+({\displaystyle \frac{d}{dx}},x{\displaystyle \frac{d}{dx}},x^2{\displaystyle \frac{d}{dx}},\mathrm{},x^n{\displaystyle \frac{d}{dx}},\mathrm{})N_2+L_2`$
where
$$M_1=\left(\begin{array}{cccccccc}0& 0& 6& 0& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 12& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& 20& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& 0& \mathrm{}& n(n1)& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right),$$
(41)
$$N_1=\left(\begin{array}{cccccc}2& 0& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 4& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 6& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 2n& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right),$$
(42)
$$L_1=(0,2,\mathrm{\hspace{0.33em}\hspace{0.33em}0},\mathrm{\hspace{0.33em}\hspace{0.33em}0},\mathrm{},\mathrm{\hspace{0.33em}\hspace{0.33em}0},\mathrm{})$$
(43)
$$M_2=\left(\begin{array}{cccccccc}2\alpha & 0& 4\widehat{H}& 0& 0& & \mathrm{}& \mathrm{}\\ 3\beta & 4\alpha & 0& 4\widehat{H}& 0& & \mathrm{}& \mathrm{}\\ 4\gamma & 5\beta & 6\alpha & 0& 8\widehat{H}& & \mathrm{}& \mathrm{}\\ 0& 6\gamma & 7\beta & 8\alpha & 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 8\gamma & 9\beta & 10\alpha & \mathrm{}& 2n\widehat{H}& \mathrm{}\\ 0& 0& 0& 10\gamma & 11\beta & \mathrm{}& 0& \mathrm{}\\ 0& 0& 0& 0& 12\gamma & \mathrm{}& 2(n+1)\alpha & \mathrm{}\\ 0& 0& 0& 0& 0& \mathrm{}& 2(n+\frac{3}{2})\beta & \mathrm{}\\ 0& 0& 0& 0& 0& \mathrm{}& 2(n+2)\gamma & \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}\end{array}\right),$$
(44)
$$N_2=\left(\begin{array}{cccccccc}0& 0& 2& 0& 0& & \mathrm{}& \mathrm{}\\ 0& 0& 0& 6& 0& & \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& 12& & \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& 0& \mathrm{}& n(n+1)& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}\end{array}\right),$$
(45)
and
$$L_2=(0,\mathrm{\hspace{0.33em}\hspace{0.33em}2}\widehat{H},\mathrm{\hspace{0.33em}\hspace{0.33em}0},\mathrm{\hspace{0.33em}\hspace{0.33em}0},\mathrm{},\mathrm{\hspace{0.33em}\hspace{0.33em}0},\mathrm{})$$
(46)
It is interesting to note that for other types of potentials, all matrices except for $`M_2`$ are exactly the same.
From Eqs.(38,40), we find
$$\begin{array}{ccc}[\widehat{H},(\widehat{x},\widehat{x}^2,\widehat{x}^3,\mathrm{},\widehat{x}^n,\mathrm{})R+(\frac{d}{dx},x\frac{d}{dx},x^2\frac{d}{dx},\mathrm{},x^n\frac{d}{dx},\mathrm{})S]\hfill & & \\ =(\widehat{x},\widehat{x}^2,\widehat{x}^3,\mathrm{},\widehat{x}^n,\mathrm{})(M_1R+M_2S)+(\frac{d}{dx},x\frac{d}{dx},x^2\frac{d}{dx},\mathrm{},x^n\frac{d}{dx},\mathrm{})(N_1R+N_2S)+(L_1R+L_2S)\hfill & & \end{array}.$$
(47)
Thus, if the following relations exist
$`M_1R+M_2S`$ $`=`$ $`RT`$ (48)
$`N_1R+N_2S`$ $`=`$ $`ST`$ (49)
then we have
$$\begin{array}{c}[\widehat{H},(\widehat{x},\widehat{x}^2,\widehat{x}^3,\mathrm{},\widehat{x}^n,\mathrm{})R+(\frac{d}{dx},x\frac{d}{dx},x^2\frac{d}{dx},\mathrm{},x^n\frac{d}{dx},\mathrm{})S]\hfill \\ =\{(\widehat{x},\widehat{x}^2,\widehat{x}^3,\mathrm{},\widehat{x}^n,\mathrm{})R+(\frac{d}{dx},x\frac{d}{dx},x^2\frac{d}{dx},\mathrm{},x^n\frac{d}{dx},\mathrm{})S\}T+(L_1R+L_2S)\hfill \end{array}.$$
(50)
To diagonalize the coefficient matrix $`T`$, let us suppose that $`U`$ is the transformation matrix needed. The eigenvalues and the corresponding shift operators can therefore be written as
$$\mathrm{\Lambda }=U^1TU=\text{diag}(\lambda _1,\lambda _2,\lambda _3,\mathrm{},\lambda _n,\mathrm{})$$
(51)
and
$$\begin{array}{c}(\widehat{A}_1,\widehat{A}_2,\widehat{A}_3,\mathrm{},\widehat{A}_n,\mathrm{})\hfill \\ =\{(\widehat{x},\widehat{x}^2,\widehat{x}^3,\mathrm{},\widehat{x}^n,\mathrm{})R+(\frac{d}{dx},x\frac{d}{dx},x^2\frac{d}{dx},\mathrm{},x^n\frac{d}{dx},\mathrm{})S\}U+(L_1R+L_2S)U\mathrm{\Lambda }^1\hfill \end{array}$$
(52)
respectively. In fact, Eq.(50) can be written into a more succinct form as
$$[\widehat{H},(\widehat{A}_1,\widehat{A}_2,\widehat{A}_3,\mathrm{},\widehat{A}_n,\mathrm{})]=(\widehat{A}_1,\widehat{A}_2,\widehat{A}_3,\mathrm{},\widehat{A}_n,\mathrm{})\left(\begin{array}{cccccc}\lambda _1& 0& 0& & \mathrm{}& \mathrm{}\\ 0& \lambda _2& 0& & \mathrm{}& \mathrm{}\\ 0& 0& \lambda _3& & \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& \lambda _n& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}\end{array}\right),$$
(53)
or
$`[\widehat{H},\widehat{A}_1]`$ $`=`$ $`\widehat{A}_1\lambda _1`$
$`[\widehat{H},\widehat{A}_2]`$ $`=`$ $`\widehat{A}_2\lambda _2`$
$`\mathrm{}`$ $`\mathrm{}`$ $`\mathrm{}`$
$`[\widehat{H},\widehat{A}_n]`$ $`=`$ $`\widehat{A}_n\lambda _n`$
$`\mathrm{}`$ $`\mathrm{}`$ $`\mathrm{}.`$
Thus, in general we get an unclosed algebra with infinitely many shift operators. In case of the usual harmonic oscillator, we have a pair of shift operators which can generate the entire spectrum. In the present case, we have infinitely many pairs of shift operators in which each pair is responsible for raising and lowering the corresponding energy level. The whole spectrum can only be generated by the infinite set of the shift operators acting on the ground state. Hence, it is natural to see that all the eigenvalues of the coefficient matrix are only dependent on the energy of ground state. That is, $`\lambda _i,(i=1,2,3,\mathrm{},n,\mathrm{})`$ are the functions of the ground state energy. Therefore, we can identify the operator $`\widehat{H}`$ in the coefficient matrix $`T`$ as the ground state energy! Based on this observation, we can then get the analytic expression for the energy of the ground state provided certain consistency relations hold.
It remains to solve for the matrices $`R,S`$ and $`T`$ given the matrices $`M_1,M_2,N_1`$ and $`N_2`$. Once we have obtained the matrices $`R,S,T`$ and diagonalize $`T`$, we effectively obtained all the shift operators and we can then reconstruct the full spectrum. It is instructive to note that one of the matrices $`R`$ or $`S`$ is redundant and we set it to unity. For convenience, let us set $`S`$ in Eqs.(48,49) to unity so that the problem is reduce to the solution of $`R,T`$ from the following equations
$`M_1R+M_2`$ $`=`$ $`RT`$ (54)
$`N_1R+N_2`$ $`=`$ $`T`$ (55)
and the diagonalization of $`T`$. Note that although we have considered a specific form for the potential, the procedure developed here is very general and can be applied to more complicated cases.
## IV Analytic Expression for the Ground State Energy: a Class of Anharmonic Potentials
To solve Eq.(54) and Eq.(55), we first make the following observation. If we define
$$G=\left(\begin{array}{cccccc}1& 0& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 2& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 3& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& n& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right),$$
(56)
$$P=\left(\begin{array}{ccccccc}0& 1& 0& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 1& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 1& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& \mathrm{}& 1& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)Q=\left(\begin{array}{cccccc}0& 0& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 1& 0& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 1& 0& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 1& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 1& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right),$$
(57)
we have
$$\begin{array}{ccc}M_1\hfill & =\hfill & PGPG\hfill \\ M_2\hfill & =\hfill & 2\alpha G+2\widehat{H}PGP2\gamma QGQ\beta (2G1)Q\hfill \\ N_1\hfill & =\hfill & 2G\hfill \\ N_2\hfill & =\hfill & GPGP.\hfill \end{array}$$
(58)
Here, the matrices $`G,P`$ and $`Q`$ are nothing but representations of particle number, creation and annihilation operators for the usual harmonic oscillator in Fock space. It is easy to check that
$$[G,P]=P,[G,Q]=Q,PQ=1,QP=diag(0,1,1,\mathrm{},1,\mathrm{}).$$
(59)
Furthermore, if we set
$$W=N_1R,$$
then Eq.(55) becomes
$$T=W+N_2R=N_1^1W.$$
From Eq.(58), we notice that $`N_1M_1N_1^1=N_2`$, and as a consequence, we can combine Eq.(54) and the above expression for the matrix $`T`$, to yield the expression
$$W^2+[W,N_2]N_1M_2=0.$$
(60)
Unfortunately, due to the infinite dimensionality of the matrices, the solution, except in some special cases, is not known.
In the case of the one-dimensional oscillator with cubic-quartic term, if we let $`W`$ assume the following form
$$W=2(aG+bGQ+cGP),$$
it is not difficult to show using Eq.(59) that
$`W^2+[W,N_2]N_1M_2`$ $`=`$ $`4\{(a^2+2bc\alpha )G^2+(ab{\displaystyle \frac{1}{2}}\beta )(2G^2G)Q+(b^2\gamma )GQGQ`$ (62)
$`(ac+b)(2G^2+G)P+(c^2+a+\widehat{H})GPGP\}`$
so that consistency naturally requires the following conditions
$`a^2+2bc\alpha `$ $`=`$ $`0`$ (63)
$`ab{\displaystyle \frac{1}{2}}\beta `$ $`=`$ $`0`$ (64)
$`b^2\gamma `$ $`=`$ $`0`$ (65)
$`ac+b`$ $`=`$ $`0`$ (66)
$`c^2+a+\widehat{H}`$ $`=`$ $`0`$ (67)
in order that Eq.(60) be satisfied. In particular, Eq.(67) gives the ground state energy of the system.
As explained in the previous section, we have equated $`\widehat{H}`$ to one of its eigenvalues, namely the ground state energy. Thus, in order that the full spectrum be generated from the infinitely many raising operators acting on the ground state, so that each raising operator generates its own corresponding energy level via its action on the ground state, the consistency relations in Eq.(63)-Eq.(67) must hold.
It is easy to see that the solution to Eq.(63)-Eq.(67) is
$$a=\frac{\beta }{2\sqrt{\gamma }}b=\sqrt{\gamma }c=2\frac{\gamma }{\beta }$$
(68)
$$\widehat{H}=4\frac{\gamma ^2}{\beta ^2}\frac{\beta }{2\sqrt{\gamma }}$$
(69)
$$\frac{\beta ^2}{4\gamma }4\frac{\gamma ^{3/2}}{\beta }\alpha =0$$
(70)
and
$$a=\frac{\beta }{2\sqrt{\gamma }}b=\sqrt{\gamma }c=2\frac{\gamma }{\beta }$$
(71)
$$\widehat{H}=4\frac{\gamma ^2}{\beta ^2}+\frac{\beta }{2\sqrt{\gamma }}$$
(72)
$$\frac{\beta ^2}{4\gamma }+4\frac{\gamma ^{3/2}}{\beta }\alpha =0.$$
(73)
The additional constraint Eq.(70) (or Eq.(73)) seems to indicate that the ansatz for $`W`$ may be too simplistic. However, we do not have a direct solution of $`W`$ from Eq.(60) for the given $`N_1,N_2`$ and $`M_2`$. Despite all these, we can still obtain an analytic result for exact energy of the ground state for the given anharmonic potential and this result can prove to be valuable for analyzing the accuracy of existing numerical methods. Besides, as mentioned before, it is anticipated that the operator method can provide us with an analytic tool for investigating the ground state, something which is not possible using existing numerical approaches.
To compute the energy levels of the excited states, we need to diagonalize the following infinite dimensional matrix
$$T=\left(\begin{array}{cccccccc}2a& 2c& 2& 0& 0& & \mathrm{}& \mathrm{}\\ 4b& 4a& 4c& 6& 0& & \mathrm{}& \mathrm{}\\ 0& 6b& 6a& 6c& 12& & \mathrm{}& \mathrm{}\\ 0& 0& 8b& 8a& 8c& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 10b& 10a& \mathrm{}& n(n+1)& \mathrm{}\\ 0& 0& 0& 0& 12b& \mathrm{}& 2nc& \mathrm{}\\ 0& 0& 0& 0& 0& \mathrm{}& 2na& \mathrm{}\\ 0& 0& 0& 0& 0& & 2nb& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}\end{array}\right)$$
with $`a,b`$ and $`c`$ are given by Eq.(68) or Eq.(71). Once we have diagonalized this matrix, using the result of the ground state energy, all the other energy levels can be obtained. Unfortunately, we do not know any method for diagonalizing it at present moment despite the apparent simplicity and symmetry in the matrix.
## V Comparison of Numerical with the Analytic Results for Ground State Energy
In this section, we use our analytic results to check against previous numerical computations. For numerical simulation, we use the method proposed by Ho et al. using the state-dependent diagonalization method. As claimed in ref. , this method is very accurate and efficient compare to other numerical methods for calculating the energy eigenvalues and eigenfunctions of the one-dimensional harmonic oscillator with anharmonic potentials.
Due to the consistency requirements in Eq.(70) or Eq.(73), the three coefficients $`\alpha ,\beta `$ and $`\gamma `$ in
$$V(x)=\alpha x^2+\beta x^3+\gamma x^4$$
are not mutually independent. It is not difficult to solve for $`\beta `$ using Eq.(70) or Eq.(73) for given values of $`\alpha `$ and $`\gamma `$. The final solution is
$$\beta =2\sqrt{\gamma }\left[\frac{\frac{\alpha }{3}}{\left(\gamma +\sqrt{\gamma ^2(\frac{\alpha }{3})^3}\right)^{\frac{1}{3}}}+\left(\gamma +\sqrt{\gamma ^2(\frac{\alpha }{3})^3}\right)^{\frac{1}{3}}\right]$$
(74)
corresponding to Eq.(70) and $`\beta `$ to Eq.(73) with the ground state energy being same in both cases. Specifically, the ground state energy is given by
$$E_0=4\frac{\gamma ^2}{\beta ^2}\frac{\beta }{2\sqrt{\gamma }}.$$
(75)
The dependence of $`\beta `$ on $`\alpha `$ and $`\gamma `$ is shown in Fig.(1)
Setting $`\alpha =1,2`$, the values of $`E_0`$ as a function of $`\gamma `$ for the analytic formula and the numerical simulation using the state dependent diagonalization are shown respectively in Fig.(2) and Fig.(3)) with the corresponding potential for particular values of $`\gamma `$ as a function of $`x`$ shown in Fig.(4)). The difference in the energy computed from the numerical simulation and the analytic formula as a function of $`\gamma `$ is also plotted in Fig.(5) and Fig.(6) for $`\alpha =1`$ and $`\alpha =2`$ respectively. As shown in the figures, the numerical simulation using state-dependent diagonalization is in excellent agreement with our analytical formula. Moreover, we see that the state-dependent-diagonalization method provides sufficiently high accuracy for all practical purposes.
## VI Conclusions
In summary,we have formulated a general procedure for extracting the full spectrum of a physical system with arbitrary potentials using the operator method. We have applied it to solve the one-dimensional harmonic oscillator with anharmonic problem. The analytic expression for the ground state energy is obtained for a large class of anharmonic potentials. Our results can be used to verify the accuracy of existing numerical methods. However, in order to get the full spectrum, we need to diagonalize an infinite dimensional matrix. This last problem is to our knowledge an open mathematical question except for some very special cases.
Our method has also confirmed the notion that once the Hamiltonian of certain system is given, all the energy levels can be determined through the operator method even without any prior knowledge of a reference state, provided certain consistency relations hold, a feat considered impossible hitherto.
Finally, from Eq.(74) and Eq.(75), it easy to see that in the limit when $`\gamma \mathrm{}`$, $`E_0\frac{3}{2}(2\gamma )^{1/3}`$, confirming earlier analysis on the mathematical properties of ground-state energy obtained in ref. .
## VII Acknowledgment
This work has been supported by NUS Research Grant No. RP3982713. In addition, we would also like to thank the anonymous referee for his invaluable remarks and comments.
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# 1 (a) Confining double-well potential in the longitudinal direction of the coupled quantum dot structure; (b) Two-electron ground state versus the longitudinal coordinates of the two electrons. Notice that the ground state is delocalized, ie. both electrons occupy different dots.
Localization and entanglement of two interacting electrons in a quantum-dot molecule
P. I. Tamborenea and H. Metiu
Center for Quantized Electronic Structures (QUEST) and Department of Chemistry
University of California
Santa Barbara, CA 93106-9510
The localization of two interacting electrons in a coupled-quantum-dots semiconductor structure is demonstrated through numerical calculations of the time evolution of the two-electron wave function including the Coulomb interaction between the electrons. The transition from the ground state to a localized state is induced by an external, time-dependent, uniform electric field. It is found that while an appropriate constant field can localize both electrons in one of the wells, oscillatory fields can induce roughly equal probabilities for both electrons to be localized in either well, generating an interesting type of localized and entangled state. We also show that shifting the field suddenly to an appropriate constant value can maintain in time both types of localization.
PACS: 73.23.Hk, 73.61.-r, 78.66.-w, 78.47.+p
Coherent control of quantum systems is at the heart of promising disciplines like femtochemistry and quantum information processing. A basic operation of quantum control, namely, the localization of a single electron in coupled quantum wells, has been extensively studied in the last decade. In two early publications, conditions to maintain existing localization with an AC field, and to create and maintain localization with a semi-infinite AC field were identified. Thereafter, localization in two-level systems, multilevel systems, induced by ultrashort laser pulses, in dissipative environments, in molecular systems, in trapped Bose-Einstein condensates, by means of circularly polarized fields, with bichromatic fields, including the effect of Coulomb charging energy, and with quantized electromagnetic fields, has been studied.
When two or more interacting particles are present the possibility of entanglement of the many-body wave function arises. Entanglement is an essential ingredient in any scheme of quantum information processing like quantum cryptography and quantum computation, and therefore it is a problem of great current interest to find or design systems where entanglement can be manipulated.
In the case of a single-electron, two-level system consisting of the lowest symmetric and antisymmetric states of the double well potential, perfect localization can be achieved with a strong periodic electric field that causes the two Floquet quasi-energies to be degenerate. This scheme is not applicable to the two-electron system, because it requires a pair of states whose superposition results in a localized state, a condition that is not met in this system. An alternative and trivial way of inducing localization in a two-electron double-well system is to adiabatically tilt the potential with a slowly increasing electric field. Trying to localize the electrons on a short time scale is more difficult: switching the field on rapidly excites higher electronic states causing a coherent motion from one well to another. In this Letter we examine this regime of very fast localization.
We investigate the localization and entanglement of two interacting electrons in a system of coupled quantum dots, induced by spatially uniform electric fields with a simple time dependence. Systems of coupled quantum dots, sometimes referred to as quantum-dot molecules or artificial molecules, have been actively investigated in the last five years, both experimentally and theoretically. Our main findings are as follows: (i) A constant electric field can bring the two electrons from their ground state (highly delocalized) to a state of high degree of localization in one of the dots, at a certain time; (ii) At that time, a step in the field to another constant value can maintain the localization essentially indefinitely; (iii) An oscillatory field produces a different type of localization, where both electrons are likely to be found together in either dot with roughly equal probabilities. This type of localization is a purely many-body phenomenon, and arises due to the Coulomb interaction between the electrons, which causes entanglement of the two-body wave function; (iv) The entangled/localized states can also be maintained in time by changing from the oscillatory field to an appropriate constant field. In all cases, localization takes place in a time scale of a few picoseconds.
Our system is a quasi one-dimensional, double-quantum-dot structure with two electrons in it. The transversal size of the dots is taken to be $`L=50\text{Å}`$, and the double-well potential in the longitudinal direction, $`V(z)`$, is shown in Fig. 1(a). The energies associated with the transverse dimensions are, due to the narrow lateral confinement, high compared to those of the longitudinal motion. Therefore, the lateral degrees of freedom do not participate in the dynamics and the two-electron wave function can be written as (we discuss below the spin part of the wave function)
$$\mathrm{\Psi }(𝐫_1,𝐫_2,t)=\varphi (x_1)\varphi (y_1)\varphi (x_2)\varphi (y_2)\mathrm{\Phi }(z_1,z_2,t).$$
(1)
where $`\varphi (x)=\sqrt{2/L}\mathrm{sin}(\pi x/L)`$. The time-dependent Schrödinger equation becomes
$`i\mathrm{}{\displaystyle \frac{\mathrm{\Phi }}{t}}`$ $`=`$ $`[{\displaystyle \frac{\mathrm{}^2}{2m^{}}}({\displaystyle \frac{^2}{z_1^2}}+{\displaystyle \frac{^2}{z_1^2}})+V(z_1)+V(z_2)`$ (2)
$`+`$ $`V_{1D}(|z_1z_2|)e(z_1+z_2)E(t)]\mathrm{\Phi },`$ (3)
where $`E(t)`$ is an external time-dependent electric field, and $`m^{}`$ is the effective mass. $`V_{1D}`$ is the Coulomb interaction given by
$`V_{1D}(|z_1z_2|)`$ $`=`$ $`{\displaystyle _0^L}𝑑x_1𝑑y_1𝑑x_2𝑑y_2`$ (5)
$`{\displaystyle \frac{e^2\varphi ^2(x_1)\varphi ^2(y_1)\varphi ^2(x_2)\varphi ^2(y_2)}{ϵ|𝐫_1𝐫_2|}}.`$
We use the effective mass $`m^{}`$ and dielectric constant $`ϵ`$ of GaAs.
In all our calculations the ground state is the initial state. The spin part of the ground state of the two interacting electrons is the singlet state, and correspondingly (the fermionic wave function is antisymmetric) the spatial part of the wave function is symmetric under particle exchange. We calculate the spatial part of the ground state by numerical diagonalization of the energy eigenvalue problem of the interacting electrons. A plot of the ground state is shown in Fig. 1(b) as a function of the $`z`$-coordinates of the two electrons. Since the Hamiltonian of the system including the external electric field is spin independent, the spin wave function is the singlet at all times. Therefore, the evolving spatial wave function remains symmetric (under particle exchange) at all times. We emphasize that the absence of triplet states in the expansion of the wave function (Eq. (6) below) is not an approximation, but a consequence of our choice of initial state and the lack of spin dependence in the Hamiltonian.
We calculate the evolution of the two-electron wave function by using the configuration interaction (CI) method, in which one expands the wave function as
$$\mathrm{\Phi }(z_1,z_2,t)=\underset{i,j}{}c_{ij}(t)[\phi _i(z_1)\phi _j(z_2)+\phi _j(z_1)\phi _i(z_2)],$$
(6)
where $`\phi _i(z)`$ are eigenstates of $`P_z^2/2m^{}+V(z)`$. The sum runs over $`1ij`$. We take $`j=1,\mathrm{},N`$, where $`N=12`$ is the number of bound states of the potential $`V(z)`$ shown in Fig. 1(a). The two-particle basis set of symmetric states has $`N(N+1)/2=78`$ states. We find that this basis set is large enough to achieve convergence. To obtain the time evolution, the expansion (6) is substituted in the Schrödinger equation (3) and the coefficients $`c_{ij}(t)`$ are calculated numerically with the fourth-order Runge-Kutta method.
In order to describe the localization of the electrons, we introduce the probabilities
$$P_{RL}(t)=2_R𝑑z_1_L𝑑z_2|\mathrm{\Phi }(z_1,z_2,t)|^2,$$
(7)
that one electron is in the right and the other one is in the left dot,
$$P_{RR}(t)=_R𝑑z_1_R𝑑z_2|\mathrm{\Phi }(z_1,z_2,t)|^2,$$
(8)
that both electrons are in the right dot, and $`P_{LL}(t)`$, that both electrons are in the left dot. Since the probability of ionization is kept small at all times $`P_{RL}(t)+P_{RR}(t)+P_{LL}(t)1`$.
For the ground state, shown in Fig. 1(b), $`P_{RL}0.9988`$ and $`P_{RR}=P_{LL}0.0006`$. That is to say that in the ground state the two electrons have a very small probability to be found in the same well, which is expected due to their Coulomb repulsion. In the rest of this Letter we explore the question of whether an external electric field (time dependent but spatially uniform) can induce localization of the two electrons on a very fast (picosecond) time scale.
We start the search for localization with the simplest case, a constant electric field, $`E(t)=E_0`$. Before $`t=0`$ the system is in the ground state, and at $`t=0`$ the field $`E_0`$ is switched-on suddenly. We compute the evolution of the two-electron wave function as well as the probability $`P_{RL}(t)`$ during a simulation interval of 9 ps. For each field $`E_0`$ we find the lowest value of $`P_{RL}`$ that occurs within that time interval, and we plot the result as a function of $`E_0`$ in Fig. 2(a). The minimal $`P_{RL}`$ shows a few peaks but for only one value of $`E_0`$ ($`E_0=5.18`$ kV/cm) does it drop below 0.1. To take a closer look at the peak with strongest localization, we calculate the three probabilities $`P_{RL},P_{RR}`$, and $`P_{LL}`$ for that field, and plot them versus time in Fig. 3(a). In this case, $`P_{RL}`$ and $`P_{RR}`$ show an oscillatory behavior, while $`P_{LL}`$ remains negligible at all times. Physically, the two electrons oscillate between a state in which they are highly localized in the right well and another state which is completely delocalized.
In Fig. 3(b) we illustrate two operations of control of the wave function that can be performed with piecewise-constant electric fields. First, the thick-line curve of $`P_{RL}(t)`$ in Fig. 3(b) (produced by the electric field shown by a thick line in Fig. 3(c)) indicates that the low value of $`P_{RL}`$ obtained earlier (at $`t3\text{ps}`$ in Fig. 3(a)) can be maintained permanently by switching the field to another special value. The field value that locks $`P_{RL}`$ at its minimum was found through a systematic search. The second operation consists of un-locking the localization and resuming oscillations similar to those produced by the initial constant field. The timing of this second switching is found to be immaterial. The new field value is not arbitrary, but we find that many fields have a similar effect in the evolution of $`P_{RL}`$. In Figs. 3(b)-(c) we choose to go back to the initial field value, and switch the field at the time when $`P_{RL}`$ in 3(a) is maximum (6.1 ps), so that the subsequent oscillation is almost 180 degrees out of phase with the oscillation in 3(a).
We saw above that an appropriate constant electric field can localize to a large extent both electrons in a well of choice at certain times, and the localization can be locked by switching the field to another appropriate value. We will see next that the nature of the localization is different when oscillatory fields are applied. We consider a sinusoidal field of the form $`E(t)=E_0\mathrm{cos}\omega _0t`$ (we discuss below the effect of a slow switching-on.) Fig. 2(b) shows a contour plot of the minimum probability $`P_{RL}`$ achieved in a simulation interval of $`9\text{ps}`$, as a function of $`E_0`$ and $`\omega _0`$. Darker areas correspond to lower values of $`P_{RL}`$, i.e., to stronger localization. The most prominent feature in this plot is the existence of a number of “resonant” frequencies, which lead to localization for wide ranges of $`E_0`$. The lowest value of $`P_{RL}`$, 0.027, is obtained for $`E_0=5\text{kV/cm}`$ and $`\omega _0=5.6\text{meV}`$. Other combinations of $`E_0`$ and $`\omega _0`$ also yield values of $`P_{RL}`$ below 0.1. We computed $`P_{RL}(t),P_{RR}(t)`$ and $`P_{LL}(t)`$ for the frequencies that lead to strong localization and for different values of $`E_0`$, and in Fig. 4(a) we show the result for the case of strongest localization. Fig. 4(b) shows the corresponding field. We notice that the localization, or reduction of $`P_{RL}`$, results in an increase of both $`P_{RR}`$ and $`P_{LL}`$, as opposed to the case of a constant field, where only $`P_{RR}`$ increased at the expense of $`P_{RL}`$. This feature was observed for all the AC fields ($`\omega _0`$s and $`E_0`$s) we looked at. Physically, the localization with an AC field is such that both electrons are (with high probability) together in one of the wells, with roughly equal probabilities to be found in either well. The two electrons are here in quantum states that are both localized (to a large extent) and entangled. We mention that all the states occupied by the two interacting electrons in our simulations are entangled, in the usual sense that they are not factorizable into single-particle states. We use the label “localized/entangled” for the low $`P_{RL}`$ states of Fig. 4(a) to emphasize that, in these states, while each individual electron is delocalized (it can be found in either dot with roughly 50% probability), the two electrons are correlated and very likely to be found in the same dot.
A sudden switching-on of the field introduces high frequencies and therefore population of higher lying excited states, which results in a rugged evolution of $`P_{RL},P_{RR}`$ and $`P_{LL}`$. A smoother evolution of these probabilities (Fig. 4(c)) is produced by the field plotted in Fig. 4(d), which is switched on more slowly.
Once the two electrons are in the localized/entangled state obtained at $`t13.5\text{ps}`$ in Fig. 4(c), they can be forced to stay localized by suddenly shifting the field to an appropriate constant value. In Fig. 4(e) we show this effect, along with the electric field that produces it (shown in (f)). The control of entangled quantum states in solid state systems is of great current interest. Loss and Sukhorukov studied entanglement in coupled quantum dots involving the spin degree of freedom. Here we have identified a possibly complementary method that creates localization with entanglement in the spatial wave function of two electrons in coupled quantum dots.
In summary, we have found ways to create rapidly and to maintain localization in a two-electron coupled-quantum dots system with uniform electric fields with a simple time dependence. While a constant electric field creates pure localization of both electrons in one well, oscillatory fields induce entangled states that exhibit localization in either dot with roughly equal probabilities. This localization-with-entanglement is a purely many-body phenomenon brought about by the Coulomb interaction between the electrons, and the method we propose to create it is potentially useful in future applications of the physics of entangled states in solid state systems.
We thank Ataç Imamoglu for useful conversations. This research was supported in part by QUEST (NSF Grant No. DMR 91-20007). We made use of Crunch, the parallel computer at UCSB (NSF Grant No. CDA 96-01954).
Current address: Departamento de Física, FCEN, Universidad de Buenos Aires, (1428) Buenos Aires, Argentina. E-mail: pablot@df.uba.ar.
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# Rotations and 𝑒, 𝜈 Propagators, Part II
## 1 Introduction
It is necessarily true that spacetime quantities can be obtained from rotation considerations because rotations form a subgroup of the group of spacetime transformations. The overall process is to (i) define suitable rotation-founded quantities, i.e. 2-vectors and rotation matrices, etc, (ii) specify a basis for pairs of 2-vectors in Sec. 2, (iii) obtain rotation invariant projection operators for the basis in Sec. 3, and (iv) in Sec. 4 identify energy and momentum as functions of the parameters of the basis and find spacetime symmetric combinations of the projection operators.
Energy-momentum appears both as a factor with space and time in a phase angle and it appears as a factor with gamma matrices in a matrix composed of outer products of basis pairs. For the phase, a wide range of functions would suffice, but the outer products give specific functions of the basis parameters. Therefore, by changing the basis from Part I, we generate new outer products and obtain different energy-momentum functions.
The energy-momentum functions have null magnitude and the spacetime invariant projection operators are the Green’s functions for a relativistic spin 1/2 single particle wave equation. Thus the projection operators are propagators for a free neutrino or antineutrino.
## 2 A Basis with Mixed Eigenvectors
The choice of the $`uv`$ basis pairs in Part I started the process that lead to the electron propagator. So, to obtain a new propagator, we choose a different basis.
In each of the $`uv`$ basis pairs $`u_+^+,`$ $`u_{}^{},`$ $`v_+^+,`$ and $`v_{}^{},`$ see I(9) in Part I, the upper and lower 2-vectors are proportional to the same eigenvector, either $`u^+`$ or $`u^{},`$ where the eigenvectors are displayed in I(3). Consider now a new basis obtained by exchanging the spin up eigenvectors $`u^+`$ with the spin down eigenvectors $`u^{}`$ in the lower 2-vectors. We get
$$f_{}^+=\left(\begin{array}{c}e^{w/2}u^+\\ e^{w/2}u^{}\end{array}\right)f_+^{}=\left(\begin{array}{c}e^{w/2}u^{}\\ e^{w/2}u^+\end{array}\right),$$
$$g_{}^+=\left(\begin{array}{c}e^{w/2}u^+\\ e^{w/2}u^{}\end{array}\right)g_+^{}=\left(\begin{array}{c}e^{w/2}u^{}\\ e^{w/2}u^+\end{array}\right).$$
(1)
We rotate the 2-vectors so that the pairs acquire a common phase,
$$R_{}^+f_{}^+=e^{+i\theta /2}f_{}^+R_+^{}f_+^{}=e^{+i\theta /2}f_+^{}R_{}^+g_{}^+=e^{+i\theta /2}g_{}^+R_+^{}g_+^{}=e^{+i\theta /2}g_+^{},$$
(2)
where $`R_{}^+`$ means applying $`R`$ to the upper 2-vector and $`R^1`$ to the lower 2-vector. One can show that the rotated basis pairs (2) are normalized to $`2\mathrm{cosh}w`$ and they are mutually orthogonal,
$$i^{}i=2\mathrm{cosh}wi^{}j=0,i,j\{R_{}^+f_{}^+,R_+^{}f_+^{},R_{}^+g_{}^+,R_+^{}g_+^{}\},$$
(3)
where there is no sum over $`i`$ in the left expression and $`i`$ is not the same as $`j`$ in the middle expression.
## 3 Projection Operators for the Mixed Basis
The process of obtaining the projection operators for the $`fg`$ basis goes through just as with the $`uv`$ basis in Sec. I4. Any four component object $`\psi `$ can be expressed in the $`Rfg`$ basis (2),
$$\psi =e^{+i\theta /2}\psi _0=e^{+i\theta /2}(af_{}^++bf_+^{}+cg_{}^++dg_+^{}),$$
(4)
where $`a`$ = $`f_{}^{+}{}_{}{}^{}\psi _0/(2\mathrm{cosh}w)`$ and $`\psi _0`$ is $`\psi `$ for $`\theta `$ = 0. We need to obtain four projection operators, each operator selecting the part of $`\psi `$ that is proportional to one of the basis pairs. We split the phase into separate terms, thereby defining $`\mathrm{\Delta },`$
$$\frac{\theta }{2}=\mathrm{\Delta }+p^kx^k,$$
(5)
where the $`p^k`$ are three as-yet-unidentified functions of the ratio parameter $`w`$ and the unit vector $`n^k.`$ Together with the $`x^k`$ the $`p^k`$ constitute the parameters of a three dimensional delta function for selecting $`w`$ and $`n^k.`$ $`\mathrm{\Delta }`$ is the difference between the rotated basis phase $`\theta /2`$ and the delta function phase. For details see Sec. I4.
Following the same steps that lead to I(22), the projection operator $`K(2,1,R_{}^+f_{}^+)\gamma ^4,`$ for $`f_{}^+`$ is found to be
$$K(2,1,R_{}^+f_{}^+)\gamma ^4\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}(w)^{}}e^{i\theta _2^{}/2}f_{}^+f_{}^{+}{}_{}{}^{}e^{i\theta _1^{}/2}.$$
(6)
By straightforward calculation one can verify that
$$d^3x_1K(2,1,R_{}^+f_{}^+)\gamma ^4\psi (1)=e^{i\theta _2/2}(af_{}^+)=\psi (2)_{b=c=d=0},$$
(7)
where $`\psi (1)`$ is $`\psi `$ with $`\theta `$ = $`\theta _1`$ and $`\theta _i/2`$ = $`\mathrm{\Delta }_i^{}+p_{}^{k}{}_{}{}^{}x_i^k.`$ Comparing (7) with (4), we see that the projection operator recovers the desired portion of $`\psi `$ with a new phase $`\theta _2/2.`$ The projection operators for the other basis pairs differ from $`K(2,1,R_{}^+f_{}^+)\gamma ^4`$ in the matrices by substituting $`f_+^{},`$ $`g_{}^+,`$ and $`g_+^{}`$ in place of $`f_{}^+`$ in (6).
The appearance of the scalar product $`p^kx^k`$ and the expansion of $`f_{}^+f_{}^{+}{}_{}{}^{}`$ in (12) and Problem 1 below shows that $`K(2,1,R_{}^+f_{}^+)\gamma ^4`$ is a rotation invariant projection operator and the rotation angle $`\theta `$ is also a rotation invariant when $`\mathrm{\Delta }`$ and $`\gamma ^4`$ are rotation invariants. We turn now to making the definitions and finding the combinations that give spacetime invariance.
## 4 Space-Time Symmetry, the Neutrino Propagator
In this section it is found that the projection operators can have space-time symmetry, but only after combining two projection operators. There are three parts to this, just as in Part I. The three dimensional integral over $`x_1`$ in (7) becomes a surface integral in four dimensional space-time. The phase difference $`\mathrm{\Delta }`$ is shown to be the energy times the time, thereby making the rotation angle an invariant. The remaining factor is a non-invariant matrix for one projection operator, but two such matrices make a space-time invariant quantity for a particular choice of momentum function.
Time. Time can now be introduced by applying an integral expression for the phase factor $`e^{i\mathrm{\Delta }}`$, see I(25). By (6) and (7), we have
$$\psi (2)_{b=c=d=0}=d^3x_1K(2,1,R_{}^+f_{}^+)\gamma ^4\psi (1)$$
$$=d^3x_1\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}(w)^{}}e^{i\theta _2^{}/2}f_{}^+f_{}^{+}{}_{}{}^{}e^{i\theta _1^{}/2}e^{i\theta _1/2}\gamma ^4\gamma ^4\psi _0$$
$$=d^3x_1[\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}(w)^{}}e^{i\theta _2^{}/2}f_{}^+f_{}^{+}{}_{}{}^{}\gamma ^4]e^{i(\mathrm{\Delta }_1^{}\mathrm{\Delta }_1)}e^{i(p^kp^k)x_1^k}\gamma ^4\psi _0.$$
(8)
Let $`I`$ be the quantity in (8) in brackets. By (5) and I(24) and for $`\mathrm{\Delta }_2^{}>`$ 0, we get
$$I=\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}(w)^{}}e^{i\mathrm{\Delta }_2^{}}e^{ip^kx_2^k}f_{}^+f_{}^{+}{}_{}{}^{}\gamma ^4$$
$$=\frac{d^3p^{}}{(2\pi )^3}\frac{1}{2\mathrm{cosh}(w)^{}}(\frac{i}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑a\frac{e^{ia\mathrm{\Delta }_2^{}}}{a^21+iϵ})e^{ip^kx_2^k}f_{}^+f_{}^{+}{}_{}{}^{}\gamma ^4$$
$$=i\frac{d^4p^{}}{(2\pi )^4}e^{i(p^4t_2p^kx_2^k)}\frac{(\pm m)f_{}^+f_{}^{+}{}_{}{}^{}\gamma ^4}{p_{}^{4}{}_{}{}^{2}m^2\mathrm{cosh}^2w^{}+iϵ}$$
(9)
where we introduce a fourth component of momentum and a fourth component of displacement, i.e. energy and time,
$$p_{}^{4}{}_{}{}^{}\pm m\mathrm{cosh}(w)^{}aa\mathrm{\Delta }_2^{}=p_{}^{4}{}_{}{}^{}t_2,$$
(10)
where the sign is to be determined below and $`m`$ is the ‘mass,’ a term to be discussed below.
Surface Integral. Next, in (8) identify the integral over the 3-space $`x_1^k`$ with a surface integral in space-time,
$$\psi (2)_{b=c=d=0}=_Sd^4x_1Ie^{i(\mathrm{\Delta }_1^{}\mathrm{\Delta }_1)}e^{i(p^kp^k)x_1^k}N_\mu \gamma ^\mu \psi _0$$
(11)
where $`x^4`$ = $`t,`$ $`\mu `$ $`\{1,2,3,4\},`$ $`N^\mu `$ = $`\{N^k,N^4\}`$ is the unit normal to the three dimensional surface of integration $`S`$ in four dimensional space-time, and the space-time summation convention is used, $`N_\mu \gamma ^\mu `$ = $`N^4\gamma ^4N^k\gamma ^k.`$ In (7) and (8), the integration is over the three dimensional surface $`x_1^k`$ in the special space-time reference frame with the normal in the time direction, $`N^\mu `$ = $`\{0,0,0,1\}.`$
Matrices. The projection operator $`K(2,1,R_{}^+f_{}^+)\gamma ^4`$ rewritten with (8) and (9) fails to have space-time symmetry because of the quantities $`mf_{}^{+}{}_{}{}^{}f_{}^{+}{}_{}{}^{}{}_{}{}^{}\gamma ^4`$ and $`p_{}^{4}{}_{}{}^{2}(m\mathrm{cosh}w^{})^2`$ in (9). We can express the matrix $`mf_{}^{+}{}_{}{}^{}f_{}^{+}{}_{}{}^{}{}_{}{}^{}\gamma ^4`$ as a sum of sixteen linearly independent $`4\times 4`$ matrices such as the set of gamma matrices in Part I, Appendix A. By (1) and I(49), we get
$$mf_{}^+f_{}^{+}{}_{}{}^{}\gamma ^4=\frac{1}{2}[m\mathrm{cosh}(w)\gamma ^4m\mathrm{cosh}(w)n^j\gamma ^j+i\mathrm{sinh}(w)n^k\gamma ^k\gamma ^5i\mathrm{sinh}(w)\gamma ^4\gamma ^5]+$$
$$+\left(\begin{array}{ccc}u^+u_{}^{}{}_{}{}^{}& & 0\\ 0& & u^{}u_{}^{+}{}_{}{}^{}\end{array}\right),$$
(12)
where we drop the primes and the last term on the right is messy when written in terms of the $`\gamma ^A`$s. It can be shown that the matrix is not a space-time invariant when the four gammas $`\gamma ^\mu `$ in (12) transform as a space-time 4-vector.
Similarly one finds the expansion of $`mg_{}^{+}{}_{}{}^{}g_{}^{+}{}_{}{}^{}{}_{}{}^{}\gamma ^4`$ over the same set of gammas. We have
$$mg_{}^+g_{}^{+}{}_{}{}^{}\gamma ^4=\frac{1}{2}[m\mathrm{cosh}(w)\gamma ^4m\mathrm{cosh}(w)n^j\gamma ^ji\mathrm{sinh}(w)n^k\gamma ^k\gamma ^5+i\mathrm{sinh}(w)\gamma ^4\gamma ^5]+$$
$$\left(\begin{array}{ccc}u^+u_{}^{}{}_{}{}^{}& & 0\\ 0& & u^{}u_{}^{+}{}_{}{}^{}\end{array}\right),$$
(13)
When we add the two expressions, we get
$$m(f_{}^+f_{}^{+}{}_{}{}^{}+g_{}^+g_{}^{+}{}_{}{}^{})\gamma ^4=m\mathrm{cosh}(w)\gamma ^4m\mathrm{cosh}(w)n^k\gamma ^k.$$
(14)
Let $`\gamma ^k`$ and $`\gamma ^4`$ form a 4-vector, i.e. transform as in I(28). If $`m\mathrm{cosh}(w)n^k`$ and $`m\mathrm{cosh}w`$ also transform as the components of a 4-vector, then the expression (14) is a space-time invariant.
Energy-momentum. Energy and momentum have a special role because the energy-momentum 4-vector occurs (i) with $`x^\mu `$ in the invariant delta function phase $`p_\mu x^\mu `$ and also (ii) with the gamma matrices in the invariant polarization matrices $`p_\mu \gamma ^\mu .`$ Hence, as in Part I, we define the momentum and energy functions by comparing components. We get
$$p^k=m\mathrm{cosh}(w)n^kp^4=+m\mathrm{cosh}(w)a.$$
(15)
Then the pole (its at $`a`$ = 1) in (9) occurs when $`p^4`$ is the energy $`E,`$
$$E+m\mathrm{cosh}w\mathrm{\Delta }=Et.$$
(16)
Unlike the $`m`$ in Part I for an electron where $`m`$ is the magnitude of the electron energy-momentum space-time four vector, the quantity $`m`$ may not be made an invariant as easily. The difficulty is that the smallest $`\mathrm{cosh}w`$ can be is 1 while the energy $`E`$ can be made as small as desired by a suitable red-shifting boost. See Problem 3.
$`E>0`$ Propagator. Combining (8), (9), (14), and (15) gives an expression that is a space-time invariant. We have
$$\psi (2)_{b=d=0}=i\frac{d^4x_1d^4p^{}}{(2\pi )^4}e^{ip_\mu ^{}x_2^\mu }\frac{p_{\nu }^{}{}_{}{}^{}\gamma ^\nu }{p_\tau ^{}p^\tau iϵ}e^{ip_\eta ^{}x_1^\eta }N_\sigma \gamma ^\sigma \psi (1).$$
(17)
$`E>0`$ Propagator. The projection operators for the remaining two $`fg`$ basis pairs gives
$$m(f_{+}^{}{}_{}{}^{}f_{+}^{}{}_{}{}^{}{}_{}{}^{}+g_{+}^{}{}_{}{}^{}g_{+}^{}{}_{}{}^{}{}_{}{}^{})\gamma ^4=m\mathrm{cosh}(w)^{}\gamma ^4m\mathrm{cosh}(w)^{}n^k\gamma ^k.$$
(18)
The new matrix gives the same momentum function but negative the energy function in (15). Hence we now define the energy-momentum functions to be
$$p^km\mathrm{cosh}(w)n^kp^4m\mathrm{cosh}(w)a.$$
(19)
Then the pole (its at $`a`$ = 1) in (17) occurs when $`p^4`$ is the negative energy $`E,`$
$$Em\mathrm{cosh}w\mathrm{\Delta }=Et.$$
(20)
Note that (15) and (19) differ by the sign of $`m\mathrm{cosh}w`$ for $`p^4.`$ This is important because the wave function $`\psi `$ which is originally given, by (1), (2), and (4), as a function of the rotation angle $`\theta ,`$ ratio parameter $`w,`$ and unit 3-vector $`n^k,`$ is rewritten in terms of $`p^\mu `$ so that the delta functions can be applied. That step is necessary because the delta functions are in terms of $`p^\mu ,`$ see Sec. I4.
The sum of the projection operators for $`f_+^{}`$ and $`g_+^{}`$ can be written as
$$\psi (2)_{a=c=0}=i\frac{d^4x_1d^4p^{}}{(2\pi )^4}e^{ip_\mu ^{}x_2^\mu }\frac{p_{\nu }^{}{}_{}{}^{}\gamma ^\nu }{p_\tau ^{}p^\tau iϵ}e^{ip_\eta ^{}x_1^\eta }N_\sigma \gamma ^\sigma \psi (1).$$
(21)
Experimental confirmation. The fractions in (17) and (21) are sometimes written as inverse matrices, $`(p_{\mu }^{}{}_{}{}^{}\gamma ^\mu )^1`$ because
$$(p_{\mu }^{}{}_{}{}^{}\gamma ^\mu )(p_{\nu }^{}{}_{}{}^{}\gamma ^\nu )=p_{}^{4}{}_{}{}^{2}p_{}^{k}{}_{}{}^{2}=p_\tau ^{}p^\tau $$
(22)
On this basis the propagator for positive energy (17) and the propagator for negative energy (21) can be shown to be the Green’s functions for the matrix operator $`(p_{\mu }^{}{}_{}{}^{}\gamma ^\mu ).`$ In the representation of the gammas here, the upper and lower 2-vectors become separated in the wave equations obtained for the Green’s functions. Such wave functions conventionally describe free neutrinos. (The problem of whether or not all four components of a neutrino take part in interactions is immaterial for the ‘free’ neutrino.) We conclude that the propagators (17) and (21) can describe the free neutrino and antineutrino.
Time, energy, and rotation angle. By (5), (16) and (20) and for both positive and negative energy, we have
$$t=\frac{1}{2E}(2p^kx^k\theta )$$
(23)
We make $`t_1`$ = 0 on the surface of integration $`x_1^k.`$ Next we arrange that $`t_2`$ = 0 on the surface $`x_2^k`$ = $`x_1^k.`$ By (23), at any point $`x^k,`$
$$\mathrm{positive}\mathrm{energy}E>0\mathrm{and}\theta \mathrm{decreasing}t_2>0\mathrm{and}t_2\mathrm{increasing}$$
$$\mathrm{negative}\mathrm{energy}E<0\mathrm{and}\theta \mathrm{decreasing}t_2<0\mathrm{and}t_2\mathrm{decreasing}.$$
(24)
The primes have been dropped. As we found in Part I for the electron, we find here for the neutrino that the rotation angle $`\theta `$ decreases as the particle moves in time away from the surface $`x^k`$ at $`t`$ = 0.
## Appendix A Problems
1. (a) Derive (12), (13), and (14). (b) Express the matrix displayed in (12) and (13) linearly in terms of the gammas, I(47). (c) Show that the matrices (12) and (13) are rotation invariants.
2. For a massive particle, the proper time $`\tau `$ along a path $`x^\mu `$ that is proportional to $`p^\mu `$ is the ordinary time $`t`$ in a space-time reference frame with $`p^k`$ = 0. How is $`\tau `$ related to the rotation angle $`\theta \mathrm{?}`$ (b) For the neutrino, find straight line paths $`x^\mu `$ in space-time with $`x^k`$ proportional to $`p^k`$ and that satisfy (23). How does the proper time on these paths depend on the internal spin space rotation angle $`\theta \mathrm{?}`$ (See )
3. Suppose the parameter $`m`$ in the expressions (16) and (20) for energy $`E`$ is a constant for neutrinos independent of reference frame. Consider a source of neutrinos each with energy $`E>m>0`$ and momentum $`p^k`$ = $`\{E,0,0\},`$ both measured in the rest frame of the source. (a) Find the velocity of the frame with minimum energy $`E_0`$ = $`m`$ and $`p^k`$ = $`\{m,0,0\}.`$ (b) Show that boosts from the frame in (a) along $`+x`$ are not allowed. (c) Free neutrinos would move at the speed of light in the frame in (a) (at least on average over long distances). Would neutrino scattering be affected? If so, how?
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# 1 Introduction
## 1 Introduction
It is well known that $`B`$ physics is important to determine the elements of the Cabibbo-Kobayashi-Maskawa (CKM) matrix and physics beyond the standard model. Recently, the interest has been focused on the rare $`B`$ meson decays induced by the flavor changing neutral current (FCNC) due to the CLEO measurement of the radiative $`bs\gamma `$ decay . In the standard model, these rare decays occur at loop level and provide us information on the parameters of the CKM matrix elements as well as various hadronic form factors, such as the $`B`$ meson decay constant $`f_B`$.
As in the decays of $`B^+l^+\nu _l`$ , the helicity suppression effect is also expected in the flavor changing neutral current processes of $`B_{s,d}l^+l^{}`$. These decays are sensitive probes of top quark couplings such as the CKM elements $`V_{ts(d)}`$. The decay widths of these leptonic decay modes are given by:
$`\mathrm{\Gamma }(B_ql^+l^{})={\displaystyle \frac{\alpha ^2G_F^2f_{B_q}^2}{16\pi ^3}}m_{B_q}^3\left({\displaystyle \frac{m_l^2}{m_{B_q}^2}}\right)|V_{tb}V_{tq}^{}|^2C_{10}^2,`$ (1)
where $`G_F`$ is the fermi constant, $`M_{B_q}`$ and $`m_l`$ are $`B_q`$ meson and lepton masses and $`C_{10}`$ is the Wilson coefficient. Form Eq.(1), one has that $`B(B_se^+e^{},\mu ^+\mu ^{},\tau ^+\tau ^{})(6\times 10^8,2.6\times 10^1,1.0)\times 10^6`$ and $`B(B_de^+e^{},\mu ^+\mu ^{},\tau ^+\tau ^{})(4.2\times 10^7,1.8\times 10^2,4.5)\times 10^8`$ by taking $`|V_{tb}V_{ts}^{}|=0.04`$, $`|V_{tb}V_{td}^{}|=0.01`$, $`f_{B_{s,d}}200MeV`$, $`\tau _{B_s}1.61ps`$ and $`\tau _{B_d}1.5ps`$ . It is clear that the rates for the light lepton modes are too small to be measured due to the helicity suppressions, while that for the $`\tau `$ channel, although there is no suppression, it is hard to be observed experimentally because of the low efficiency.
It has been pointed out that the radiative leptonic $`B`$ decays of $`B^+ł^+\nu _l\gamma `$ ($`ł=e,\mu `$) have larger decay rates than that purely leptonic ones . Similar enhancements were also found for the radiative decays of $`B_{s,d}ł^+ł^{}\gamma `$ in the constituent and light-cone sum rule quark models . In fact, the decay amplitudes of $`B_{s,d}ł^+ł^{}\gamma `$ can be divided into the “internal-bremsstrahlung” (IB) parts, where the photon emits from the external charged leptons, which are still helicity suppressed for the light charged lepton modes, and the “structure-dependent” (SD) ones, in which one of the photon comes from intermediate states of $`B_ql^+l^{}`$, which are free of the helicity suppression. Therefore, the decay rates of $`B_ql^+l^{}\gamma `$ ( $`l=e,\mu `$ ) might have an enhancement with respect to the purely leptonic modes of $`B_ql^+l^{}`$ if the SD contributions to the decays are dominant.
In this paper, we will use the light front quark model to evaluate the hadronic matrix elements and study the decay rates for $`B_{s,d}l^+l^{}\gamma `$. It is known that as the recoil momentum increases, we have to start considering relativistic effects seriously. The light front quark model is the widely accepted relativistic quark model in which a consistent and relativistic treatment of quark spins and the center-of-mass motion can be carried out. In this work, we calculate for the first time the tensor form factors in $`P\gamma `$ ($`P`$ is a pseudoscalar meson) directly at time-like momentum transfers by using the relativistic light-front hadronic wave function. Within the framework of light-front quark model, one can calculate in the frame where the momentum transfer is purely longitudinal, $`i.e`$, $`p_{}=0`$ and $`p^2=p^+p^{}`$, which covers the entire range. Thus, we give their dependence on the momentum transfer $`p^2`$ in whole kinematic region of $`0p^2p_{\mathrm{max}}^2`$.
The paper is organized as follows. In Sec. 2, we present the decay amplitudes of $`B_{s,d}l^+l^{}\gamma `$. We study the form factors in the $`B\gamma `$ transition within the light front framework in Sec. 3. In Sec. 4, we calculate the decay branching ratios. We also compare our results with those in literature . We give our conclusions in Sec. 5.
## 2 Decay Amplitudes
The main contributions for the processes of $`B_ql^+l^{}\gamma `$ ($`l=e,\mu `$, $`\tau `$) arise from the effective Hamiltonian that induces the purely leptonic modes of $`B_ql^+l^{}`$. The QCD-corrected amplitude for $`bl^+l^{}q`$ in the SM is given by :
$``$ $`=`$ $`{\displaystyle \frac{G_F\alpha }{\sqrt{2}\pi }}V_{tb}V_{tq}^{}\{C_9^{eff}(\overline{q}\gamma _\mu P_Lb)\overline{l}\gamma ^\mu l+C_{10}\overline{q}\gamma _\mu P_Lb\overline{l}\gamma ^\mu \gamma _5l`$ (2)
$`2{\displaystyle \frac{C_7}{p^2}}\overline{q}i\sigma _{\mu \nu }p^\nu (m_bP_R+m_qP_L)b\overline{l}\gamma ^\mu l\},`$
where $`q=d`$ or $`s`$, $`P_{L,R}`$ = $`(1\gamma _5)/2`$, and $`p`$ is the momentum transfer, which is equal to the momentum of the lepton pair. The Wilson coefficients of $`C_7,C_9^{eff}`$ and $`C_{10}`$ can be found, for example, in , respectively.
The amplitude for $`B_ql^+l^{}\gamma `$ can be written as:
$`(B_ql^+l^{}\gamma )=_{IB}+_{SD}`$ (3)
where $`_{IB}`$ and $`_{SD}`$ represent the IB and SD contributions, respectively. For the IB part, the amplitude is clearly proportional to the lepton mass $`m_l`$ and it is found to be
$`_{IB}`$ $`=`$ $`ie{\displaystyle \frac{G_F\alpha }{\sqrt{2}\pi }}V_{tb}V_{tq}^{}f_{B_q}C_{10}m_ł\left[\overline{ł}({\displaystyle \frac{\overline{)}ϵ\overline{)}P_B}{2p_1q_\gamma }}{\displaystyle \frac{\overline{)}P_B\overline{)}ϵ}{2p_2q_\gamma }})\gamma _5ł\right],`$ (4)
where $`P_B`$, $`p_1`$, $`p_2`$ and $`q_\gamma `$, are the momentum of $`B_q`$, $`ł^+`$, $`ł^{}`$ and $`\gamma `$, respectively. In Eq.(4), the $`f_{B_q}`$ is the $`B_q`$ meson decay constant, defined by:
$`<0|\overline{s}\gamma ^\mu \gamma _5b|B_q(p)>`$ $`=`$ $`if_{B_q}p^\mu .`$ (5)
When a photon emitted from the charged internal line, the amplitude is suppressed by a factor $`m_b^2/m_W^2`$ in the Wilson coefficients and the main contribution comes from the photon radiating from the initial quark line. Therefore $`_{SD}`$ can be written as
$`_{SD}`$ $`=`$ $`{\displaystyle \frac{G_F\alpha }{\sqrt{2}\pi }}V_{tb}V_{tq}^{}\{C_9^{eff}\gamma (q_\gamma )|\overline{q}\gamma _\mu P_Lb|B(p+q_\gamma )\overline{l}\gamma ^\mu l`$ (6)
$`+C_{10}\gamma (q_\gamma )|\overline{q}\gamma _\mu P_Lb|B(p+q_\gamma )\overline{l}\gamma ^\mu \gamma _5l`$
$`2{\displaystyle \frac{C_7m_b}{p^2}}\gamma (q_\gamma )|\overline{q}i\sigma _{\mu \nu }p^\nu P_Rb|B(p+q_\gamma )\overline{l}\gamma ^\mu l\}.`$
From the amplitude in Eq.(6), we see that to find the decay rates, one has to evaluate the hadronic matrix elements. These elements can be parameterized as follows:
$`\gamma (q_\gamma )|\overline{q}\gamma _\mu \gamma _5b|B_q(p+q_\gamma )`$ $`=`$ $`e{\displaystyle \frac{F_{VA}}{M_{B_q}}}\left[(pq_\gamma )ϵ_\mu ^{}(ϵ^{}p)q_{\gamma \mu }\right],`$
$`\gamma (q_\gamma )|\overline{q}\gamma _\mu b|B_q(p+q_\gamma )`$ $`=`$ $`ie{\displaystyle \frac{F_{VV}}{M_{B_q}}}\epsilon _{\mu \alpha \beta \nu }ϵ^\alpha p^\beta q_\gamma ^\nu ,`$ (7)
$`\gamma (q_\gamma )|\overline{q}i\sigma _{\mu \nu }p^\nu \gamma _5b|B_q(p+q_\gamma )`$ $`=`$ $`ieF_{TA}[(pq_\gamma )ϵ_\mu ^{}(ϵ^{}p)q_{\gamma }^{}{}_{\mu }{}^{})],`$
$`\gamma (q_\gamma )|\overline{q}i\sigma _{\mu \nu }p^\nu b|B_q(p+q_\gamma )`$ $`=`$ $`eF_{TV}\epsilon _{\mu \nu \alpha \beta }ϵ^\nu q_\gamma ^\alpha p^\beta ,`$ (8)
where the $`ϵ_\mu `$ is the photon polarization vector, $`q_\gamma `$ and $`p+q_\gamma `$ are the four momenta of the photon and the $`B_q`$ meson, and $`F_{VA}`$, $`F_{VV}`$, $`F_{TA}`$ and $`F_{TV}`$ are the form factors of axial-vector, vector, axial-tensor and tensor, respectively. Form Eqs.(7) and (8), we rewrite the amplitude of Eq.(6) as,
$`_{SD}`$ $`=`$ $`{\displaystyle \frac{\alpha G_F}{2\sqrt{2}\pi }}V_{tb}V_{ts}^{}\{ϵ_{\mu \nu \alpha \beta }ϵ_{}^{}{}_{}{}^{\nu }p^\alpha q_{\gamma }^{}{}_{}{}^{\beta }[A\overline{ł}\gamma ^\mu ł+C\overline{ł}\gamma ^\mu \gamma _5ł]`$ (9)
$`+i[(pq_\gamma )ϵ_\mu ^{}(ϵ^{}p)q_{\gamma }^{}{}_{\mu }{}^{}][B\overline{ł}\gamma ^\mu ł+D\overline{ł}\gamma ^\mu \gamma _5ł]\},`$
where the factors of $`A`$-$`D`$ are defined by
$`A`$ $`=`$ $`{\displaystyle \frac{C_9^{eff}}{M_B}}F_{VA}(p^2)2C_7{\displaystyle \frac{m_b}{p^2}}F_{TA}(p^2),`$
$`B`$ $`=`$ $`{\displaystyle \frac{C_9^{eff}}{M_B}}F_{VV}(p^2)2C_7{\displaystyle \frac{m_b}{p^2}}F_{TV}(p^2),`$
$`C`$ $`=`$ $`{\displaystyle \frac{C_{10}}{M_B}}F_{VA}(p^2),`$
$`D`$ $`=`$ $`{\displaystyle \frac{C_{10}}{M_B}}F_{VV}(p^2),`$ (10)
respectively. The form factors of $`F_{VA}`$ and $`F_{VV}`$ have been evaluated previously in the light-front model , while that of $`F_{TA}`$ and $`F_{TV}`$ shall be studied in the next section.
## 3 Form Factors in the Light Front Model
In this section, we will use the light-front approach to calculate the tensor type form factors in Eq. (8) for $`B_q\gamma `$ ($`q=s`$ or $`d`$) transition. In this approach, the $`B`$ meson bound state consists of an anti-quark $`\overline{b}`$ and a quark $`q`$ with the total momentum of $`(p+q_\gamma )`$ and it is given by
$`|B(p+q_\gamma )>`$ $`=`$ $`{\displaystyle \underset{\lambda _1\lambda _2}{}}{\displaystyle [dk_1][dk_2]2(2\pi )^3\delta ^3(p+q_\gamma k_1k_2)}`$ (11)
$`\times \mathrm{\Phi }_B^{\lambda _1\lambda _2}(x,k_{})b_{\overline{b}}^+(k_1,\lambda _1)d_q^+(k_2,\lambda _2)|0>,`$
where $`k_{1(2)}`$ is the on-mass shell light front momentum of the internal quark $`\overline{b}(q)`$, the light front relative momentum variables $`(x,k_{})`$ are defined by
$`k_1^+=x(p+q_\gamma )^+,k_1=x(p+q_\gamma )_{}+k_{},`$ (12)
and
$`\mathrm{\Phi }_B^{\lambda _1\lambda _2}(x,k_{})=\left({\displaystyle \frac{2k_1^+k_2^+}{M_0^2\left(m_qm_b\right)^2}}\right)^{\frac{1}{2}}\overline{u}(k_1,\lambda _1)\gamma ^5v(k_2,\lambda _2)\varphi (x,k_{}),`$ (13)
with $`\varphi (x,k_{})`$ being the momentum distribution amplitude. The amplitude of $`\varphi (x,k_{})`$ can be solved in principle by the light-front QCD bound state equation. Here, we use the Gaussian type wave function:
$`\varphi (x,k_{})=N\sqrt{{\displaystyle \frac{dk_z}{dx}}}\mathrm{exp}\left({\displaystyle \frac{\stackrel{}{k}^2}{2\omega _B^2}}\right),`$ (14)
where
$`[dk_1]={\displaystyle \frac{dk^+dk_{}}{2(2\pi )^3}},N=4\left({\displaystyle \frac{\pi }{\omega _B^2}}\right)^{\frac{3}{4}},`$
$`k_z=\left(x{\displaystyle \frac{1}{2}}\right)M_0+{\displaystyle \frac{m_b^2m_q^2}{2M_0}},M_0^2={\displaystyle \frac{k_{}^2+m_q^2}{x}}+{\displaystyle \frac{k_{}^2+m_b^2}{1x}},`$
$`{\displaystyle \underset{\lambda }{}}u(k,\lambda )\overline{u}(k,\lambda )={\displaystyle \frac{m+\overline{)}k}{k^+}},{\displaystyle \underset{\lambda }{}}v(k,\lambda )\overline{v}(k,\lambda )={\displaystyle \frac{m\overline{)}k}{k^+}}.`$ (15)
We note that the wave function in Eq.(14) could be also applied to many other of hadronic transitions. For the gauged photon state, one has
$`|\gamma (q_\gamma )>`$ $`=`$ $`N^{}\{a^+(q_\gamma ,\lambda )+{\displaystyle \underset{\lambda _1\lambda _2}{}}{\displaystyle }[dk_1][dk_2]2(2\pi )^3\delta ^3(q_\gamma k_1k_2)`$ (16)
$`\times \mathrm{\Phi }_{q\overline{q}}^{\lambda _1\lambda _2\lambda }(q_\gamma ,k_1,k_2)b_q^+(k_1,\lambda _1)d_{\overline{q}}^+(k_1,\lambda _2)\}|0>,`$
where
$`\mathrm{\Phi }_{q\overline{q}}^{\lambda _3\lambda _4\lambda }(q_\gamma ,k_1,k_2)`$ $`=`$ $`{\displaystyle \frac{e_q}{ED}}\chi _{\lambda _2}^+\{2{\displaystyle \frac{q_{\gamma }^{}{}_{}{}^{}ϵ_{}}{q_\gamma ^+}}\stackrel{~}{\sigma }_{}ϵ_{}{\displaystyle \frac{\stackrel{~}{\sigma }_{}k_2_{}+im_2}{k_2^+}}`$ (17)
$`{\displaystyle \frac{\stackrel{~}{\sigma }_{}k_1_{}+im_1}{k_1^+}}\stackrel{~}{\sigma }_{}ϵ_{}\}\chi _{\lambda _1},`$
with
$`ED`$ $`=`$ $`{\displaystyle \frac{q_{\gamma }^{}{}_{}{}^{2}}{q_\gamma ^+}}{\displaystyle \frac{k_1_{}^2+m_1^2}{k_1^+}}{\displaystyle \frac{k_2_{}^2+m_2^2}{k_2^+}}.`$ (18)
In Eq. (17), we have used the two-component form of the light-front quark fields . Since the physically accessible kinematic region is $`0p^2p_{\mathrm{max}}^2`$ where $`p_{max}^2=M_B^2`$, to calculate the matrix elements in Eq. (8), we choose a frame where the transverse momentum $`p_{}`$ = $`0`$. Then $`p^2=p^+p^{}0`$ which can cover the entire range of the momentum transfers. By considering the “good” component of $`\mu =+`$, the tensor current in Eq. (8) can be rewritten as:
$`<\gamma (q_\gamma )|(q_+^+\gamma ^0\gamma _5b_{}q_{}^+\gamma ^0\gamma _5b_+)|B(p+q_\gamma )>`$ $`=`$ $`eF_{TA}\left(ϵ_{}^{}q_{\gamma }^{}{}_{}{}^{}\right),`$
$`<\gamma (q_\gamma )|(q_+^+\gamma ^0b_{}q_{}^+\gamma ^0b_+)|B(p+q_\gamma )>`$ $`=`$ $`ieF_{TV}ϵ^{ij}ϵ_i^{}q_{\gamma }^{}{}_{j}{}^{},`$ (19)
where $`q_+(b_+)`$ and $`q_{}(b_{})`$ are the light-front up and down component of the quark fields. In the two-component form , they are expressed by
$`q_+`$ $`=`$ $`\left(\begin{array}{c}\chi \\ 0\end{array}\right)`$ (22)
and
$`q_{}`$ $`=`$ $`{\displaystyle \frac{1}{i^+}}(i\alpha _{}_{}+\beta m_q)q_+=\left(\begin{array}{c}0\\ \frac{1}{^+}(\stackrel{~}{\sigma }_{}_{}+m_q)\chi _q\end{array}\right).`$ (25)
In Eq. (21), $`\chi _q`$ is a two-component spinor field and $`\sigma `$ is the Pauli matrix. The form factors of $`F_{TA}`$ and $`F_{TV}`$ in Eqs. (19) are then found to be:
$`F_{TA}(p^2)`$ $`=`$ $`{\displaystyle \frac{dxd^2k_{}}{2(2\pi )^3}\mathrm{\Phi }(x^{},k_{}^2)}`$ (26)
$`\times \left\{{\displaystyle \frac{1}{3}}{\displaystyle \frac{A_1+A_2k_{}^2\mathrm{\Theta }}{m_b^2+k_{}^2}}+{\displaystyle \frac{1}{3}}{\displaystyle \frac{B_1+B_2k_{}^2\mathrm{\Theta }}{m_q^2+k_{}^2}}\right\}`$
and
$`F_{TV}(p^2)`$ $`=`$ $`{\displaystyle \frac{dxd^2k_{}}{2\left(2\pi \right)^3}\mathrm{\Phi }(x^{},k_{}^2)}`$ (27)
$`\times \left\{{\displaystyle \frac{1}{3}}{\displaystyle \frac{C_1+C_2k_{}^2\mathrm{\Theta }}{m_b^2+k_{}^2}}+{\displaystyle \frac{1}{3}}{\displaystyle \frac{D_1+D_2k_{}^2\mathrm{\Theta }}{m_q^2+k_{}^2}}\right\}`$
where
$`A_1`$ $`=`$ $`{\displaystyle \frac{2}{xx^2(1x^{})(1x)}}\{(x^{}+x2x^{}x)[x^{}(x1)x(2x1)]k_{}^2`$
$`+x[(xx^{})+2x^{}x(1x)]m_b^2+2x^2(1x^{})^2m_bm_q\},`$
$`A_2`$ $`=`$ $`{\displaystyle \frac{2(xx^{})}{xx^2(1x^{})(1x)}}\{(x^{}+x2x^{}x)(12x)k_{}^2`$
$`+2x(1x^{})m_qm_bx(12x)(1x^{})^2m_q^2+x^{}(1+x^{}x2xx^2)m_b^2\},`$
$`B_1`$ $`=`$ $`{\displaystyle \frac{2}{x^{}x(1x)(1x^{})^2}}\{(x^{}+x2x^{}x)(12x+2x^2x^{}x)k_{}^2`$
$`+2xx^{}(1x^{})(1x)m_qm_b+(1x^{})(x^{}+x4x^{}x+2x^{}x^2)m_q^2\},`$
$`B_2`$ $`=`$ $`{\displaystyle \frac{2(xx^{})}{x^{}x(1x)(1x^{})^2}}\{(x^{}+x2x^{}x)(12x)k_{}^2`$
$`[(12x)x^2(1x)m_b^2+(1x^{})(x^{}+x3x^{}x2x^2(1x^{}))m_q^2]\},`$
$`C_1`$ $`=`$ $`{\displaystyle \frac{2}{xx_{}^{}{}_{}{}^{2}(1x^{})(1x)}}\{(xx^{}+xx^{})(x^{}+x2x^{}x)k_{}^2`$
$`2x^2(1x^{})^2m_bm_q+x^{}[(xx^{})2(1x^{})x^2]m_b^2\},`$
$`C_2`$ $`=`$ $`{\displaystyle \frac{2(xx^{})}{xx_{}^{}{}_{}{}^{2}(1x^{})(1x)}}\{(x^{}+x2x^{}x)k_{}^2+x^{}(12x+xx^{})m_b^2`$
$`m_q^2x(1x^{})^22x(1x^{})m_qm_b\},`$
$`D_1`$ $`=`$ $`{\displaystyle \frac{2}{x(1x)x^{}(1x^{})^2}}\{(1x)(12x+x^{})(x^{}+x2x^{}x)k_{}^2`$
$`[(1x^{})(x^{}+x2x^22x^{}x+2x^2x^{})m_q^2+2x_{}^{}{}_{}{}^{2}(1x)^2m_qm_b]\},`$
$`D_2`$ $`=`$ $`{\displaystyle \frac{2(xx^{})}{x^{}(1x)x(1x^{})^2}}\{(x^{}+x2x^{}x)k_{}^2+x^2(1x)m_b^2`$
$`+(1x^{})(x^{}xx^{}x)m_q^2\},`$
$`\mathrm{\Phi }(x,k_{}^2)`$ $`=`$ $`N\left({\displaystyle \frac{2x(1x)}{M_0^2(m_qm_b)^2}}\right)^{1/2}\sqrt{{\displaystyle \frac{dk_z}{dx}}}\mathrm{exp}\left({\displaystyle \frac{\stackrel{}{k}^2}{2\omega _B^2}}\right),`$
$`\mathrm{\Theta }`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Phi }(x,k_{}^2)}}{\displaystyle \frac{d\mathrm{\Phi }(x,k_{}^2)}{dk_{}^2}},`$
$`x^{}`$ $`=`$ $`x\left(1{\displaystyle \frac{p^2}{M_B^2}}\right),\stackrel{}{k}=(\stackrel{}{k}_{},\stackrel{}{k}_z).`$ (28)
To illustrate numerical results, we input $`m_q=m_s=0.4,`$ $`m_b=4.5`$ in $`GeV`$, and $`\omega =0.55`$ $`GeV`$ which is determined by fitting $`f_{B_s}=200GeV`$ via Eq. (5). The values of $`F_{TA}`$ and $`F_{TV}`$ in the entire range of $`p^2`$ are shown in Fig. 2. We note that the tails of $`F_{TV,TA}`$ at the large momentum transfers go down may be the long distance contribution associated with $`BB^{}\gamma `$ vertex, which is not included in the present work. It is interesting to note that the formulas in Eqs. (26) and (27) can be used for other pseudoscalars, such as pions and kaons, to the photon transitions as well once we put in the corresponding masses.
## 4 Decay Branching Ratios
The partial decay width for $`B_ql^+l^{}\gamma `$ in the $`B`$ rest frame is given by
$`d\mathrm{\Gamma }={\displaystyle \frac{1}{2M_B}}||^2(2\pi )^4\delta ^4(P_Bp_1p_2q_\gamma ){\displaystyle \frac{d\stackrel{}{q_\gamma }}{(2\pi )^32E_\gamma }}{\displaystyle \frac{d\stackrel{}{p}_1}{(2\pi )^32E_1}}{\displaystyle \frac{d\stackrel{}{p}_2}{(2\pi )^32E_2}},`$ (29)
where the square of the matrix element can be written as
$`||^2=|_{IB}|^2+|_{SD}|^2+2Re(_{IB}_{SD}^{}).`$ (30)
To describe the kinematic of $`B_ql^+l^{}\gamma `$, we defined two variables of $`x_\gamma =2P_Bq_\gamma /M_B`$ and $`y=2P_Bp_1/M_B`$. One can easily write the transfer momentum $`p^2`$ in term of $`x_\gamma `$ as
$`p^2=M_B^2(1x_\gamma ).`$ (31)
The double differential decay rate is found to be
$`{\displaystyle \frac{d^2\mathrm{\Gamma }^l}{dx_\gamma d\lambda }}`$ $`=`$ $`{\displaystyle \frac{M_B}{256\pi ^3}}\left|M\right|^2=C\rho (x_\gamma ,\lambda ),`$ (32)
where
$`C=\alpha |{\displaystyle \frac{\alpha V_{tb}V_{ts}^{}}{8\pi ^2}}|^2G_F^2M_B^5,`$ (33)
and
$`\rho (x_\gamma ,\lambda )`$ $`=`$ $`\rho _{IB}(x_\gamma ,\lambda )+\rho _{SD}(x_\gamma ,\lambda )+\rho _{IN}(x_\gamma ,\lambda ),`$ (34)
with
$`\rho _{IB}`$ $`=`$ $`4|f_BC_{10}|^2{\displaystyle \frac{r_l}{M_B^2x_\gamma ^2}}\left\{{\displaystyle \frac{x_\gamma ^22x_\gamma +24r_l}{\lambda (1\lambda )}}2r_l({\displaystyle \frac{1}{\lambda ^2}}+{\displaystyle \frac{1}{(1\lambda )^2}})\right\},`$
$`\rho _{SD}`$ $`=`$ $`{\displaystyle \frac{M_B^2}{8}}x_\gamma ^2\{(|A|^2+|B|^2)[(1x_\gamma +2r_l)2(1x_\gamma )(\lambda \lambda ^2)]`$
$`+(|C|^2+|D|^2)\left[(1x_\gamma 2r_l)2(1x_\gamma )(\lambda \lambda ^2)\right]`$
$`+2Re(B^{}C+A^{}D)(1x_\gamma )(2\lambda 1)\},`$
$`\rho _{IN}`$ $`=`$ $`f_BC_{10}r_l\left\{Re(A){\displaystyle \frac{x_\gamma }{\lambda (1\lambda )}}+Re(D){\displaystyle \frac{x_\gamma (12\lambda )}{\lambda (1\lambda )}}\right\}.`$ (35)
Here $`\lambda =(x_\gamma +y1)/x_\gamma `$ and $`r_l=m_l^2/M_B^2`$ and the physical regions for $`x_\gamma `$ and $`\lambda `$ are given by:
$`0`$ $``$ $`x_\gamma 14r_l,`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{4r_l}{1x_\gamma }}}`$ $``$ $`\lambda {\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{4r_l}{1x_\gamma }}}.`$ (36)
For the Wilson coefficients $`C_7`$ and $`C_{10}`$, we use the results given by Refs. and they are
$`C_7=0.315,C_{10}=4.642.`$
The analytic expressions for $`C_9^{eff}`$ in the next-to-leading order approximation is given by:
$`C_9^{eff}`$ $`=`$ $`C_9+0.124w(s)+g(\widehat{m_c},s)(3C_1+C_2+3C_3+C_4+3C_5+C_6)`$ (37)
$`{\displaystyle \frac{1}{2}}g(\widehat{m_q},s)(C_3+3C_4){\displaystyle \frac{1}{2}}g(\widehat{m_b},s)(4C_3+4C_4+3C_5+C_6)`$
$`+{\displaystyle \frac{2}{9}}(3C_3+C_4+3C_5+C_6),`$
with
$`C_9`$ $`=`$ $`4.227,`$
$`(3C_1+C_2+3C_3+C_4+3C_5+C_6)`$ $`=`$ $`0.359,`$
$`(C_3+3C_4)`$ $`=`$ $`0.0659,`$
$`(4C_3+4C_4+3C_5+C_6)`$ $`=`$ $`0.0675,`$
$`(3C_3+C_4+3C_5+C_6)`$ $`=`$ $`0.00157,`$ (38)
and $`s=p^2/m_b^2`$. In the Eq.(37), the function of $`w(\widehat{s})`$ comes from the single-gluon correction to the matrix element of $`O_9`$ and its form can be found in Refs., while that of $`g(\widehat{m_i},s)`$ from the one loop contributions of the four quark operators $`O_1`$ \- $`O_6`$, given by :
$`g(\widehat{m_i},s)`$ $`=`$ $`{\displaystyle \frac{8}{9}}\mathrm{ln}\widehat{m_i}+{\displaystyle \frac{8}{27}}+{\displaystyle \frac{4}{9}}y_i{\displaystyle \frac{2}{9}}(2+y_i)\sqrt{|1y_i|}`$ (41)
$`\times \{\begin{array}{cc}(\mathrm{ln}\frac{1+\sqrt{1y_i}}{1\sqrt{1y_i}}i\pi )& \hfill \text{for }y_i<1\\ 2\mathrm{arctan}\frac{1}{\sqrt{y_i1}}& \hfill \text{for }y_i>1\end{array}`$
with $`y_i=\widehat{m_i}^2/s`$ and $`\widehat{m_i}=m_i/m_b`$. By taking into account the long distance effects mainly due to the $`J/\psi `$ family resonances, one may use the replacement
$`g(\widehat{m_c},s)g(\widehat{m_c},s){\displaystyle \frac{3\pi }{\alpha ^2}}{\displaystyle \underset{V=J/\psi ,\psi ^{}}{}}{\displaystyle \frac{\widehat{m_V}Br(Vl^+l^{})\widehat{\mathrm{\Gamma }}_{tot}^V}{s\widehat{m_V}^2+i\widehat{m_V}\widehat{\mathrm{\Gamma }}_{tot}^V}}`$ (42)
where $`\widehat{m_V}=m_V/m_b`$ and $`\widehat{\mathrm{\Gamma }}_{tot}=\mathrm{\Gamma }/m_b`$ . The masses and decay widths of the corresponding mesons in Eq.(35) are listed in table 1.
In Figs.2 and 3 we present the differential decay rates of $`B_s\mu ^+\mu ^{}\gamma `$ and $`B_s\tau ^+\tau ^{}\gamma `$ as functions of $`x_\gamma `$, with and without resonance ($`J/\psi and\psi ^{}`$) contributions. From these figures we see that the contributions from the IB parts, corresponding to small $`x_\gamma `$ region, are infrared divergence. To obtain the decay width of $`B_s\tau ^+\tau ^{}\gamma `$, a cut on the photon energy is needed. Our results for the integrated branching ratios without and with long-distance effects as well as those given by the constituent quark model and the light-cone QCD sum rule model are summarized in tables 2 and 3, respectively. Here, we have used the cut value of $`\delta =0.01`$ and $`m_c=1.5GeV`$, $`m_s=0.4GeV`$, $`|V_{tb}V_{ts}^{}|=0.045`$, $`|V_{tb}V_{td}^{}|=0.01`$, $`\tau (B_s)=1.61\times 10^{12}s`$ and $`\tau (B_d)=1.5\times 10^{12}s`$.
We now compare our results with those in the literature . As shown in Table 2, the decay branching ratios for $`B_sl^+l^{}\gamma `$ ($`ł=e,\mu `$) without long distance contributions found in our calculations are smaller than those in the constituent quark model , whereas for $`B_dl^+l^{}\gamma `$ the statement are much smaller than ones. It is mainly due to that in the constituent quark model the decay rate of $`B_ql^+l^{}\gamma `$ is proportional to the inverse of the quark mass $`m_q`$. It is interesting to see that our results are close to those in the light-cone QCD sum rule model . We note that the SD contributions for the decays in both constituent quark and light-cone QCD sum rule models are sensitive to the values of the decay constants $`f_{B_{s,d}}`$.
## 5 Conclusions
We have studied the decays of $`B_{s,d}l^+l^{}\gamma `$ within the light-front model. We have calculated the tensor type form factors and used these form factors to evaluate the decay branching ratios. We have found that, in the standard model, the branching ratios of $`B_{s(d)}e^+e^{}\gamma `$, $`B_{s(d)}\mu ^+\mu ^{}\gamma `$ and $`B_{s(d)}\tau ^+\tau ^{}\gamma `$ are $`7.1\times 10^9`$ ($`1.5\times 10^{10}`$), $`8.3\times 10^9`$ ($`1.8\times 10^{10}`$) and $`1.6\times 10^8`$ ($`6.2\times 10^{10}`$), respectively. Comparing with the purely leptonic decays of $`B_{s,d}l^+l^{}`$, we have shown that the branching ratios of $`B_{s,d}l^+l^{}\gamma `$ have the same order of magnitude for the $`\mu `$ channel but that of $`B_{s,d}e^+e^{}\gamma `$ are much larger. We conclude that some of the radiative leptonic decays of $`B_{s,d}ł^+ł^{}\gamma `$ could be measured in the $`B`$ factories as well as LHC-B experiments, where approximately, $`6\times 10^{11}(2\times 10^{11})B_d(B_s)`$ mesons are expected to be produced per year.
Acknowledgments
This work was supported in part by the National Science Council of the Republic of China under contract number NSC-89-2112-M-007-013 and NSC-89-2112-M-006-026 .
## Figure Captions
* The values of the form factors $`F_{TA}`$ (solid curve) and $`F_{TV}`$ (dashed curve) as functions of the momentum transfer $`p^2`$ for $`B_s\gamma `$.
* The differential decay branching ratio $`dB(B_s\mu ^+\mu ^{}\gamma )/dx_\gamma `$ as a function of $`x_\gamma =2E_\gamma /M_{B_s}`$ with(solid curve) and without (dashed curve) long distance.
* The differential decay branching ratio $`dB(B_s\tau ^+\tau ^{}\gamma )/dx_\gamma `$ as a function of $`x_\gamma =2E_\gamma /M_{B_s}`$ with (solid curve) and without (dashed curve) long distance.
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# July 2000 UM-P-030/2000 RCHEP-006/2000 Further studies on relic neutrino asymmetry generation II: a rigorous treatment of repopulation in the adiabatic limit
## I Introduction
The computationally convenient and rewarding static approximation may be rigorously derived from the quantum kinetic equations (QKEs) in the adiabatic Boltzmann limit , assuming that the repopulation or refilling function is insignificantly small. In actual numerical studies, however, we also adopt the view that a depleted neutrino momentum state is instantaneously refilled from the background plasma. The latter condition naturally implies that while the repopulation function may be small, it cannot be identically zero.
In principle, simultaneous suppositions of a vanishing refilling function and instantaneous repopulation are somewhat contradictory, or at least uncontrolled. However, their combined validity in the collision dominant epoch is strongly supported by numerical evidence . Indeed, instantaneous repopulation without the assumption of a negligible repopulation function can be consistently implemented as a sensible approximation to the QKEs .
Commencing with the full-fledged QKEs, we demonstrate rigorously that the two aforementioned approximations, conceived originally for computational convenience, are indeed appropriate in a certain limit. Our approach differs from the fully consistent instantaneous repopulation approximation introduced in Ref. , since we do not make any a priori assumptions regarding the repopulation rate, instantaneous or otherwise. Instead, the present work shows that a “brute force” treatment of the QKEs incorporating a finite repopulation rate does indeed lead to approximate evolution equations that reduce to those following from the dual approximations of a negligible refilling function and instantaneous repopulation in the temperature regime of interest.
## II Nomenclature
Consider a two-flavour system comprising an active species $`\nu _\alpha `$ (where $`\alpha =e`$, $`\mu `$ or $`\tau `$), and a sterile species $`\nu _s`$, whose properties are characterised by the density matrix
$$\rho (p)=\frac{1}{2}[P_0(p)+𝐏(p)\sigma ],$$
(1)
in which the variables $`P_0`$ and $`𝐏(p)=P_x(p)\widehat{x}+P_y(p)\widehat{y}+P_z(p)\widehat{z}`$ are functions of both time $`t`$ and momentum $`p`$, and $`\sigma =\sigma _x\widehat{x}+\sigma _y\widehat{y}+\sigma _z\widehat{z}`$ are the Pauli matrices. The diagonal entries of $`\rho `$ represent respectively the $`\nu _\alpha `$ and $`\nu _s`$ distribution functions in momentum space, that is,
$`N_\alpha (p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[P_0(p)+P_z(p)]N^{\text{eq}}(p,0),`$ (2)
$`N_s(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[P_0(p)P_z(p)]N^{\text{eq}}(p,0),`$ (3)
where the reference distribution $`N^{\text{eq}}(p,0)`$ is defined as the equilibrium Fermi–Dirac function,
$$N^{\text{eq}}(p,\mu )=\frac{1}{2\pi ^2}\frac{p^2}{1+\mathrm{exp}\left(\frac{p\mu }{T}\right)},$$
(4)
with chemical potential $`\mu =0`$ at temperature $`T`$.
The evolution of $`P_0(p)`$ and $`𝐏(p)`$ is governed by the quantum kinetic equations (QKEs)
$`{\displaystyle \frac{𝐏}{t}}`$ $`=`$ $`𝐕(p)\times 𝐏(p)D(p)[P_x(p)\widehat{x}+P_y(p)\widehat{y}]+{\displaystyle \frac{P_0}{t}}\widehat{z},`$ (5)
$`{\displaystyle \frac{P_0}{t}}`$ $`=`$ $`R_\alpha (p).`$ (6)
Here, the quantity $`𝐕(p)=\beta (p)\widehat{x}+\lambda (p)\widehat{z}`$ is the matter potential vector , with
$`\beta (p)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }m^2}{2p}}\mathrm{sin}2\theta _0,`$ (7)
$`\lambda (p)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }m^2}{2p}}[b(p)a(p)\mathrm{cos}2\theta _0],`$ (8)
in which $`\mathrm{\Delta }m^2`$ is the mass-squared difference between the neutrino states, $`\theta _0`$ is the vacuum mixing angle, and
$`a(p)`$ $`=`$ $`{\displaystyle \frac{4\zeta (3)\sqrt{2}G_FL^{(\alpha )}T^3p}{\pi ^2\mathrm{\Delta }m^2}},`$ (9)
$`b(p)`$ $`=`$ $`{\displaystyle \frac{4\zeta (3)\sqrt{2}G_FA_\alpha T^4p^2}{\pi ^2\mathrm{\Delta }m^2m_W^2}},`$ (10)
given that $`G_F`$ is the Fermi constant, $`m_W`$ the $`W`$-boson mass, $`\zeta `$ the Riemann zeta function and $`A_e17`$, $`A_{\mu ,\tau }4.9`$. The effective total lepton number (for the $`\alpha `$-neutrino species)
$$L^{(\alpha )}=L_\alpha +L_e+L_\mu +L_\tau +\eta 2L_\alpha +\stackrel{~}{L},$$
(11)
combines all asymmetries individually defined as $`L_\alpha =(n_{\nu _\alpha }n_{\overline{\nu }_\alpha })/n_\gamma `$, where the symbol $`n`$ denotes number density, and $`\eta `$ is a small term due to the cosmological baryon asymmetry.
The decoherence function $`D(p)`$ is related to the total collision rate for $`\nu _\alpha `$, $`\mathrm{\Gamma }(p)`$, via
$$D(p)=\frac{1}{2}\mathrm{\Gamma }(p)=\frac{1}{2}\frac{p}{p_0}y_\alpha G_F^2T^5,$$
(12)
where $`p_03.15T`$ is the average momentum for a relativistic Fermi–Dirac distribution with zero chemical potential, $`y_e4`$, $`y_{\mu ,\tau }2.9`$ and $`z_e0.1`$, $`z_{\mu ,\tau }0.04`$.
The repopulation or refilling function
$$R_\alpha (p)\mathrm{\Gamma }(p)\left\{K_\alpha \frac{1}{2}[P_0(p)+P_z(p)]\right\},$$
(13)
with
$$K_\alpha \frac{N^{\text{eq}}(p,\mu )}{N^{\text{eq}}(p,0)},$$
(14)
is determined from assuming thermal equilibrium for all species in the background plasma, and that the $`\nu _\alpha `$ distribution is also approximately thermal . Physically, Eq. (13) means that as active neutrinos of some momentum are converted to sterile neutrinos, a gap is created in the Fermi–Dirac distribution. The weak interaction processes repopulate this depleted state, driving the ensemble back towards equilibrium.
A separate but equivalent set of expressions parameterises the evolution of the $`\overline{\nu }_\alpha \overline{\nu }_s`$ system. These are distinguished from their ordinary counterparts with an overhead bar.
## III Standard static/adiabatic limit: a brief outline
Together with the requirement of $`\alpha +s`$ lepton number conservation, we may derive from the QKEs an exact evolution equation for the neutrino asymmetry $`L_\alpha `$, that is ,
$$\frac{dL_\alpha }{dt}=\frac{1}{2n_\gamma }\beta [P_y(p)\overline{P}_y(p)]N^{\text{eq}}(p,0)𝑑p,$$
(15)
where the quantities $`P_y`$ and $`\overline{P}_y`$ are obtained from numerically integrating the respective QKEs for the neutrino and antineutrino systems given by Eq. (5). The role of the static/adiabatic limit approximation , therefore, is to generate approximate expressions for $`P_y`$ that are dynamically decoupled from the other variables in order to minimise the computational effort.
The first step in the formal adiabatic procedure of Ref. consists of setting the repopulation function to zero, i.e.,
$$R_\alpha 0,$$
(16)
thereby reducing the four-component QKEs to a system of three coupled homogeneous differential equations. Further manipulation produces the approximate equality
$$P_y(t)\frac{\beta D}{D^2+\lambda ^2}P_z(t),$$
(17)
in the limit $`D,|\lambda ||\beta |`$, such that Eq. (15) becomes<sup>*</sup><sup>*</sup>*The denominator in Eq. (17), and therefore Eq. (18), should properly contain an additive term $`\beta ^2`$ to ensure the equations’ validity even when the condition $`D|\beta |`$ does not hold (see companion paper Ref. ). In this work, however, we shall assume that the said condition is always met, and omit the $`\beta ^2`$ term for simplicity.
$$\frac{dL_\alpha }{dt}\frac{1}{2n_\gamma }\beta ^2\left[\frac{\overline{D}(\overline{N}_\alpha \overline{N}_s)}{\overline{D}^2+\overline{\lambda }^2}\frac{D(N_\alpha N_s)}{D^2+\lambda ^2}\right]𝑑p.$$
(18)
Detailed discussions of the standard adiabatic procedure may be found in the companion paper Ref. .
In numerical studies, an analogous expression describing sterile neutrino production for each momentum state is employed for tracking the quantity $`N_s(p)`$ in Eq. (18), while the $`\nu _\alpha `$ distribution function is taken to be
$$N_\alpha (p)N^{\text{eq}}(p,\mu ).$$
(19)
This is the so-called instantaneous repopulation approximation, which assumes that a depleted momentum state is immediately refilled from the background medium so that thermal equilibrium is always maintained.
A naïve interpretation of Eqs. (13) and (19) suggests that instantaneous repopulation leads directly to an identically zero refilling rate, which seems to argue for consistency between the central assumptions contained in Eqs. (16) and (19). However, the instantaneous repopulation limit is also related to taking the collision rate to infinity. Thus the right hand side of Eq. (13) is an a priori undetermined finite and generally nonzero function , that is apparently at odds with the assumption of a vanishing repopulation rate.
## IV Extended adiabatic approximation
The exact QKEs in Eq. (5) are fundamentally a system of four coupled first order differential equations, of both the homogeneous and inhomogeneous varieties, that is best displayed in matrix form,
$$\frac{}{t}\left(\begin{array}{c}P_x\\ P_y\\ P_z\\ P_0\end{array}\right)=\left(\begin{array}{cccc}D& \lambda & 0& 0\\ \lambda & D& \beta & 0\\ 0& \beta & D& D\\ 0& 0& D& D\end{array}\right)\left(\begin{array}{c}P_x\\ P_y\\ P_z\\ P_0\end{array}\right)+\left(\begin{array}{c}0\\ 0\\ 2DK_\alpha \\ 2DK_\alpha \end{array}\right)𝒦^{\text{rp}}𝐏^{\text{rp}}+𝐀.$$
(20)
Dealing firstly with the homogeneous part of Eq. (20), we introduce an instantaneous diagonal basis onto which we map the vector $`𝐏^{\text{rp}}(𝐏,P_0)`$ from its original fixed basis via the transformation
$$\left(\begin{array}{c}Q_1^{\text{rp}}\\ Q_2^{\text{rp}}\\ Q_3^{\text{rp}}\\ Q_4^{\text{rp}}\end{array}\right)𝐐^{\text{rp}}=𝒰_{\text{rp}}𝐏^{\text{rp}},$$
(21)
where, by definition,
$$𝒦_d^{\text{rp}}\text{diag}(\mathrm{\Lambda }_1,\mathrm{\Lambda }_2,\mathrm{\Lambda }_3,\mathrm{\Lambda }_4)=𝒰_{\text{rp}}𝒦^{\text{rp}}𝒰_{\text{rp}}^1,$$
(22)
such that the matrix $`𝒦_d^{\text{rp}}`$ is diagonal and similar to $`𝒦^{\text{rp}}`$.
The eigenvalues $`\mathrm{\Lambda }_i`$ are solutions to the quartic characteristic equation
$$\mathrm{\Lambda }^4+4D\mathrm{\Lambda }^3+(5D^2+\lambda ^2+\beta ^2)\mathrm{\Lambda }^2+2D(D^2+\lambda ^2+\beta ^2)\mathrm{\Lambda }+\beta ^2D^2=0,$$
(23)
formally given by
$$\mathrm{\Lambda }_i=D\pm \frac{\sqrt{D^2\lambda ^2\beta ^2\pm \sqrt{D^4+2D^2\lambda ^2+\lambda ^4+2\beta ^2\lambda ^22\beta ^2D^2+\beta ^4}}}{\sqrt{2}}.$$
(24)
We are primarily interested in the limit $`D,|\lambda ||\beta |`$, in which case the eigenvalues are well approximated by the leading order terms in the small $`\beta `$ power series expansion of Eq. (24),
$`\mathrm{\Lambda }_1`$ $`=`$ $`\mathrm{\Lambda }_2^{}=D+i\lambda +{\displaystyle \frac{i}{2}}{\displaystyle \frac{\beta ^2\lambda }{D^2+\lambda ^2}}+𝒪(\beta ^4),`$ (25)
$`\mathrm{\Lambda }_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\beta ^2D}{D^2+\lambda ^2}}+𝒪(\beta ^4),`$ (26)
$`\mathrm{\Lambda }_4`$ $`=`$ $`2D+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\beta ^2D}{D^2+\lambda ^2}}+𝒪(\beta ^4).`$ (27)
The transformation matrix $`𝒰_{\text{rp}}^1`$ comprises the column eigenvectors $`\kappa _i^{\text{rp}}`$, which for $`\lambda 0`$ are exactly
$$\kappa _i^{\text{rp}}=\left(\begin{array}{c}1\\ \frac{D+\mathrm{\Lambda }_i}{\lambda }\\ \frac{\beta (D+\mathrm{\Lambda }_i)^2}{\lambda \mathrm{\Lambda }_i(2D+\mathrm{\Lambda }_i)}\\ \frac{\beta D(D+\mathrm{\Lambda }_i)}{\lambda \mathrm{\Lambda }_i(2D+\mathrm{\Lambda }_i)}\end{array}\right),$$
(28)
while its inverse $`𝒰_{\text{rp}}`$ consists of the row vectors
$`v_i^{\text{rp}}`$ $`=\mathrm{\Lambda }_i(2D+\mathrm{\Lambda }_i)({\displaystyle \frac{1}{D^2}}{\displaystyle \underset{ji}{}}{\displaystyle \frac{D+\mathrm{\Lambda }_j}{\mathrm{\Lambda }_i\mathrm{\Lambda }_j}},\lambda {\displaystyle \underset{ji}{}}{\displaystyle \frac{1}{\mathrm{\Lambda }_i\mathrm{\Lambda }_j}},`$ (31)
$`{\displaystyle \frac{\lambda }{\beta D^2}}{\displaystyle \underset{ji}{}}{\displaystyle \frac{1}{\mathrm{\Lambda }_i\mathrm{\Lambda }_j}}\left[4D^3+2D^2(\mathrm{\Lambda }_l+\mathrm{\Lambda }_m+\mathrm{\Lambda }_n)+D(\mathrm{\Lambda }_l\mathrm{\Lambda }_m+\mathrm{\Lambda }_m\mathrm{\Lambda }_n+\mathrm{\Lambda }_n\mathrm{\Lambda }_l)+\mathrm{\Lambda }_l\mathrm{\Lambda }_m\mathrm{\Lambda }_n\right],`$
$`{\displaystyle \frac{\lambda }{\beta D}}{\displaystyle \underset{ji}{}}{\displaystyle \frac{1}{\mathrm{\Lambda }_i\mathrm{\Lambda }_j}}[4D^2+2D(\mathrm{\Lambda }_l+\mathrm{\Lambda }_m+\mathrm{\Lambda }_n)+\mathrm{\Lambda }_l\mathrm{\Lambda }_m+\mathrm{\Lambda }_m\mathrm{\Lambda }_n+\mathrm{\Lambda }_n\mathrm{\Lambda }_l]),`$
where the indices $`l`$, $`m`$ and $`n`$ are the three integers not equal to $`i`$, e.g., for $`i=1`$, $`l`$, $`m`$ and $`n`$ are $`2`$, $`3`$ and $`4`$ respectively.
Note that Eqs. (28) and (31) do not apply if two or more eigenvalues are degenerate, in which case the matrix $`𝒦^{\text{rp}}`$ may genuinely have less than four distinct eigenvectors thereby rendering $`𝒰_{\text{rp}}^1`$ momentarily uninvertible, or, if four linearly independent eigenvectors exist, alternative methods are needed for their evaluation. However, such circumstances are hard to come by and have practically no influence on the outcome of this work. Thus the exact QKEs \[Eq. (20)\] are equivalently
$$\frac{𝐐^{\text{rp}}}{t}=𝒦_d^{\text{rp}}𝐐^{\text{rp}}+𝒰_{\text{rp}}𝐀𝒰_{\text{rp}}\frac{𝒰_{\text{rp}}}{t}^1𝐐^{\text{rp}},$$
(32)
written in the new instantaneous basis where the transformation matrices are valid.
The term $`𝒰_{\text{rp}}\frac{𝒰_{\text{rp}}}{t}^1`$ in Eq. (32) contains explicit dependence on the derivatives of the parameters $`D`$, $`\lambda `$ and $`\beta `$. These are discarded in the adiabatic limit (see Ref. for the relevant constraintsThe bounds in Ref. are calculated assuming $`\frac{P_0}{t}0`$, and should naturally be different from those one would obtain with the incorporation of a finite repopulation function. However, we expect the difference to be minimal given the similarity between the relevant entries in the two cases’ respective transformation matrices in the appropriate limit.) so that the remaining terms form a set of four decoupled inhomogeneous differential equations given by
$$\frac{𝐐^{\text{rp}}}{t}𝒦_d^{\text{rp}}𝐐^{\text{rp}}+𝐁,$$
(33)
where $`𝐁𝒰_{\text{rp}}𝐀`$, or equivalently in index form,
$$\frac{}{t}Q_i^{\text{rp}}(t)\mathrm{\Lambda }_i(t)Q_i^{\text{rp}}(t)+B_i(t),i=1,2,\mathrm{},4.$$
(34)
Equation (34) may be formally solved to give
$$Q_i^{\text{rp}}(t)=e^{_0^t\mathrm{\Lambda }_i(t^{})𝑑t^{}}Q_i^{\text{rp}}(0)+e^{^t\mathrm{\Lambda }_i(t^{})𝑑t^{}}_i,$$
(35)
with the inhomogeneous segment of the QKEs contained entirely in the integral
$$_i=_0^te^{^t^{}\mathrm{\Lambda }_i(t^{\prime \prime })𝑑t^{\prime \prime }}B_i(t^{})𝑑t^{},$$
(36)
thus amounting to the following solution for the vector $`𝐏^{\text{rp}}`$,
$`\left(\begin{array}{c}P_x(t)\\ P_y(t)\\ P_z(t)\\ P_0(t)\end{array}\right)`$ $`=`$ $`𝒰_{\text{rp}}^1(t)[\left(\begin{array}{cccc}e^{_0^t\mathrm{\Lambda }_1(t^{})𝑑t^{}}& 0& 0& 0\\ 0& e^{_0^t\mathrm{\Lambda }_2(t^{})𝑑t^{}}& 0& 0\\ 0& 0& e^{_0^t\mathrm{\Lambda }_3(t^{})𝑑t^{}}& 0\\ 0& 0& 0& e^{_0^t\mathrm{\Lambda }_4(t^{})𝑑t^{}}\end{array}\right)\left(\begin{array}{c}Q_1^{\text{rp}}(0)\\ Q_2^{\text{rp}}(0)\\ Q_3^{\text{rp}}(0)\\ Q_4^{\text{rp}}(0)\end{array}\right)`$ (58)
$`+\left(\begin{array}{cccc}e^{^t\mathrm{\Lambda }_1(t^{})𝑑t^{}}& 0& 0& 0\\ 0& e^{^t\mathrm{\Lambda }_2(t^{})𝑑t^{}}& 0& 0\\ 0& 0& e^{^t\mathrm{\Lambda }_3(t^{})𝑑t^{}}& 0\\ 0& 0& 0& e^{^t\mathrm{\Lambda }_4(t^{})𝑑t^{}}\end{array}\right)\left(\begin{array}{c}_1\\ _2\\ _3\\ _4\end{array}\right)],`$
in the adiabatic limit.
Observe in Eq. (27) that the real components of the eigenvalues $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$ and $`\mathrm{\Lambda }_4`$ are of the order $`D`$, while $`\mathrm{\Lambda }_3`$ is proportional to $`\beta ^2`$. This implies that the exponentials $`\mathrm{exp}[_0^t\mathrm{\Lambda }_i(t^{})𝑑t^{}]`$, where $`i=1,2,4`$, are rapidly damped relative to the “decay” time scale of their $`\mathrm{\Lambda }_3`$ counterpart in the homogeneous part of Eq. (58) (for a full discussion, see Ref. ). We may therefore implement in Eq. (58) the collision dominance approximation
$$e^{_0^t\mathrm{\Lambda }_i(t^{})𝑑t^{}}0,i=1,2,4,$$
(59)
to obtain
$`P_y(t)`$ $``$ $`{\displaystyle \frac{D+\mathrm{\Lambda }_3}{\lambda }}e^{_0^t\mathrm{\Lambda }_3(t^{})𝑑t^{}}Q_3^{\text{rp}}(0){\displaystyle \underset{i=1}{\overset{4}{}}}{\displaystyle \frac{D+\mathrm{\Lambda }_i}{\lambda }}e^{^t\mathrm{\Lambda }_i(t^{})𝑑t^{}}_i,`$ (60)
$`P_z(t)`$ $``$ $`{\displaystyle \frac{\beta (D+\mathrm{\Lambda }_3)^2}{\lambda \mathrm{\Lambda }_3(2D+\mathrm{\Lambda }_3)}}e^{_0^t\mathrm{\Lambda }_3(t^{})𝑑t^{}}Q_3^{\text{rp}}(0){\displaystyle \underset{i=1}{\overset{4}{}}}{\displaystyle \frac{\beta (D+\mathrm{\Lambda }_i)^2}{\lambda \mathrm{\Lambda }_i(2D+\mathrm{\Lambda }_i)}}e^{^t\mathrm{\Lambda }_i(t^{})𝑑t^{}}_i,`$ (61)
$`P_0(t)`$ $``$ $`{\displaystyle \frac{\beta D(D+\mathrm{\Lambda }_3)}{\lambda \mathrm{\Lambda }_3(2D+\mathrm{\Lambda }_3)}}e^{_0^t\mathrm{\Lambda }_3(t^{})𝑑t^{}}Q_3^{\text{rp}}(0)+{\displaystyle \underset{i=1}{\overset{4}{}}}{\displaystyle \frac{\beta D(D+\mathrm{\Lambda }_i)}{\lambda \mathrm{\Lambda }_i(2D+\mathrm{\Lambda }_i)}}e^{^t\mathrm{\Lambda }_i(t^{})𝑑t^{}}_i,`$ (62)
and consequently, from combining the above expressions for $`P_y`$ and $`P_z`$,
$`P_y(t)=`$ $`{\displaystyle \frac{\mathrm{\Lambda }_3(2D+\mathrm{\Lambda }_3)}{\beta (D+\mathrm{\Lambda }_3)}}P_z(t)`$ (64)
$`+{\displaystyle \underset{j=1,2,4}{}}\left({\displaystyle \frac{\mathrm{\Lambda }_3(2D+\mathrm{\Lambda }_3)}{\beta (D+\mathrm{\Lambda }_3)}}{\displaystyle \frac{\mathrm{\Lambda }_j(2D+\mathrm{\Lambda }_j)}{\beta (D+\mathrm{\Lambda }_j)}}\right){\displaystyle \frac{\beta (D+\mathrm{\Lambda }_j)^2}{\lambda \mathrm{\Lambda }_j(\mathrm{\Lambda }_j+2D)}}e^{^t\mathrm{\Lambda }_j(t^{})𝑑t^{}}_j.`$
Note that the collision dominance approximation in Eq. (59) is not immediately applicable to the indefinite integrals associated with the inhomogeneous part of Eq. (58) since the term $`_i`$ contains the factor $`\mathrm{exp}[^t\mathrm{\Lambda }_i(t^{})𝑑t^{}]`$.
With the quantity $`B_i`$ in Eq. (36) evaluating to
$$B_i=2K_\alpha \frac{\lambda \mathrm{\Lambda }_i(2D+\mathrm{\Lambda }_i)}{\beta D}\underset{ji}{}\frac{2D+\mathrm{\Lambda }_j}{\mathrm{\Lambda }_i\mathrm{\Lambda }_j},$$
(65)
courtesy of Eq. (31), Eq. (64) simplifies to
$`P_y(t)=`$ $`{\displaystyle \frac{\beta D}{D^2+\lambda ^2}}P_z(t)e^{^t(Di\lambda )𝑑t^{}}{\displaystyle _0^t}e^{^t^{}(Di\lambda )𝑑t^{\prime \prime }}{\displaystyle \frac{\beta DK_\alpha }{D+i\lambda }}𝑑t^{}`$ (68)
$`e^{^t(D+i\lambda )𝑑t^{}}{\displaystyle _0^t}e^{^t^{}(D+i\lambda )𝑑t^{\prime \prime }}{\displaystyle \frac{\beta DK_\alpha }{Di\lambda }}𝑑t^{}`$
$`+{\displaystyle \frac{4D}{\lambda }}e^{^t2D𝑑t^{}}{\displaystyle _0^t}e^{^t^{}2D𝑑t^{\prime \prime }}{\displaystyle \frac{\beta \lambda DK_\alpha }{D^2+\lambda ^2}}𝑑t^{}+𝒪(\beta ^3),`$
as a power series in the small expansion parameter $`\beta `$. Clearly, the first term in the above expression is simply the standard adiabatic result given by Eq. (17).
The second item in Eq. (68) involving the inhomogeneous parameter $`K_\alpha `$ may be partially solved using integration by parts,
$`e^{^t(Di\lambda )𝑑t^{}}{\displaystyle _0^t}e^{^t^{}(Di\lambda )𝑑t^{\prime \prime }}{\displaystyle \frac{\beta DK_\alpha }{D+i\lambda }}𝑑t^{}`$ (69)
$`=`$ $`e^{^t(Di\lambda )𝑑t^{}}{\displaystyle _0^t}(Di\lambda )e^{^t^{}(Di\lambda )𝑑t^{\prime \prime }}{\displaystyle \frac{\beta DK_\alpha }{D^2+\lambda ^2}}𝑑t^{}`$ (70)
$`=`$ $`{\displaystyle \frac{\beta DK_\alpha }{D^2+\lambda ^2}}|_{t^{}=t}e^{_0^t(Di\lambda )𝑑t^{}}{\displaystyle \frac{\beta DK_\alpha }{D^2+\lambda ^2}}|_{t^{}=0}`$ (72)
$`e^{^t(Di\lambda )𝑑t^{}}{\displaystyle _0^t}e^{^t^{}(Di\lambda )𝑑t^{\prime \prime }}\left[{\displaystyle \frac{\beta D}{D^2+\lambda ^2}}{\displaystyle \frac{dK_\alpha }{dt^{}}}+K_\alpha {\displaystyle \frac{d}{dt^{}}}\left({\displaystyle \frac{\beta D}{D^2+\lambda ^2}}\right)\right]𝑑t^{}.`$
A sufficiently large damping parameter $`D`$ again allows the approximation of $`\mathrm{exp}[_0^t(Di\lambda )𝑑t^{}]0`$, thereby promptly obliterating the second term on the right hand side of the last equality. The quantity $`\frac{d}{dt}\left(\frac{\beta D}{D^2+\lambda ^2}\right)`$ in the integral turns out to be one of the elements constituting the $`4\times 4`$ matrix $`𝒰_{\text{rp}}\frac{𝒰_{\text{rp}}}{t}^1`$. Consistency with the adiabatic approximation then calls for the setting of $`\frac{d}{dt}\left(\frac{\beta D}{D^2+\lambda ^2}\right)0`$ such that solutions to $`𝐏_{\text{rp}}`$ have no explicit dependence on the time derivatives of the parameters $`D`$, $`\lambda `$ and $`\beta `$.
Further exploitation of the powerful technique of integration by parts on Eq. (69) in conjunction with the aforementioned approximations thus yields
$`e^{^t(Di\lambda )𝑑t^{}}{\displaystyle _0^t}e^{^t^{}(Di\lambda )𝑑t^{\prime \prime }}{\displaystyle \frac{\beta DK_\alpha }{D+i\lambda }}𝑑t^{}`$ (73)
$``$ $`{\displaystyle \frac{\beta DK_\alpha }{D^2+\lambda ^2}}|_{t^{}=t}+\left[{\displaystyle \underset{n=1}{\overset{N}{}}}\left({\displaystyle \frac{1}{Di\lambda }}\right)^n{\displaystyle \frac{\beta D}{D^2+\lambda ^2}}{\displaystyle \frac{d^{(n)}}{dt^{(n)}}}K_\alpha \right]_{t^{}=t}`$ (75)
$`e^{^t(Di\lambda )𝑑t^{}}{\displaystyle _0^t}e^{^t^{}(Di\lambda )𝑑t^{\prime \prime }}\left({\displaystyle \frac{1}{Di\lambda }}\right)^N{\displaystyle \frac{\beta D}{D^2+\lambda ^2}}{\displaystyle \frac{d^{(N+1)}}{dt^{(N+1)}}}K_\alpha 𝑑t^{}.`$
Above the neutrino decoupling temperature, the function $`K_\alpha `$ is well approximated by
$$K_\alpha 1+\frac{12\zeta (3)}{\pi ^2}\frac{e^{p/T}}{1+e^{p/T}}L_\alpha ,$$
(76)
for small $`L_\alpha `$, and consequently,
$$\frac{d^{(n)}}{dt^{(n)}}K_\alpha \frac{12\zeta (3)}{\pi ^2}\frac{e^{p/T}}{1+e^{p/T}}\frac{d^{(n)}}{dt^{(n)}}L_\alpha ,$$
(77)
with the recognition that the dimensionless factor $`e^{p/T}/(1+e^{p/T})`$ is independent of time. Since the quantity $`\frac{dL_\alpha }{dt}`$ is of the order of $`\beta ^2`$, the first term in the sum containing this is equivalently an $`𝒪(\beta ^3)`$ quantity that is generally negligible. The remainder of the sum may be similarly ignored as higher order time derivatives of $`L_\alpha `$ are successively smaller by factors of $`\beta ^2`$, leading to further simplification of Eq. (73) to
$`e^{^t(Di\lambda )𝑑t^{}}{\displaystyle _0^t}e^{^t^{}(Di\lambda )𝑑t^{\prime \prime }}{\displaystyle \frac{\beta DK_\alpha }{D+i\lambda }}𝑑t^{}`$ (78)
$``$ $`{\displaystyle \frac{\beta DK_\alpha }{D^2+\lambda ^2}}|_{t^{}=t}+𝒪(\beta ^3),`$ (79)
and likewise for its complex conjugate in Eq. (68).
The last real term in the Eq. (68) may be shown to reduce to
$`{\displaystyle \frac{4D}{\lambda }}e^{^t2D𝑑t^{}}{\displaystyle _0^t}e^{^t^{}2D𝑑t^{\prime \prime }}{\displaystyle \frac{\beta \lambda DK_\alpha }{D^2+\lambda ^2}}𝑑t^{}`$ (80)
$``$ $`{\displaystyle \frac{2\beta DK_\alpha }{D^2+\lambda ^2}}|_{t^{}=t}+\left[{\displaystyle \underset{n=1}{\overset{N}{}}}\left({\displaystyle \frac{1}{2D}}\right)^n{\displaystyle \frac{2\beta D}{D^2+\lambda ^2}}{\displaystyle \frac{d^{(n)}}{dt^{(n)}}}K_\alpha \right]_{t^{}=t}`$ (82)
$`{\displaystyle \frac{2D}{\lambda }}e^{^t2D𝑑t^{}}{\displaystyle _0^t}e^{^t^{}2D𝑑t^{\prime \prime }}\left({\displaystyle \frac{1}{2D}}\right)^N{\displaystyle \frac{\beta \lambda }{D^2+\lambda ^2}}{\displaystyle \frac{d^{(N+1)}}{dt^{(N+1)}}}K_\alpha 𝑑t^{}`$
$``$ $`{\displaystyle \frac{2\beta DK_\alpha }{D^2+\lambda ^2}}|_{t^{}=t}+𝒪(\beta ^3),`$ (83)
by similar arguments such that Eqs. (73) and (80) combine to produce
$$P_y(t)\frac{\beta D}{D^2+\lambda ^2}P_z(t)+𝒪(\beta ^3).$$
(84)
Equation (84) clearly shows the miraculous cancellation between terms pertaining to a nonzero repopulation function originally present in Eq. (68) to first order in $`\beta `$. The final expression for $`P_y(t)`$ is therefore identical to the standard adiabatic result given by Eq. (17) in the limit of interest.
Our next task is to verify that the assumption of instantaneous repopulation in Eq. (19) is indeed correct. For this purpose, we begin by expressing $`P_z`$ as a function of $`P_0`$ by way of Eq. (62),
$$P_z(t)=\left(1+\frac{\mathrm{\Lambda }_3}{D}\right)P_0(t)+\underset{j=1,2,4}{}\frac{\beta (D+\mathrm{\Lambda }_j)(\mathrm{\Lambda }_3\mathrm{\Lambda }_j)}{\lambda \mathrm{\Lambda }_j(2D+\mathrm{\Lambda }_j)}e^{^t\mathrm{\Lambda }_j(t^{})𝑑t^{}}_i.$$
(85)
To the lowest order in $`\beta `$, Eq. (85) is equivalently
$$P_z(t)=P_0(t)+\frac{4(D^2+\lambda ^2)}{\beta \lambda }e^{^t2D𝑑t^{}}_0^te^{^t^{}2D𝑑t^{\prime \prime }}\frac{\beta \lambda DK_\alpha }{D^2+\lambda ^2}𝑑t^{}+𝒪(\beta ^2),$$
(86)
on which we apply the same technique of integration by parts in conjunction with the rationale employed previously to eliminate the various time derivatives to obtain
$`P_z(t)`$ $``$ $`P_0(t)+2K_\alpha (t)+2\left[{\displaystyle \underset{n=1}{\overset{N}{}}}\left({\displaystyle \frac{1}{2D}}\right)^n{\displaystyle \frac{d^{(n)}}{dt^{(n)}}}K_\alpha \right]_{t^{}=t}`$ (88)
$`{\displaystyle \frac{2(D^2+\lambda ^2)}{\beta \lambda }}e^{^t2D𝑑t^{}}{\displaystyle _0^t}e^{^t^{}2D𝑑t^{\prime \prime }}\left({\displaystyle \frac{1}{2D}}\right)^N{\displaystyle \frac{\beta \lambda }{D^2+\lambda ^2}}{\displaystyle \frac{d^{(N+1)}}{dt^{(N+1)}}}K_\alpha 𝑑t^{}`$
$``$ $`P_0(t)+2K_\alpha (t)+𝒪(\beta ^2).`$ (89)
The last approximate equality in Eq. (88) is identically
$$\frac{N_\alpha (p)}{N^{\text{eq}}(p,0)}K_\alpha \frac{N^{\text{eq}}(p,\mu )}{N^{\text{eq}}(p,0)},$$
(90)
which clearly affirms the validity of the instantaneous repopulation assumption to the lowest order in $`\beta `$.
## V Conclusion
We have taken the exact four-component QKEs for a two-flavour active–sterile neutrino system and derived from them an approximate evolution equation for the relic neutrino asymmetry in which the role of repopulation is methodically embedded. The consequence of including a finite refilling function is to generate higher order terms which are readily discarded in the high temperature epoch of interest, thereby yielding a rate equation identical to that found earlier in the standard adiabatic limit where the said function is taken to be negligible. The formal procedure developed in the present work to establish this result has been labelled the extended adiabatic approximation.
Frequently adopted in numerical studies with notable accuracy, the assumption of instantaneous repopulation is also shown to arise naturally in the extended adiabatic limit. We have thus furnished a rigorous justification for the superficially incompatible assumptions of a negligible refilling function and instantaneous repopulation in the regime where collisions are the dominant asymmetry generation mode.
###### Acknowledgements.
This work was supported in part by the Australian Research Council and in part by the Commonwealth of Australia’s postgraduate award scheme.
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# Extinction Transition on a Pie
## Abstract
Extinction transition of bacteria under forced rotation is analyzed in pie geometry. Under convection, separation of the radial and the azimuthal degrees of freedom is not possible, and the linearized evolution operator is diagonalized numerically. Some characteristics scales are compared with the results of recent experiments, and the “integrable” limit of the theory at narrow growth region is analyzed.
The time evolution of bacterial colonies on a petri-dish has been studied recently both theoretically and experimentally. The colony is relatively simple biological system, and its basic component, an individual bacterium, involves only “elementary” biological processes like diffusion, food consumption, multiplication, death and perhaps some interaction like chemotaxis. Studies with some bacteria strains have reported a wide variety of complex pattern formation, in most cases due to competition for food resources and chemical interaction. With uniform, not-exhaustible background of nutrients and without the presence of mutations and chemical signaling, these simple strains are suppose to invade a region of nutrient rich agar in the form of a front propagating with some typical velocity, known as the Fisher front .
Biological problems of colony growth in inhomogeneous environment and under forced convection has been modeled recently by Nelson and Shnerb and by Dahmen, Nelson and Shnerb . These studies have focused on the spectral properties of the linearized evolution operator, which becomes non-Hermitian in the presence of convection . An experiment designed to test these predictions, has been carried out recently by Neicu et al .
In the experiment, a colony of Bacillus subtilis bacteria is forced to migrate in order to “catch up” with a shielded region on the the petri-dish, where all the other parts of the dish are exposed to an ultra violet (UV) light, which (under the experimental conditions) makes the unshilded bacteria immotile. It was assumed that the adaptation of the bacterial colony to the moving shielded region has nothing to do with information processing or mutual signaling in the colony, but is attributed solely to the combined effect of “dumb” diffusion of individual bacteria and the larger growth rate under the shelter. Theoretically, it was predicted that the adaptation of the colony to the moving environment fails as the drift is faster than the Fisher front velocity, and in this case the colony lags behind the shelter and an extinction transition takes place.
In order to get the essence of the theory, let us consider a one-dimensional example, where bacteria are diffusing on a line parametrized by $`x`$, and are subject to some environmental heterogeneity which implies fluctuating growth rate. If the bacteria diffuses, multiply and are forced to migrate with some convection velocity $`v`$, the differential equation which describes the evolution of the colony is,
$$\frac{c}{t}=D\frac{^2c}{x^2}+v\frac{c}{x}+a(x)cbc^2.$$
(1)
where $`a(x)`$ is the local growth rate. When the hostile environment outside the “oasis” causes the immediate death of any bacteria, and inside the oasis there is some positive growth rate,
$$a(x)=\begin{array}{cc}a& 0xx_0\\ \mathrm{}& elsewhere.\end{array}$$
(2)
If there is no drift, the linearized version of this problem is equivalent to the (imaginary time) evolution of quantum particle in an infinite potential well, and is determined by the eigenvalues of the evolution operator, which gives a colony localized on the oasis if it has some minimal width (which scales like the width of the Fisher front), $`x_0>\pi \sqrt{D/a}`$. The introduction of a drift term into (1) is compensated by the “gauge” of the evolution (Liouville) operator eigenfunctions
$$\varphi _n=\mathrm{sin}(n\pi x/x_0)e^{\pm vx/2D}\mathrm{sin}(n\pi x/x_0)$$
(3)
together with the eigenvalues “rigid” shift
$$\mathrm{\Gamma }_n(v)=\mathrm{\Gamma }_n(v=0)\frac{v^2}{4D}.$$
(4)
The theory, thus, predicts an extinction transition as all the eigenvalues of (4) becomes negative, i.e., as $`v_c=2\sqrt{aD}O(1/x_0)`$, which is the Fisher velocity . Right above the extinction transition, only the largest growth eigenvalue (the “ground state”) is positive, hence the nonlinear interaction between eigenmodes \[the term $`bc^2`$ in (1)\] are suppressed at the transition, and the analysis is focused on the ground state of the linearized operator.
In the experiment, part of a petri-dish was shielded from the UV source, and then this shield was given a constant angular velocity with respect to the petri-dish. The corresponding convection velocity $`v(r)=\omega r`$ was chosen to interpolate between zero (at the rotation axis) and about $`2v_c`$ at the edge of the dish. It turns out that the colony indeed fails to keep rotating with the shield at about half the radius. On the other hand, the velocity profile for the bacterial density $`c(r,\omega )`$ does not equilibrate on the time scales of the experiment ($`3days)`$. In this paper, I consider the differences between the one dimensional system (1,2) and the actual experimental setup. In particular, the two-dimensional nature of the experiment and the effect of radial diffusion are considered explicitly.
In order to capture the essential physics using the simplest geometry, the same extinction transition is considered on a pie, i.e., a section of the two dimensional disc \[ See Fig. (1)\]. Although the shielded region of the experiment was not in that shape, it turns out that even in this simple geometry there is a coupling between the radial and the azimuthal degrees of freedom, and the spectrum becomes “chaotic” when convection takes place. The results for this case are, accordingly, relevant also to the more complicated geometry of the experiment.
The basic equation for the bacterial growth problem on a non-uniform substrate, in the absence of mutation and chemical interactions, is :
$$\frac{c(𝐱,𝐭)}{t}=D^2c(𝐱,𝐭)+𝐚(𝐱)𝐜(𝐱,𝐭)+𝐯𝐜\mathrm{𝐛𝐜}^\mathrm{𝟐}.$$
(5)
With no convection and homogeneous, positive $`a`$ this equation supports Fisher front propagation with velocity $`2\sqrt{Da}`$. The experimental situation corresponds to $`D10^6cm^2/s`$ and $`a10^3/s`$. The Fisher velocity is of order $`0.11\mu m/s`$, as has been found experimentally. The Fisher width, which is the characteristic scale of spatial correlations in their system, is $`\sqrt{D/a}10^2cm`$, much smaller than the petri-dish radius of few centimeters.
In cylindrical geometry, Eq. (5) takes the form,
$$\frac{c(r,\theta ,t)}{t}=D^2c(r,\theta ,t)+a(\theta )c(r,\theta ,t)+vcbc^2,$$
(6)
and for rotating petri-dish the convection term is ,
$$vc=\omega \frac{c}{\theta }.$$
(7)
Pie geometry is defined by,
$$a(\theta )=\begin{array}{cc}a& 0\theta \theta _0\\ \mathrm{}& elsewhere\end{array},$$
(8)
i.e., we have absorbing boundary conditions:
$$c(r,\theta _0,t)=c(r,0,t)=0.$$
(9)
As for the boundary conditions on the petri-dish edge at $`r=R,`$ it is reasonable to take Von-Neumann boundary and to impose the no-slip condition on the bacterial density at the surface. However, the data seems to indicate extinction of the colony at the edge of the dish. This is perhaps due to the fact that the width of the boundary layer (which is expected due to the no slip condition) is about the Fisher width, which has been shown above to be very small. Accordingly, we further simplify the problem by using,
$$c(R,\theta t)=0.$$
(10)
Dropping the term $`bc^2`$ at Eq. (5), one has the linearized evolution operator, and for the no-drift $`(\omega =0)`$ case, the density of bacteria at time $`t`$ is given by:
$$c(r,\theta t)=\underset{m,n}{}A_{m,n}e^{a\mathrm{\Gamma }_{m,n}}\varphi _{m,n}(r,\theta ),$$
(11)
with the eigenstates of the evolution operator,
$$\varphi _{m,n}(r,\theta )=\eta _{m,n}J_{\frac{n\pi }{\theta _0}}(r/\sqrt{D/\mathrm{\Gamma }_{m,n}})\mathrm{sin}(\frac{n\pi \theta }{\theta _0}),$$
(12)
where $`\eta _{m,n}`$ is normalization factor
$$\eta _{n,m}=\frac{2}{R\sqrt{\theta _0}}\frac{1}{J_{[\frac{n\pi }{\theta _0}+1]}(R/\sqrt{D/\mathrm{\Gamma }_{m,n}})},$$
(13)
and the constants $`A_{m,n}`$ are determined by the initial density distribution $`c(r,\theta ,t=0).`$ The eigenvalues of the Hermitian problem are:
$$\mathrm{\Gamma }_{m,n}=D\left(\frac{j_{n\pi /\theta _0}^m}{R}\right)^2,$$
(14)
where $`j_{n\pi /\theta _0}^m`$ is the m-th zero of the corresponding Bessel function.
Let us get an order of magnitude estimate for the time scales which are relevant to the experiment . The characteristic times needed for the “ground state” to control the system is given by the typical difference between two eigenvalues. In our case, since the first zeroes of the Bessel functions are of order 1, the times involved are $`\frac{R^2}{D}`$. For an experimental system with $`R0.01m`$ and $`D10^{10}m^2/s`$, the typical relaxation times are $`O(10^6sec)11days`$, which is larger than the typical time of the actual experiment.
Now I look at the non-Hermitian case, where $`\omega 0.`$ Unlike its Cartesian analogy , no simple gauge solves the problem, and separation of variables is impossible. Spanning the space of normalizable functions by a set of Hermitian eigenstates, the perturbative term $`\omega _\theta `$ mixes both quantum numbers $`m`$ and $`n`$. The matrix elements of the convection term are:
$$n,m|\omega _\theta |k,l=2\omega R^2\gamma _{nmkl}.$$
(15)
where $`\gamma _{nmkl}`$ is:
$$\gamma _{nmkl}=\{\begin{array}{cc}k+n=even& 0\\ k+n=odd& \frac{2kn}{n^2k^2}\eta _{n,m}\eta _{k,l}I_{nmkl}\end{array}$$
(16)
and
$$I_{nmkl}=_0^1J_{\frac{n\pi }{\theta _0}}(j_{n\pi /\theta _0}^my)J_{\frac{k\pi }{\theta _0}}(j_{k\pi /\theta _0}^ly)y𝑑y.$$
(17)
In order to get the eigenvalues and eigenfunctions at finite angular velocity one should diagonalize the full non-Hermitian Liouville operator, and the extinction transition takes place as the ground state (smallest) eigenvalue, $`\mathrm{\Gamma }_{1,1}`$, becomes larger than the growth rate $`a`$ on the pie.
As the rotating system is not integrable, it should be studied numerically using some computer diagonalization of the linearized evolution operator. Essentially, one should look at the ground state of this operator, since this state dominates the system close to the extinction transition.
This numerical analysis, however, may lead to erroneous results if the continuum limit is not taken carefully. In the most general case, a discretized version of a model with local growth rate and hopping between sites may be realized numerically as a matrix, where the growth rates are the coefficients on the diagonal and the hopping process gives the off-diagonal terms. As any hopping term is positive semi-definite, the only negative terms are the local growth rate, and for any finite matrix, by adding an appropriate multiplication of the unit matrix, one may get a positive semi-definite matrix with the same eigenvectors. Perron-Frobenius theorem then implies that ground state should be a nodeless, positive eigenvector. There is a simple physical interpretation to this result: since the ground state dominates the system at long times, and the number of bacteria should not become negative, Perron-Frobenius theorem should hold. Diagonalizing numerically the evolution operator, one may get a ground state with nodes, which is physically impossible.
In order to solve this problem one should carefully take the discrete limit of the continuum theory. For our case, as $`\theta `$ is discretized to quanta of $`\mathrm{\Delta }\theta `$, the hopping rate due to diffusion becomes $`D/(r^2\mathrm{\Delta }\theta ^2)`$ and the hopping rate due to the drift is $`\pm \omega /\mathrm{\Delta }\theta `$. In order to avoid the (physically impossible) negative hopping rates, one should keep $`\mathrm{\Delta }\theta `$ small enough. If the effective discretization is given by $`\mathrm{\Delta }\theta =\theta _0/n`$, one should truncate the matrix (16) only as
$$n\frac{\omega R^2\theta _0}{D}.$$
(18)
Although (18) seems to indicate that the numerical diagonalization of (15) becomes simpler as $`\theta _00,`$ this, in fact, is not the case. As the eigenvalues of the unperturbed problem are related to the zeroes of the corresponding Bessel functions, and the rotation operator admits matrix elements only between eigenvalues related to Bessel functions of different order, it is much simpler to diagonalize (15) as $`\theta _0>>0`$. As $`\theta _00,`$ the higher $`m`$ zeroes of any Bessel of order $`n`$ are smaller than the $`m=1`$ zero of the $`n+1`$ state and the condition (18) implies the diagonalization of an infinite matrix. Accordingly, I present here the numerical results for the case $`\theta _0=\pi .`$ This situation does not coincide with the experimental conditions at , but there seems to be no prevention to perform the same experiment with large shielded area.
In Fig. (2), contour plots of the ground state for different angular velocities are shown. One may identify clearly the large deviations from the ground state from its shape at $`\omega =0.`$ The largest 100 spectral points for any case are shown in Fig. (3).
Fig. (4) presents the ground state eigenvalue, $`\mathrm{\Gamma }_{0,}`$ in units of $`D/R^2,`$ as a function of the angular velocity of the dish. The extinction transition takes place as this eigenvalue is larger than the growth rate on the “pie”, $`\frac{a}{D/R^2},`$ as has been found above.
Let us show now how to get a problem equivalent to (1,2) on a rotating petri-dish. In order to do that, the geometry should be taken on a narrow shell as in Fig. (5), i.e., the boundary conditions are,
$$\begin{array}{c}c(r,\theta _0,t)=c(r,0,t)=0.\\ c(R_1,\theta _0,t)=c(R_2,0,t)=0.\end{array}$$
(19)
with $`\mathrm{\Delta }R=R_2R_1`$. In the limit $`R_1\mathrm{}`$ at constant $`n`$, the asymptotic expansion of the Bessel functions $`J_\nu `$ and $`Y_\nu `$ at large argument gives the eigenfunctions of the unperturbed Liouville operator,
$$\varphi _{m,n}(r,\theta )\frac{(\mathrm{\Gamma }_{n,m}/D)^{1/4}}{\mathrm{\Delta }R\sqrt{R_1\theta _0}}\mathrm{sin}(\frac{m\pi r}{\mathrm{\Delta }R}+\alpha _n)\mathrm{sin}(\frac{n\pi \theta }{\theta _0}),$$
(20)
where the phase $`\alpha _n`$ ensures the boundary conditions at $`R_1`$ and the eigenvalues, $`\mathrm{\Gamma }_{n,m}=D\frac{m^2\pi ^2}{(\mathrm{\Delta }R)^2}`$, are independent of $`n.`$ The matrix elements of the operator $`\omega _\theta `$ are given by
$$n,m|\omega _\theta |k,l=\omega \delta _{m,l}\gamma _{nk}.$$
(21)
with
$$\gamma _{nk}=\{\begin{array}{cc}k+n=even& 0\\ k+n=odd& \frac{2kn}{n^2k^2}\end{array}.$$
(22)
where the approximation
$`{\displaystyle _{R_1}^{R_2}}\mathrm{sin}({\displaystyle \frac{m\pi r}{\mathrm{\Delta }R}}+\alpha _n)\mathrm{sin}({\displaystyle \frac{l\pi r}{\mathrm{\Delta }R}}+\alpha _k)\sqrt{r}𝑑r`$ (23)
$`\sqrt{R_1}{\displaystyle _{R_1}^{R_2}}\mathrm{sin}({\displaystyle \frac{m\pi r}{\mathrm{\Delta }R}}+\alpha _n)\mathrm{sin}({\displaystyle \frac{l\pi r}{\mathrm{\Delta }R}}+\alpha _k)𝑑r`$ (24)
for $`\frac{\mathrm{\Delta }R}{R_1}<<1`$ has been used. Accordingly, for any $`m`$ sector, both the diagonal and the off diagonal matrix elements are identical with the corresponding one dimensional problem, and the results should be the same.
In conclusion, the mathematical problem which corresponds to the experiment by has been found to be non-integrable, and no simple gauge transformation connects the eigenvectors of the static and the dynamic problems. The actual critical velocity and ground state properties have to be studied numerically, and the limit of very narrow sector ($`\theta _00)`$ involves diverging numerical load. The time scales needed for the ground state to dominate the system are larger than the duration of the actual experiment, and this explains the observed inequilebration.
I wish to thank A. Kudrolli, D. R. Nelson and K. Dahmen for helpful discussions and comments.
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# Ferromagnetism from Undressing
## I Introduction
The concept of a quasiparticle is central to our understanding of the physics of many-electron systems. A quasiparticle is what remains of a particle (electron) after taking into account its interaction with surrounding particles. It may be understood as a particle carrying with it a ’cloud’ of other particles with which it interacts. This cloud can be visualized as ’clothing’, or ’dressing’, of the original particle, and naturally it will generally lead to an increased effective mass of the quasiparticle.
In addition to the effective mass, another central element of the concept of a quasiparticle is its ’weight’. Paradoxically, the ’weight’ of a quasiparticle is usually the $`inverse`$ of its effective mass. If a particle is heavily dressed by interactions, its effective mass is large and its ’quasiparticle weight’ is small. The ’quasiparticle weight’, $`Z`$, expresses how much of the particle still remains intact, with a well-defined energy versus momentum relation. If $`Z`$ is small, the particle has lost most of its identity; by a sum rule, this lost weight reappears elsewhere, in an incoherent background.
To recall the origin of the key connection between quasiparticle weight and its effective mass let us remember some basic concepts of many-body theory. Consider the Green’s function for a spin $`\sigma `$ electron in a many-body system
$$G(k,\tau )=<Tc_{k\sigma }(\tau )c_{k\sigma }^{}(0)>.$$
(1)
G gives the probability amplitude to find an electron of momentum $`k`$ at time $`\tau `$ after it was created at time $`0`$. Qualitatively, if the electron does not interact with anything it will be found in its entirety at time $`\tau `$ with the same momentum $`k`$ as when it was created at time 0. However if the electron interacts strongly with other degrees of freedom, we will find very little of it back at any later time. In that case the electron has lost its ’coherence’, most or all of it has been lost in an incoherent background. This is why the Green’s function in a many-body system is generally written as
$$G(k,\omega )=G_{coh}(k,\omega )+G^{}(k,\omega )$$
(2)
where $`G_{coh}`$, the coherent part of the Green’s function, represents the quasiparticle, and $`G^{}`$ describes the incoherent background. This can be understood as follows. The exact Green’s function can be written as
$$G(k,\omega )=\frac{1}{\omega ϵ_k\mathrm{\Sigma }(k,\omega )}$$
(3)
where $`\mathrm{\Sigma }`$ is the self-energy, and $`ϵ_k`$ is its kinetic enery measured from the chemical potential, inversely proportional to its mass. Assume for simplicity that the self-energy has no $`k`$-dependence, and we have for its real part for small $`\omega `$
$$\mathrm{\Sigma }_{re}(\omega )=\mathrm{\Sigma }_{re}(0)+\omega \frac{\mathrm{\Sigma }_{re}}{\omega }$$
(4)
$`\mathrm{\Sigma }_{re}(0)`$ just renormalizes the chemical potential. From general phase space arguments one knows that the imaginary part of $`\mathrm{\Sigma }`$ goes to zero as $`\omega ^2`$ for small $`\omega `$. Hence we can separate the low frequency real part of the Green’s function and obtain
$$G(k,\omega )=\frac{1}{\omega (1\frac{\mathrm{\Sigma }_{re}}{\omega })ϵ_k}+G^{}=\frac{Z}{\omega Zϵ_k}+G^{}$$
(5)
with
$$Z=\frac{1}{1\frac{\mathrm{\Sigma }_{re}}{\omega }}$$
(6)
The term $`G^{}`$ contains the imaginary part of the self-energy and gives rise to the incoherent contribution. Thus, Eq. (5) shows that the same factor $`Z`$, the wave function renormalization factor, determines the quasiparticle spectral weight and the effective mass renormalization.
This deep connection between quasiparticle weight and its effective mass is well known. However it has not been stressed in the recent literature. For example, there is a vast recent literature on the phenomenon of colossal magnetoresistance in manganites, where the transition to the ferromagnetic state is thought to be accompanied by a reduction in the carrier’s effective mass. Yet there has been no discussion to our knowledge of any corresponding change in the quasiparticle weight. Here we will make such a connection, for the manganites and for ferromagnetic metals in general.
We have pointed out in the past that within a class of model Hamiltonians both superconductivity and ferromagnetism may be understood as driven by a lowering of the carrier’s effective mass as the ordered state develops, or equivalently a lowering of the carrier’s kinetic energy, and proposed that this common aspect of the physics may be essential to both phenomena. Only recently however have we focused on the fundamental connection between the lowering of effective mass and the corresponding expected increase in the quasiparticle weight, or ’undressing’, for the case of superconductivity(, hereafter referred to as I). This connection was brought to the limelight by the beautiful experimental results and insightful analysis of Ding et al on photoemission in cuprates, as well as experimental work by Feng et al and Basov et al. Here we make this connection for the case of ferromagnetism. It leads to a remarkably simple picture of metallic ferromagnetism, and to the understanding that both ferromagnetism and superconductivity may be driven by the same physical principle: undressing.
## II The Physics
Consider the process of creating an electron of spin $``$ at site $`i`$. Imagine there is an electron of spin down at the bond connecting site $`i`$ to a neighboring site $`j`$, as shown in Fig. 1. The $``$ electron is in state $`|0>`$, its ground state, in the absence of occupation of site $`i`$. The strong Coulomb repulsion between like charges will affect the state of that bond charge. When the $``$ electron is created at site $`i`$, the $``$ electron will make a transition to one of the bond states $`|1^l>`$, the eigenstates of the bond $``$ electron in the presence of an $``$ electron at site $`i`$. Let the ground state of that manifold be $`|1>|1^0>`$, and we denote by
$$S=<0|1>$$
(7)
the overlap matrix element of the ground states of the $``$ bond electron in the absence and in the presence of an $``$ electron at site $`i`$. We can express this mathematically as
$$c_i^{}|0>|0>=|>|0>=\underset{l}{}|>|1^l><1^l|0>=|>|1>S+\underset{l0}{}|>|1^l><1^l|0>.$$
(8)
Here, the first ket denotes the electronic state of the site $`i`$, and the second ket denotes the state of the $``$ electron at the bond. The first term in Eq. (8), where the $``$ electron ends up in its ground state $`|1>`$, represents a coherent process that preserves the phase of the wave function. It is a ’diagonal transition’ in the language of small polaron theory. If instead the $``$ electron at the bond ends up in an excited state, it represents an incoherent process (non-diagonal transition). The overlap matrix element $`S`$ represents what fraction of the $``$ electron created at site $`i`$ remains coherent, and gives rise to its reduced quasiparticle ’weight’. If $`S`$ is very small it means most of the effect of the $``$ electron creation has been dissipated in incoherent processes that left the ’background’ $``$ electron in excited states.
If instead we create the $``$ electron at site $`i`$ when there is no electron at bond $`ij`$ we have simply
$$c_i^{}|0>=|>$$
(9)
and no incoherent piece is generated, because there was no ’background’ degree of freedom to excite; hence, the weight of the quasiparticle is $`1`$ in this case.
More generally, we could instead have an $``$ electron on the bond, or both $``$ and $``$ electrons. We argue that Eq. (8) still applies, with the second ket denoting the state of the total charge at the bond. If there is no charge at the bond, Eq. (9) instead applies. We define then a quasiparticle operator $`\stackrel{~}{c}_i`$ through the relation
$$c_i^{}=[1(1S)\frac{(\stackrel{~}{n}_{ij}+\stackrel{~}{n}_{ij})}{2}]\stackrel{~}{c}_i^{}$$
(10)
where $`\stackrel{~}{n}_{ij\sigma }`$ is the bond occupation number ($`0`$ or $`1`$) of the spin $`\sigma `$ electron. Eq. (10) is only the coherent part of the electron operator, as it does not generate the second part of Eq. (8) (we use the same operator notation on the left for simplicity only). Eq. (10) restates Eqs. (8) and (9), that the weight of $`c_i^{}`$ is $`1`$ if the neighboring bond is unoccupied, and it is maximally reduced to $`S<1`$ if the neighboring bond is occupied by spin $``$ and $``$ electrons.
The physics resulting from these equations is shown schematically in Fig. 2. As in I, we use an ’independent boson model’ with an Einstein oscillator to describe the coupling of the electron at site $`i`$ to the excited states of the charge (in this case the bond charge) shown in Fig. 1. The coherent part of the spectral function (quasiparticle peak, labeled q.p.) arises from the ground-state to ground-state transition of the oscillator when the electron is created at the site, and its height is the quasiparticle weight $`Z`$. $`Z`$ increases as the bond charge decreases, and correspondingly weight in the spectral function shifts from the incoherent part to the quasiparticle peak. As we will show in the next section, the bond charge decreases when spin polarization develops.
If instead of focusing on the bond charge we were to focus on the site charge the equation analogous to (10) is
$$c_i^{}=[1(1S)\stackrel{~}{n}_i]\stackrel{~}{c}_i^{}$$
(11)
where $`\stackrel{~}{n}_i`$ is the site charge occupation. It was shown in I that Eq. (11) leads to superconductivity through undressing.
Next we wish to express the bond charge in terms of electron operators. We use the operator representation
$$\stackrel{~}{n}_{ij\sigma }=\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+\stackrel{~}{c}_{j\sigma }^{}\stackrel{~}{c}_{i\sigma }$$
(12)
which has eigenvalue $`1`$ operating on the low energy bonding state
$$\frac{|\sigma >_i|0>_j+|0>_i|\sigma >_j}{\sqrt{2}}$$
(13)
and zero if the bonding state is empty. Eq. (10) is then
$$c_i^{}=[1(1S)\frac{1}{2}\underset{\sigma }{}(\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+\stackrel{~}{c}_{j\sigma }^{}\stackrel{~}{c}_{i\sigma })]\stackrel{~}{c}_i^{}$$
(14)
Eq. (14) is the analog, for the effect of creating the $``$ electron on site $`i`$ on the neighboring bond charge, to Eq. (11) for the effect of creating the $``$ electron on site $`i`$ on the site charge. Note an important difference: under a particle-hole transformation
$$c_{i\sigma }^{}(1)^ic_{i\sigma }$$
(15)
(on a bipartite lattice) Eq. (14) is invariant, while Eq. (11) changes to
$$c_i^{}=[S+(1S)\stackrel{~}{n}_i]\stackrel{~}{c}_i^{}$$
(16)
Eq. (11) implies that increasing $`electron`$ site concentration leads to increased $`dressing`$ of electrons, and conversely Eq. (16) implies that increasing $`hole`$ site concentration leads to $`undressing`$ of holes. Instead, Eq. (14) and its identical form in hole representation imply that increasing bond occupation leads to increased dressing, both for electrons and holes. This difference between the dressing effects of site and bond charges lies at the root of the difference between superconductivity and ferromagnetism.
Finally, we consider a d-dimensional hypercubic lattice and add the contributions from all the bonds connecting to a given site, and Eq. (14) becomes
$$c_{i\sigma }^{}=[1(1S)\frac{1}{2}\underset{\delta ,\sigma ^{}}{}(\stackrel{~}{c}_{i\sigma ^{}}^{}\stackrel{~}{c}_{i+\delta \sigma ^{}}+\stackrel{~}{c}_{j\sigma ^{}}^{}\stackrel{~}{c}_{i\sigma ^{}})]\stackrel{~}{c}_{i\sigma }^{}.$$
(17)
We explore its consequences in the next sections.
## III Quasiparticle Hamiltonian
Consider the kinetic energy operator on a lattice
$$H_{kin}=\underset{i,j,\sigma }{}t_{ij}c_{i\sigma }^{}c_{j\sigma }$$
(18)
Replacing the bare electron operators in Eq. (18) by the quasiparticle operators Eq. (17) yields the low energy effective Hamiltonian for quasiparticles
$`H_{kin}`$ $`=`$ $`{\displaystyle \underset{i,j,\sigma }{}}t_{ij}\times `$ (19)
$`[`$ $`1{\displaystyle \frac{(1S)}{2}}{\displaystyle \underset{\delta ,\sigma ^{}}{}}(\stackrel{~}{c}_{i\sigma ^{}}^{}\stackrel{~}{c}_{i+\delta \sigma ^{}}+\stackrel{~}{c}_{i+\delta \sigma ^{}}^{}\stackrel{~}{c}_{i\sigma ^{}})]\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }[1{\displaystyle \frac{(1S)}{2}}{\displaystyle \underset{\delta ,\sigma ^{}}{}}(\stackrel{~}{c}_{j\sigma ^{}}^{}\stackrel{~}{c}_{j+\delta \sigma ^{}}+\stackrel{~}{c}_{j+\delta \sigma ^{}}^{}\stackrel{~}{c}_{j\sigma ^{}})]`$ (20)
In expanding this expression we will ignore terms involving more than two centers for simplicity, as well as terms with more than four fermion operators. The latter can certainly be rigurously justified if the electron (or hole) density is low. Eq. (19) then becomes
$$H_{kin}=\underset{i,j,\sigma }{}t_{ij}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+(1S)\underset{<i,j>}{}t_{ij}(\underset{\sigma }{}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+h.c.)^2$$
(22)
which can also be written as
$$H_{kin}=\underset{i,j,\sigma }{}t_{ij}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+2(1S)\underset{i,j,\sigma \sigma ^{}}{}t_{ij}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }\stackrel{~}{c}_{j\sigma ^{}}^{}\stackrel{~}{c}_{i\sigma ^{}}+2(1S)\underset{i,j,\sigma \sigma ^{}}{}t_{ij}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }\stackrel{~}{c}_{i\sigma ^{}}^{}\stackrel{~}{c}_{j\sigma ^{}}$$
(23)
In the form Eq. (20a), the interaction term generated can be simply understood as bond-charge Coulomb repulsion. In the form Eq. (20b) it can be seen that the two interaction terms are precisely of the same form as the exchange and pair hopping terms that result from considering off-diagonal matrix elements of the Coulomb interaction in a tight binding representation:
$$J_{ij}=d^3rd^3r^{}\varphi _i^{}(r)\varphi _j^{}(r^{})\frac{e^2}{|rr^{}|}\varphi _i(r^{})\varphi _j(r)$$
(25)
$$J_{ij}^{}=d^3rd^3r^{}\varphi _i^{}(r)\varphi _i^{}(r^{})\frac{e^2}{|rr^{}|}\varphi _j(r^{})\varphi _j(r)$$
(26)
In general, we will have $`J=J^{}`$ from Eq. (21) if the wavefunctions can be assumed to be real. Eq. (20) implies
$$J_{ij}=J_{ij}^{}=2t_{ij}(1S)$$
(27)
Supplementing the kinetic energy with an on-site Coulomb repulsion leads to the low energy effective Hamiltonian
$`H`$ $`=`$ $`{\displaystyle \underset{<ij>,\sigma }{}}t_{ij}(\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+h.c.)+U{\displaystyle \underset{i}{}}\stackrel{~}{n}_i\stackrel{~}{n}_i`$ (28)
$`+`$ $`{\displaystyle \underset{i,j,\sigma \sigma ^{}}{}}J_{ij}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }\stackrel{~}{c}_{j\sigma ^{}}^{}\stackrel{~}{c}_{i\sigma ^{}}+{\displaystyle \underset{i,j,\sigma \sigma ^{}}{}}J_{ij}^{}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }\stackrel{~}{c}_{i\sigma ^{}}^{}\stackrel{~}{c}_{j\sigma ^{}}`$ (29)
to describe the dynamics of the quasiparticles.
We have extensively studied the properties of this Hamiltonian for nearest neighbor hopping $`t_{ij}=t`$ and interactions $`J_{ij}=J`$, $`J_{ij}^{}=J^{}`$, in particular for the cases $`J^{}=0`$ and $`J^{}=J`$. The ’exchange term’ involving $`J`$ can also be written as
$$H_J=2J\underset{<ij>}{}(\stackrel{}{S}_i\stackrel{}{S}_j+\frac{1}{4}n_in_j)$$
(31)
with
$$(\stackrel{}{S_i})_\alpha =\frac{1}{2}(c_i^{},c_i^{})\sigma _\alpha \left(\begin{array}{c}c_i\\ c_i\end{array}\right)$$
(32)
and $`\sigma _\alpha `$ a Pauli matrix ($`\alpha =x,y,z`$). In the form Eq. (24) it looks like a ’Heisenberg exchange’ term. However, as emphasized earlier, the origin of ferromagnetism here is $`not`$ quantum-mechanical exchange of localized spins, as in Heisenberg’s case. The combination of the $`J`$ and $`J^{}`$ terms in the form Eq. (20a) displays the origin of these interactions as bond-charge Coulomb repulsion. Ferromagnetism in this model is driven by reduction of bond-charge Coulomb repulsion as spin polarization develops and accompanying kinetic energy lowering, rather than quantum-mechanical exchange.
The properties of the model Eq. (23) for nearest neighbor hoppings and interactions
$$H=t\underset{<ij>,\sigma }{}(\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+h.c.)+U\underset{i}{}\stackrel{~}{n}_i\stackrel{~}{n}_i+J\underset{i,j,\sigma \sigma ^{}}{}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }\stackrel{~}{c}_{j\sigma ^{}}^{}\stackrel{~}{c}_{i\sigma ^{}}+J^{}\underset{i,j,\sigma \sigma ^{}}{}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{i,\sigma ^{}}^{}\stackrel{~}{c}_{j\sigma ^{}}\stackrel{~}{c}_{j\sigma }$$
(33)
are similar for different values of $`J^{}/J`$. In particular, within mean field theory and for a model with constant density of states, the conditions on the parameters for ferromagnetism to occur are
$$j>\frac{1u}{2m^2(1n)^2}$$
(35)
$$j>\frac{1u}{\frac{5}{3}\frac{3}{2}m^2\frac{1}{2}(1n)^2}$$
(36)
for $`J^{}=0`$ and for $`J^{}=J`$ respectively, with
$$u=U/D$$
(38)
$$j=zJ/D$$
(39)
$$j^{}=zJ^{}/D$$
(40)
Here, $`m`$ is the magnetization per site, $`n`$ the total occupation per site, $`z`$ the number of nearest neighbors to a site and $`D=2zt`$ the bare bandwidth. In particular, for the half-filled band ($`n=1`$) the condition for full spin polarization ($`m=n`$) is
$$j=\frac{J}{2t}>1u$$
(41)
in both cases. For $`U=0`$, this condition is achieved in the limit $`S0`$ according to Eq. (22), while for increasingly larger $`U`$ smaller values of $`J`$ are required, and hence larger values of $`S`$ are sufficient. Exact diagonalization studies of the Hamiltonian Eq. (25) show that the mean field conditions Eq. (26) are qualitatively correct and reasonably accurate particularly for the half-filled band and not too large values of $`U`$.
In what follows we will for simplicity consider the model with $`J`$ only:
$$H=t\underset{<ij>,\sigma }{}(\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }+h.c.)+U\underset{i}{}\stackrel{~}{n}_i\stackrel{~}{n}_i+J\underset{i,j,\sigma \sigma ^{}}{}\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{j\sigma }\stackrel{~}{c}_{j\sigma ^{}}^{}\stackrel{~}{c}_{i\sigma ^{}}$$
(42)
One argument for dropping the $`J^{}`$ term is that its importance is suppressed due to on-site Coulomb repulsion. However, in treating the Hamiltonian in mean field theory this effect may not be properly taken into account. The effective kinetic energy that results from Eq. (29) within mean field theory is
$$ϵ_{k\sigma }=(12j(I_{}+I_{}))ϵ_k$$
(44)
with
$$I_\sigma =<\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{i+\delta ,\sigma }>$$
(45)
(one half of) the average bond charge for spin $`\sigma `$. $`ϵ_k`$ is the Fourier transform of the bare hopping amplitude $`t_{ij}`$. For a model with a constant density of states, the average bond charge at zero temperature is given by
$$I_{}+I_{}I=\frac{1m^2(1n)^2}{2}$$
(46)
so that it decreases with increasing magnetization, as expected.
The properties of the model are simplest in the half-filled band case, and we will restrict ourselves to that case in what follows. For that case,
$$I_{}=I_{}=\frac{I}{2}$$
(47)
even in the presence of spin polarization. As the temperature is lowered below $`T_c`$ and ferromagnetism develops, the average bond charge decreases for both majority and minority spins, and the effective bandwidth
$$D_{eff}=(12jI)D$$
(48)
broadens. In the normal state, as the temperature decreases the bond charge occupation Eq. (30b) increases and hence the bandwidth narrows. The temperature and magnetization dependence of the bond charge leads to a variety of interesting properties of the mean field solution of this model that are not found in the Stoner model, i.e. the mean field solution of the repulsive Hubbard model, and that describe experimental observations, as discussed in the references.
What is the significance of having derived the Hamiltonians Eq. (23) or Eq. (25) in this new way? Twofold. First, if one assumes that the interactions $`J`$ and $`J^{}`$ in Eq. (23) arise from off-diagonal matrix elements of the Coulomb interaction as given by Eq. (21), their value is expected to be rather small. This is because one has to use properly orthogonalized atomic orbitals in Eq. (21). When the Mulliken approximation holds, which is usually the case, off-diagonal matrix elements such as $`J`$ and $`J^{}`$ are very small for orthogonalized orbitals, and this is the justification for the ’zero differential overlap’ approximation in quantum chemistry. However, there is a more fundamental reason why the present derivation of the Hamiltonian Eq. (23) is more satisfactory than the one using the Coulomb matrix elements argument. This is discussed in the next section.
## IV Ferromagnetism and spectral weight transfer
The single particle Green’s function (for spin $``$ electrons) is given by
$$G_{ij}(\tau )=<Tc_i(\tau )c_j^{}(0)>=G_{ij}^{coh}(\tau )+G_{ij}^{incoh}(\tau )$$
(49)
with $`T`$ the time ordering operator. The coherent and incoherent parts of the Green’s function arise from the first and other terms in Eq. (8) respectively. For the coherent part, we replace the electron operators in terms of quasiparticle operators and obtain
$$G_{ij}^{coh}(\tau )=<T[1\frac{(1S)}{2}\stackrel{~}{n}_{bond,i}(\tau )]\stackrel{~}{c}_i(\tau )\stackrel{~}{c}_i^{}(0)[1\frac{(1S)}{2}\stackrel{~}{n}_{bond,j}(0)]>$$
(50)
where $`n_{bond,i}`$ represents the bond charge adjacent to site $`i`$. A mean field decoupling leads to
$$G_{ij}^{coh}(\tau )=[1(1S)<\stackrel{~}{n}_{bond}>]<T\stackrel{~}{c}_i(\tau )\stackrel{~}{c}_i^{}(0)>Z<T\stackrel{~}{c}_i(\tau )\stackrel{~}{c}_i^{}(0)>$$
(51)
Equation (36) defines the quasiparticle weight $`Z`$. We take for the average bond charge
$$<\stackrel{~}{n}_{bond}>=\underset{\sigma }{}<c_{i\sigma }^{}c_{j\sigma }+h.c.>=2(I_{}+I_{})=2I$$
(53)
which can be written as
$$I=_{D/2}^{D/2}dϵg(ϵ)(\frac{ϵ}{D/2})[f(ϵ_{}(ϵ)+f(ϵ_{}(ϵ)]I(T,m)$$
(54)
and is a function of temperature and magnetization. Here, $`ϵ_\sigma (ϵ)`$ are the quasiparticle energies and $`g(ϵ)`$ the density of states. Hence the quasiparticle weight is simply
$$Z=Z(T,m)=12(1S)I(T,m)=12jI(T,m)$$
(55)
and will depend on temperature and magnetization through the temperature and magnetization dependence of the bond charge. Note that for a more general case with a non-half-filled band, one would have different quasiparticle weights $`Z_\sigma `$ for spin up and down electrons in the spin-polarized state. Such a situation is also easily treated within this framework.
We consider the mean field solution of the model Eq. (29). The quasiparticle energies are given by
$$ϵ_{k\sigma }=ϵ_\sigma (ϵ_k)=(12jI)ϵ_k\sigma \frac{U+Jz}{2}m\mu $$
(56)
where the magnetization $`m`$ and chemical potential $`\mu `$ are determined by the conditions
$$m=_{D/2}^{D/2}dϵg(ϵ)[f(ϵ_{}(ϵ)f(ϵ_{}(ϵ)]$$
(58)
$$n=_{D/2}^{D/2}dϵg(ϵ)[f(ϵ_{}(ϵ)+f(ϵ_{}(ϵ)]$$
(59)
with $`n`$ the carrier concentration. From Eq. (39), the effective bandwidth is given by
$$D_{eff}=(12jI)D=[12(1S)I]D$$
(60)
where for the last equality we have used Eqs. (22) and (27b). The effective mass is given by
$$\frac{m^{}}{m_0}=\frac{1}{12jI}$$
(61)
where $`m_0`$ is the bare mass determined by the bare hopping amplitude in Eq. (18). From Eqs. (38) and (42) we have simply
$$\frac{m^{}}{m}=\frac{1}{Z(T,m)}$$
(62)
as expected from Eq. (5).
The behavior of $`I(T,m)`$ is discussed in the references. In Fig. 3 we reproduce a representative case. As the temperature is lowered in the normal state $`I(T,m)`$ increases, and it decreases again as spin polarization develops. Correspondingly, the bandwidth $`D_{eff}`$ decreases in the normal state upon cooling and expands again as the ordered state develops; similarly the quasiparticle weight $`Z`$ decreases above $`T_c`$ as $`T`$ decreases, and increases as spin polarization develops. If no magnetization were to develop, for the parameters in Fig. 3 the effective bandwidth would shrink to zero as the temperature goes to zero.
The coherent part of the spectral function is given by
$$A_{\sigma coh}(k,\omega )=\frac{1}{\pi }ImG_{coh}(k,\omega +i\delta )=Z\delta (\omega ϵ_{k\sigma })$$
(63)
We can model the full spectral function by assuming a harmonic oscillator spectrum of frequency $`\omega _0`$ associated with the bond charge excitations at each bond in Eq. (8). The result is
$`A_\sigma (k,\omega )`$ $`=`$ $`Z\delta (\omega ϵ_{k\sigma })`$ (64)
$`+`$ $`Z{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(ln\frac{1}{Z})^l}{l!}}{\displaystyle \frac{1}{N}}{\displaystyle \underset{k^{}}{}}[n_k^{}\delta (\omega +l\omega _0ϵ_{k^{}\sigma })+(1n_k^{})\delta (\omega l\omega _0ϵ_{k^{}\sigma })]`$ (65)
which is easily seen to satisfy the sum rule
$$_{\mathrm{}}^{\mathrm{}}𝑑\omega A_\sigma (k,\omega )=1$$
(66)
and describes the transfer of spectral weight from high to low frequencies as the quasiparticle weight $`Z`$ (Eq. (38) increases when spin polarization develops. Similarly, the optical sum rule states
$$_0^{\omega _m}𝑑\omega \sigma _1(\omega )=\frac{\pi e^2n}{2m^{}}$$
(67)
for the intra-band spectral weight of the optical conductivity $`\sigma _1(\omega )`$. In Eq. (47), $`\omega _m`$ is a high frequency cutoff that excludes transitions to other bands. Using Eq. (43),
$$_0^{\omega _m}𝑑\omega \sigma _1(\omega )=\frac{\pi e^2n}{2m_0}Z(T,m)$$
(68)
so that as the system becomes more coherent with increasing $`Z`$, spectral weight is also transfered into the intra-band part of the optical conductivity. If Eq. (47) is integrated to infinity however
$$_0^{\mathrm{}}𝑑\omega \sigma _1(\omega )=\frac{\pi e^2n}{2m^{}}+_{\omega _m}^{\mathrm{}}\sigma _1^{incoh}(\omega )=\frac{\pi e^2n}{2m_0}$$
(69)
with $`m_0`$ the bare mass. Thus, the extra spectral weight that goes into intraband optical absorption has to be compensated by a corresponding decrease in the incoherent contribution so as to leave Eq. (49) invariant.
In the presence of a magnetic field the quasiparticle energies are
$$ϵ_\sigma (ϵ)=(12jI)ϵ\sigma (\frac{U+Jz}{2}m+Dh)\mu $$
(70)
with $`h`$ a dimensionless magnetic field. Increasing $`h`$ gives rise to increasing magnetization and decreasing bond charge $`I`$, as seen in Fig. 3. Hence, the quasiparticle weight increases and the effective mass decreases. The magnetoresistance in this model is given by
$$\frac{\mathrm{\Delta }\rho }{\rho }=\frac{\rho (h)\rho (0)}{\rho (0)}=2j\frac{I(T,m(h))I(T,m(0))}{12jI(t,m(h))}$$
(71)
and its behavior with temperature and magnetic field resembles that seen in ferromagnets. From the Drude form for the intra-band optical conductivity
$$\sigma _1(\omega )=\frac{ne^2}{m^{}}\frac{\tau }{1+\omega ^2\tau ^2}=\frac{ne^2}{m_0}Z(T,m)\frac{\tau }{1+\omega ^2\tau ^2}$$
(72)
we conclude that the intra-band conductivity will increase with application of a magnetic field, and correspondingly the high frequency conductivity from incoherent processes will decrease. To model both parts of the conductivity we take as a simple Ansatz the spectral density Eq. (45) ignoring the momentum dependence
$$\sigma _1(\omega )=\frac{ne^2}{m_0}Z(T,m)[\frac{\tau }{1+\omega ^2\tau ^2}+\frac{\pi }{2}\underset{l=1}{\overset{\mathrm{}}{}}\frac{(ln\frac{1}{Z})^l}{l!}\delta (\omega l\omega _0)]$$
(73)
which properly satisfies the sum rules Eqs. (48) and (49). In figure 4 we show examples of the behavior expected under variation of magnetic field and temperature. Qualitatively similar behavior is seen in the optical properties of colossal magnetoresitive manganites and of europium hexaboride.
Similarly, we expect the enhanced coherence in the ferromagnetic state to be displayed in angle-resolved photoemission experiments: under application of a magnetic field or lowering the temperature in the ferromagnetic state, quasiparticle peaks should become stronger reflecting the enhanced quasiparticle weight $`Z(T,m)`$. We will present quantitative analysis and comparison with experiment elsewhere.
## V Puzzles with the optical sum rule
The optical sum rule in tight binding models needs to be treated with some care. Consider the quasiparticle Hamiltonian Eq. (29) . According to the discussion in the previous section, within mean field theory Eq. (42) gives rise to a lowering of effective mass as spin polarization develops, hence to an increased intra-band optical spectral weight according to Eq. (47).
However, the polarization operator on the lattice is given by
$$\stackrel{}{P}=e\underset{i}{}\stackrel{}{R}_in_i$$
(74)
with $`R_i`$ the position vector for site $`i`$. The current operator (in direction $`\delta `$) is obtained from its time derivative
$$J_\delta =\frac{dP_\delta }{dt}=\frac{i}{\mathrm{}}[H,P_\delta ]$$
(75)
and is easily seen to be $`independent`$ $`of`$ $`J`$, because the exchange term in Eq. (29) carries no current. Hence the $`exact`$ intra-band optical sum rule for this Hamiltonian is
$$_0^{\omega _m}𝑑\omega \sigma _1(\omega )=\frac{\pi a^2e^2}{2\mathrm{}^2}<T_t^\delta >$$
(77)
$$<T_t^\delta >=t<\stackrel{~}{c}_{i\sigma }^{}\stackrel{~}{c}_{i+\delta \sigma }+h.c.>$$
(78)
Eq. (56) predicts that as spin polarization develops and the bond charge decreases the intra-band spectral weight will $`decrease`$. This qualitatively contradicts the prediction of mean field theory for this very same model, as well as the expectation based on the physical considerations of the previous section. We are thus led to the remarkable conclusion that the exact solution of the model Eq. (29) does worse than its mean field solution in capturing essential aspects of its physics.
The situation can be remedied to some extent by including the pair hopping term in the Hamiltonian. That term does carry a current, and the sum rule Eq. (56) becomes
$$_0^{\omega _m}𝑑\omega \sigma _1(\omega )=\frac{\pi a^2e^2}{2\mathrm{}^2}[<T_\delta ^t>+4<T_\delta ^J^{}>]$$
(80)
$$<T_\delta ^J^{}>=J^{}\underset{i,\sigma }{}<\stackrel{~}{c}_{i+\delta ,\sigma }^{}\stackrel{~}{c}_{i+\delta ,\sigma ^{}}^{}\stackrel{~}{c}_{i\sigma ^{}}^{}\stackrel{~}{c}_{i\sigma }>$$
(81)
In the normal state the expectation value Eq. (57b) is negative and as spin polarization develops it will decrease in magnitude leading to an increase in the optical spectral weight, in accordance with qualitative expectation. These considerations illustrate that the optical sum rule places severe constraints on what acceptable effective low energy Hamiltonians are and on the relative magnitude of their parameters. The subject clearly needs further investigation which is outside the scope of this paper.
## VI Conclusions
This paper started from the assumption that the electron creation operator at site $`i`$ can be represented as
$$c_i^{}=[1(1S)\stackrel{~}{n}(bond)]\stackrel{~}{c}_i^{}+incoherentpart$$
(83)
where the ’incoherent part’ contains excitations of the local bond charge, and $`S`$ describes the overlap of the bond-charge configuration ground state in the presence and absence of the $``$ electron at site $`i`$. Similarly, the work of I was based on the assumption that the operator can be represented as
$$c_i^{}=[1(1S)\stackrel{~}{n}(site)]\stackrel{~}{c}_i^{}+incoherentpart$$
(84)
More generally then we conclude that a general representation should be
$$c_i^{}=[1(1S)\stackrel{~}{n}(local)]\stackrel{~}{c}_i^{}+incoherentpart$$
(85)
where $`\stackrel{~}{n}(local)`$ includes nearby site and bond charges. We have shown in I that Eq. (58b) leads to superconductivity, and here that Eq. (58a) leads to ferromagnetism. More generally, Eq. (58c) will lead to a unified description of superconductivity and ferromagnetism where one or the other (or neither) will dominate depending on characteristics of the system such as nature of the orbitals, lattice structure and band filling.
In connection with superconductivity we note that in I the nature of the boson degree of freedom was left somewhat unspecified. The analogy with the situation discussed here makes it unambiguous that the boson degree of freedom there is also electronic, as described e.g. by the electronic model with two orbitals per site, rather than e.g. a high frequency phonon.
The physics of ferromagnetism that results from Eq. (58a), and that of superconductivity that results from Eq. (58b), are remarkably alike: onset of the ordered state leads to increased quasiparticle coherence. Spectral weight, both in the one-particle properties such as photoemission, and in two-particle properties such as optical absorption, is transfered from the incoherent high frequency background to the low frequency coherent response. Quasiparticles undress, and become more like free particles, as the ordered state develops.
Remarkably, the quasiparticle Hamiltonians that result from Eqs. (58a) and (58b) lead respectively to the various off-diagonal matrix elements of the Coulomb interaction in a local representation: Eq. (58a) to exchange and pair hopping (or equivalently bond-bond-charge repulsion), and Eq. (58b) to correlated hopping (or bond-site charge repulsion). In our earlier work we had shown that these off-diagonal matrix elements lead to metallic ferromagnetism and hole superconductivity respectively, and that the common aspect of the physics of both instabilities induced by these interactions is effective mass reduction, which through the optical sum rule implies transfer of optical spectral weight to low frequencies. However, thinking about those interaction terms as simply derived from static matrix elements of the Coulomb interaction does not lead to an understanding of the optical spectral weight transfer process, nor to the understanding that spectral weight transfer should also occur in the one-particle Green’s function. Instead, the point of view presented in this paper and in I does.
What is the evidence that the physics discussed here takes place in ferromagnetic metals? We have already mentioned that at least in some ferromagnetic metals there is evidence for optical spectral weight transfer from high to low frequencies as the ferromagnetic state develops, either by decreasing the temperature or by increasing the magnetic field. We are not aware that evidence for undressing physics has been seen yet in the single particle spectral function of ferromagnets, as would be detected in photoemission experiments, but expect that it should be observable at least in manganites and hexaborides. For ferromagnets that are more ’metallic’ in the normal state, i.e. have higher coherence, it will be more difficult to detect these effects since the changes should be comparatively smaller.
As suggested earlier we believe the universal properties of negative magnetoresistance of ferromagnets and anomalously large decrease of resistivity below $`T_c`$ are evidence for effective mass reduction due to ’undressing’, as described by this paper, rather than for reduction of spin disorder scattering as usually assumed. The difference is $`not`$ semantics: in the Drude form for the optical conductivity, Eq. (52), changes in $`\tau `$ and in $`m^{}`$ will lead to different behavior at non-zero frequencies, and it should be possible to decide this question experimentally.
We also mention that another argument in favor of the picture of ferromagnetism discussed here and in our earlier work is the anomalous thermal expansion seen in ferromagnets below $`T_c`$. This is clearly due to reduction of the $`bond`$ $`charge`$ as the systems become ferromagnetic, and points to the importance of the bond charge in the phenomenon of ferromagnetism as described by our theory.
We believe that an experimental effort to detect the existence of undressing in itinerant electron systems that become ferromagnetic should be undertaken. There will of course be many system where such evidence may be too small to be experimentally detectable. However, if the evidence is found in a variety of different itinerant ferromagnets it will provide convincing evidence that the universal principle governing the transition to ferromagnetism in metals is undressing. The fact that it may be possible to also understand superconductivity with the same physical principle lends further support to the possibility that ’undressing’ may capture the essential physics of both phenomena.
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# Spin-triplet superconducting pairing due to local (Hund’s rule, Dirac) exchange
## Abstract
We discuss general implications of the local spin-triplet pairing among fermions induced by local ferromagnetic exchange, example of which is the Hund’s rule coupling. The quasiparticle energy and their wave function are determined for the three principal phases with the gap, which is momentum independent. We utilize the Bogolyubov-Nambu-De Gennes approach, which in the case of triplet pairing in the two-band case leads to the four-components wave function. Both gapless modes and those with an isotropic gap appear in the quasiparticle spectrum. A striking analogy with the Dirac equation is briefly explored. This type of pairing is relevant to relativistic fermions as well, since it reflects the fundamental discrete symmetry-particle interchange. A comparison with the local interband spin-singlet pairing is also made.
PACS Nos. 74.20Mn, 03.65.Pm
———————————————————–
$`{}_{}{}^{})`$ Electronic address: ufspalek@if.uj.edu.pl
I. Introduction
The discovery of superconductivity in the orbitally degenerate system $`Sr_2RuO_4`$, which is closely related to both ferromagnetic $`SrRuO_3`$ and antiferromagnetic and Mott insulating $`Ca_2RuO_4`$, poses a question about the role of short-range Coulomb and exchange interactions in stabilizing the spin-triplet superfluid state . In the case of Mott-Hubbard insulators the kinetic exchange interaction plays an essential role in stabilizing antiferromagnetism. This interaction is also instrumental in the form of real space pairing in driving the system close to the Mott-Hubbard boundary towards spin-singlet superconducting state. In the case of orbitally degenerate systems, the ferromagnetic and antiferromagnetic kinetic exchange interactions compete with each other for the number of electrons per atom $`n>1`$. Ferromagnetism (with a possible orbital ordering) usually wins for $`n1`$, whereas the antiferromagnetism takes over when $`nd`$, where $`d`$ is the orbital degeneracy. This type of competition should also be present in $`Sr_2RuO_4`$, in which $`4d^4`$ configuration of $`Ru^{4+}`$ contains two holes in $`t_{2g}`$ shell composed of nominally triply degenerate $`d_ϵ=(d_{xy},d_{yz},d_{zx})`$ orbitals. The two-dimensional antiferromagnetic spin fluctuations have been indeed observed in $`Sr_2RuO_4`$ system . From the symmetry point of view $`d_{xy}`$ does not mix with $`d_{yz}`$ and $`d_{zx}`$, so the fluctuations can be ascribed as solely due to the electrons in $`d_{xy}`$ band. The Hund’s rule coupling between $`d_{xy}`$ and the remaining two bands $`(d_{yz},d_{zx})`$ must than suppress the formation of the antiferromagnetic state. In effect, we are left with two electronic liquids: the doubly degenerate and hybridized $`d_{yz}d_{zx}`$ band containing approximately one hole and the $`d_{xy}`$ band containing the other. It must be underlined that all $`t_{2g}`$ holes are delocalized, since one observes a well defined Baber-Landau-Pomeranchuk $`(T^2)`$ contribution to the resistivity in both $`xy`$ $`(RuO_2)`$ plane and in c direction .
From what has been said above it is important to formulate first the model of local pairing represented a doubly degenerate (or almost degenerate) band coupled by the Hund’s rule and characterize the possible spin-triplet solutions induced by the Hund’s rule (ferromagnetic) exchange. This type of model has been formulated by us recently . We have shown there that sizeable (of the order of bare bandwidth) Coulomb correlations renormalize the system properties, i.e. lead to an almost localized Fermi liquid with a nonretarded real-space and spin-triplet pairing. The renormalized Fermi-liquid nature of our fermionic system will be a starting point in this paper, in which we consider basic features of the superconducting state such as the quasiparticle wave function (in the Fock space) and their energies. We list the possible solutions for our effective model with interorbital pairing. The question of coexistence of the A1 state with ferromagnetism, as well as the competition with the orbitally ordered-spin ferromagnetic state has been discussed separately . We believe that the present two-band model stands on its own ground, independently of the detailed nature of $`Sr_2RuO_4`$ superconductivity (which should include the third band and the anisotropic interband hybridization) and must be considered separately, to amplify the physical plausibility of this mechanism of spin-triplet pairing (see also the discussion at the end). This is particularly so because the Hund’s rule and associated with it ferromagnetic fluctuations represent probably the most natural determinants of spin-triplet pairing under these circumstances. Also, the present real-space pairing represents is formally analogous to the spin-singlet pairing and additionally, reflects a fundamental symmetry - the particle interchange. So, it contains fundamental physics in the sense, that the nature of the ground state, i.e. that of the spin-triplet superconductor, can appear instead of or together with an itinerant ferromagnetism.
II. Nambu-De Gennes method for the triplet pairing in the two-band case
We consider a degenerate two-band Fermi-liquid system coupled by a local triplet pairing. The corresponding effective Hamiltonian is of the simple form
$$=\underset{𝐤\sigma l=1,2}{}E_{𝐤l}a_{𝐤l\sigma }^{}a_{𝐤l\sigma }2\stackrel{~}{J}\underset{im}{}A_{im}^{}A_{im},$$
(1)
where $`E_{𝐤l}`$ are the quasiparticle energies with enhanced masses by the band narrowing factor $`q^1`$ (calculated self-consistently ) in the bands $`l=1,2`$, $`\stackrel{~}{J}Jt^2`$ is the effective Hund’s rule coupling (the local interorbital exchange), and $`t^2`$ is the probability of having interorbital local spin-triplet configurations, characterized by the creation operators $`A_1^{}=a_{il}^{}a_{il^{}}^{}`$, $`A_1^{}=a_{il}^{}a_{il^{}}^{}`$, and $`A_0^{}=\frac{1}{\sqrt{2}}(a_{il}^{}a_{il^{}}^{}+a_{il}^{}a_{il^{}}^{})`$ for $`ll^{}`$. The local exchange origin of the second term derives from the exact relation between the pairing operators in real space and the full exchange operator projecting the corresponding two-particle state onto the spin-triplet configuration.
$$\underset{m=1}{\overset{1}{}}A_{im}^{}A_{im}=𝐒_{il}𝐒_{il^{}}+\frac{3}{4}n_{il}n_{il^{}},$$
(2)
where $`𝐒_{il}`$ and $`n_{il}`$ are respectively the spin and the particle number operators for electron on site $`i`$ and orbital $`l`$. Explicitly $`n_{il}=_\sigma n_{il\sigma }`$, $`n_{il\sigma }=a_{il\sigma }^{}a_{il\sigma }`$, whereas the spin operators $`𝐒_{il}(S_{il}^+,S_{il}^{},S_{il}^z)(a_{il}^{}a_{il},a_{il}^{}a_{il},(1/2)(n_{il}n_{il}))`$. The right-hand side of (3) represents thus the full exchange operator.
After making the BCS-type approximation in the local form
$$A_{im}^{}A_{im}A_{im}^{}<A_{im}>+<A_{im}^{}>A_{im}<A_{im}^{}><A_{im}>$$
(3)
we can cast Hamiltonian (1) into the four-component form, which in the reciprocal $`(𝐤)`$ space takes the form
$$_{BCS}=\underset{𝐤}{}𝐟_𝐤^{}𝐇_𝐤𝐟_𝐤+\underset{𝐤}{}E_{𝐤2},$$
(4)
where the corresponding Nambu operators take the form: $`𝐟_𝐤^{}=(f_{𝐤1}^{},f_{𝐤1},f_{𝐤2},f_{𝐤2})`$, $`𝐟_𝐤=(𝐟_𝐤^{})^{}`$, and the Hamiltonian matrix for selected $`𝐤`$ state reads
$$𝐇_𝐤=\left(\begin{array}{cccc}E_{𝐤1}\mu ,& 0,& \mathrm{\Delta }_1,& \mathrm{\Delta }_0\\ 0,& E_{𝐤1}\mu ,& \mathrm{\Delta }_0,& \mathrm{\Delta }_1\\ \mathrm{\Delta }_1^{},& \mathrm{\Delta }_0^{},& E_{𝐤2}+\mu ,& 0\\ \mathrm{\Delta }_0^{},& \mathrm{\Delta }_1^{},& 0,& E_{𝐤2}+\mu \end{array}\right)\left(\begin{array}{cc}E_{𝐤1}\widehat{\sigma }_0,& \widehat{\mathrm{\Delta }}\\ \widehat{\mathrm{\Delta }}^{},& E_{𝐤2}\widehat{\sigma }_0\end{array}\right),$$
(5)
where $`\widehat{\sigma }_0\mathrm{𝟏}`$ is the unit $`2\times 2`$ matrix, and $`\mu `$ is the chemical potential. The superconducting gap is parametrized as $`\mathrm{\Delta }_m2\stackrel{~}{J}_𝐤<f_{𝐤1\sigma }^{}f_{𝐤2\sigma ^{}}^{}>`$, with $`m=(\sigma +\sigma ^{})/2`$, and $`\sigma ,\sigma ^{}=\pm 1`$. The $`2\times 2`$ matrix $`\widehat{\mathrm{\Delta }}`$ is parametrized in the usual form
$$\widehat{\mathrm{\Delta }}=i(𝐝\stackrel{~}{\sigma })\sigma _y=\left(\begin{array}{cc}d_x+id_y,& d_z\\ d_z,& d_x+id_y\end{array}\right),$$
(6)
where $`\stackrel{~}{\sigma }`$ is composed of the three Pauli matrices, whereas the vector $`𝐝`$ in spin space has the components $`d_x=(\mathrm{\Delta }_1\mathrm{\Delta }_1)/2`$, $`d_y=(\mathrm{\Delta }_1+\mathrm{\Delta }_1)/2`$, and $`d_z=\mathrm{\Delta }_0`$. The form (5) is a generalization of the Nambu representation to the triplet case with three, in general different, gaps $`\mathrm{\Delta }_m`$.
It is straightforward to introduce the $`4\times 4`$ Dirac matrices
$$\stackrel{~}{\beta }\left(\begin{array}{cc}\mathrm{𝟏},& 0\\ 0,& \mathrm{𝟏}\end{array}\right)\mathrm{and}\stackrel{~}{\alpha }_i=\left(\begin{array}{cc}0,& \sigma _i\\ \sigma _i,& 0\end{array}\right),$$
and then rewrite (5) for the degenerate case $`E_{𝐤1}=E_{𝐤2}`$ and for $`\mathrm{\Delta }_m=\mathrm{\Delta }_m^{}`$ in the form
$$𝐇_𝐤=\stackrel{~}{\beta }(E_𝐤\mu )+i(𝐝\stackrel{~}{\alpha })\mathrm{\Sigma }_2,$$
(7)
where
$$\mathrm{\Sigma }_2=\left(\begin{array}{cc}\mathrm{𝟎},& \sigma _y\\ \sigma _y,& \mathrm{𝟎}\end{array}\right)$$
is the $`y`$ component of the relativistic spin operator. We discuss in detail the simple situation of degenerate electrons $`(E_{𝐤1}=E_{𝐤2})`$ with a real gap $`\mathrm{\Delta }_m`$ in the next Section.
One can also look at the approach from a different prospective. Let us introduce the four component wave function for a single quasiparticle in the suprconducting phase propagating in the real space as follows
$$\widehat{\mathrm{\Psi }}(𝐱,t)=\frac{1}{\sqrt{N}}\underset{𝐤}{}\left(\begin{array}{c}\psi _{1𝐤}f_{𝐤1}\\ \psi _{2𝐤}f_{𝐤1}\\ \psi _{3𝐤}f_{𝐤2}^{}\\ \psi _{4𝐤}f_{𝐤2}^{}\end{array}\right)\mathrm{exp}\left[i\left(𝐤𝐱\frac{E_𝐤}{\mathrm{}}t\right)\right],$$
(8)
where $`\psi _{\mu 𝐤}`$ are the quasiparticle amplitudes which are determined for each eigenstate (see below). In this representation the Bogolyubov-De Gennes equation for a single quasiparticle in the superconducting states reads:
$$i\mathrm{}_t\widehat{\mathrm{\Psi }}=\stackrel{~}{\beta }\{E_𝐤(𝐤\frac{}{i})\mu \}\widehat{\mathrm{\Psi }}+i(𝐝\stackrel{~}{\alpha })\mathrm{\Sigma }_2\widehat{\mathrm{\Psi }},$$
(9)
where $`E_𝐤(𝐤\frac{}{i})`$ represents now the differential operator $`(1/i)`$ replacing the wave vector $`𝐤`$ in the dispersion relation $`E_𝐤`$ for quasiparticles. In the effective-mass approximation and in the stationary case this wave equation for quasiparticles in the superconducting phase has the following form
$$\lambda \left(\begin{array}{c}\psi _1\\ \psi _2\\ \psi _3\\ \psi _4\end{array}\right)=\left(\frac{\mathrm{}^2}{2m^{}}^2+\mu \right)\left(\begin{array}{c}\psi _1\\ \psi _2\\ \psi _3\\ \psi _4\end{array}\right)+\left(\begin{array}{c}\mathrm{\Delta }_1\psi _3+\mathrm{\Delta }_0\psi _4\\ \mathrm{\Delta }_0\psi _3+\mathrm{\Delta }_1\psi _4\\ \mathrm{\Delta }_1\psi _1+\mathrm{\Delta }_0\psi _2\\ \mathrm{\Delta }_0\psi _1+\mathrm{\Delta }_1\psi _2\end{array}\right),$$
(10)
where $`\psi _\mu \psi _\mu (𝐱)`$ and $`\lambda `$ is an eigenvalue of quasiparticle state in the superconducting state with the above 4-component wave function ($`\mathrm{\Delta }_m`$ are regarded as real). The validity of this equation goes beyond the simple solution (8), as one can include the magnetic and electric fields and other inhomogeneities if they appear on the mesoscopic or macroscopic scale. In the next Section we will use explicitly the momentum representation of Eqs.(10), as we will discuss exclusively homogeneous superconducting states. We will return to Eqs.(10) when discussing the general features of this Hamiltonian in Section IV. One should also note that finding the eigenvalues for Hamiltonian in the forms (4) or (7) can be achieved by diagonalizing of the matrix $`4\times 4`$ in general case, as discussed in analytic terms in Appendix A.
III. Superconducting states and their quasiparticles
We now discuss three principal solutions of Eq.(10) by taking $`\psi _\mu (𝐱)=\psi _\mu \mathrm{exp}(i𝐤𝐱)/\sqrt{V}`$, where $`V`$ is the system volume. We also assume that $`\mathrm{\Delta }_\mu =\mathrm{\Delta }_\mu ^{}`$, (e.g. neglect the applied magnetic fields), since we consider only spatially homogeneous solutions. Namely, rewriting Eq.(10) in components we obtain the combinations
$$\{\begin{array}{c}\lambda (\psi _1+\psi _2)=(E_𝐤\mu )(\psi _1+\psi _2)+(\mathrm{\Delta }_1+\mathrm{\Delta }_0)\psi _3+(\mathrm{\Delta }_0+\mathrm{\Delta }_1)\psi _4\\ \lambda (\psi _3+\psi _4)=(E_𝐤\mu )(\psi _3+\psi _4)+(\mathrm{\Delta }_1+\mathrm{\Delta }_0)\psi _1+(\mathrm{\Delta }_0+\mathrm{\Delta }_1)\psi _2,\end{array}$$
(11)
and
$$\{\begin{array}{c}\lambda (\psi _1\psi _2)=(E_𝐤\mu )(\psi _1\psi _2)+(\mathrm{\Delta }_1\mathrm{\Delta }_0)\psi _3+(\mathrm{\Delta }_0\mathrm{\Delta }_1)\psi _4\\ \lambda (\psi _3\psi _4)=(E_𝐤\mu )(\psi _3\psi _4)+(\mathrm{\Delta }_1\mathrm{\Delta }_0)\psi _1+(\mathrm{\Delta }_0\mathrm{\Delta }_1)\psi _2.\end{array}$$
(12)
Such combinations of particle $`(\psi _1`$ and $`\psi _2)`$ and hole ($`\psi _3`$ and $`\psi _4`$) components contain basic symmetry, as we will see on example of particular solutions, which we discuss next.
A. Isotropic solution: $`\mathrm{\Delta }_0=\mathrm{\Delta }_1=\mathrm{\Delta }_1\mathrm{\Delta }`$
In that situation Eqs.(11) - (12) take a simple form
$$\{\begin{array}{c}\lambda (\psi _1+\psi _2)=(E_𝐤\mu )(\psi _1+\psi _2)+2\mathrm{\Delta }(\psi _3+\psi _4)\\ \lambda (\psi _3+\psi _4)=(E_𝐤\mu )(\psi _3+\psi _4)+2\mathrm{\Delta }(\psi _1+\psi _2),\end{array}$$
(13)
and
$$\{\begin{array}{c}\lambda (\psi _1\psi _2)=(E_𝐤\mu )(\psi _1\psi _2)\\ \lambda (\psi _3\psi _4)=(E_𝐤\mu )(\psi _3\psi _4).\end{array}$$
(14)
The first two equations lead to the modes with a gap
$$\lambda =\lambda _{𝐤1,2}=\pm \sqrt{(E_𝐤\mu )^2+4\mathrm{\Delta }^2}\pm \lambda _𝐤.$$
(15)
For those two modes $`\psi _1=\psi _2`$ and $`\psi _3=\psi _4`$ and their eigenstates are characterized by the following quasiparticle operators
$$\alpha _𝐤=u_𝐤\frac{1}{\sqrt{2}}\left(f_{𝐤1}+f_{𝐤1}\right)v_𝐤\frac{1}{\sqrt{2}}\left(f_{𝐤2}^{}+f_{𝐤2}\right),$$
(16)
and
$$\beta _𝐤^{}=v_𝐤\frac{1}{\sqrt{2}}\left(f_{𝐤1}+f_{𝐤1}\right)+u_𝐤\frac{1}{\sqrt{2}}\left(f_{𝐤2}^{}+f_{𝐤2}\right),$$
(17)
with the Bogolyubov coherence factors
$$u_𝐤=\frac{1}{\sqrt{2}}\left(1+\frac{E_𝐤\mu }{\lambda _𝐤}\right)^{1/2},v_𝐤=\frac{1}{\sqrt{2}}\left(1\frac{E_𝐤\mu }{\lambda _𝐤}\right)^{1/2}.$$
(18)
The quasiparticle operators contain symmetric combinations $`(f_{𝐤1}+f_{𝐤1})/\sqrt{2}`$ and $`(f_{𝐤2}+f_{𝐤2})/\sqrt{2}`$. The wave function is symmetric with respect to particle-spin interchange $`()`$ and describes quasiparticle states of energy $`\pm \lambda _𝐤`$, respectively.
Eqs.(14) lead to the gapless modes of the form
$$\lambda =\lambda _{𝐤3,4}=\pm (E_𝐤\mu ),$$
(19)
and correspond to the eigenstates characterized by the operators
$$\gamma _𝐤=\frac{1}{\sqrt{2}}\left(f_{𝐤1}f_{𝐤1}\right),\mathrm{and}\delta _𝐤^{}=\frac{1}{\sqrt{2}}\left(f_{𝐤2}^{}f_{𝐤2}^{}\right)$$
(20)
and constitute the antisymmetric-in-spin operators, representing the unpaired electrons. These gapless modes disappear when the gap components are not equal, as shown in Appendix A. One should note that the gapless modes appear even though the superconducting gap here is $`𝐤`$-independent.
Combining the solutions (16) - (18) and (19 \- (20) we can express the original (”old”) particle operators in terms of quasiparticle (”new”) operators in the following manner
$$\left(\begin{array}{c}f_{𝐤1}\\ f_{𝐤1}\\ f_{𝐤2}^{}\\ f_{𝐤2}^{}\end{array}\right)=\frac{1}{\sqrt{2}}\left(\begin{array}{cccc}u_𝐤,& v_𝐤,& 1,& 0\\ u_𝐤,& v_𝐤,& 1,& 0\\ v_𝐤,& u_𝐤,& 0,& 1\\ v_𝐤,& u_𝐤,& 0,& 1\end{array}\right)\left(\begin{array}{c}\alpha _𝐤\\ \beta _𝐤^{}\\ \gamma _𝐤\\ \delta _𝐤^{}\end{array}\right).$$
(21)
This transformation is necessary for determining the self-consistent equation for the gap and for the chemical potential $`\mu `$. First, we rewrite the Hamiltonian (4) in the diagonal form
$$=\underset{𝐤}{}\lambda _𝐤(\alpha _𝐤^{}\alpha _𝐤\beta _𝐤\beta _𝐤^{})+E_𝐤(\gamma _𝐤^{}\gamma _𝐤\delta _𝐤\delta _𝐤^{})+\underset{𝐤}{}E_𝐤$$
$$=\underset{𝐤}{}\lambda _𝐤(\alpha _𝐤^{}\alpha _𝐤+\beta _𝐤^{}\beta _𝐤1)+E_𝐤(\gamma _𝐤^{}\gamma _𝐤+\delta _𝐤^{}\delta _𝐤).$$
(22)
The equation for the gap e.g. $`\mathrm{\Delta }_1=<f_{𝐤1}^{}f_{𝐤2}^{}>`$ is obtained by substituting the relevant transformed operators in (21) to $`\mathrm{\Delta }_1`$. In effect, we obtain the usual BCS form $`(E_𝐤E_𝐤\mu )`$
$$<f_{𝐤1}^{}f_{𝐤2}^{}>=\frac{1}{2}\frac{\mathrm{\Delta }}{\sqrt{E_𝐤^2+4\mathrm{\Delta }^2}}\mathrm{tanh}\left(\frac{\beta \sqrt{E_𝐤^2+4\mathrm{\Delta }^2}}{2}\right),$$
(23)
where $`\beta (k_BT)^1`$. So, the gap equation has two solutions: $`1^o\mathrm{\Delta }0,`$
$`2^o`$
$$1=\frac{J}{N}\underset{𝐤}{}\frac{1}{\sqrt{E_𝐤^2+4\mathrm{\Delta }^2}}\mathrm{tanh}\left(\frac{\beta \sqrt{E_𝐤^2+4\mathrm{\Delta }^2}}{2}\right).$$
(24)
The last equation tells us that the physical gap is $`2\mathrm{\Delta }`$. The self-consistent equation for the chemical potential must include gapless modes, i.e. takes the form
$$n=\frac{1}{N}\underset{𝐤\sigma }{}<f_{𝐤1\sigma }^{}f_{𝐤1\sigma }+f_{𝐤2\sigma }^{}f_{𝐤2\sigma }>=\frac{2}{N}\underset{𝐤}{}<\alpha _𝐤^{}\alpha _𝐤+\gamma _𝐤^{}\gamma _𝐤>.$$
(25)
Normally, as we shall see, $`\mathrm{\Delta }|\mu |`$, and hence approximately half of all particles will have the spectrum gapped. The details must be analysed numerically for a concrete structure of the density of states. In the limit $`\stackrel{~}{W}\stackrel{~}{J}`$ we have the estimate of the gap value at $`T=0`$ in the form $`\mathrm{\Delta }=(\stackrel{~}{W}/2)\mathrm{exp}(\stackrel{~}{W}/(2\stackrel{~}{J})`$; this yields the value $`\mathrm{\Delta }/\stackrel{~}{W}10^310^4`$, or in the regime $`110K`$ for $`\stackrel{~}{W}1eV`$ and $`\stackrel{~}{J}0.1\stackrel{~}{W}`$.
We need also the expression for the ground state energy, as various solutions are possible. In the present case, this energy can be written as
$$\frac{E_G}{N}=\frac{2}{N}\underset{𝐤}{}\{\sqrt{E_𝐤^2+\stackrel{~}{\mathrm{\Delta }}^2}<\alpha _𝐤^{}\alpha _𝐤>+E_𝐤<\gamma _𝐤^{}\gamma _𝐤>\sqrt{E_𝐤^2+\stackrel{~}{\mathrm{\Delta }}^2}\}+\frac{\stackrel{~}{\mathrm{\Delta }}^2}{2J},$$
(26)
where $`\stackrel{~}{\mathrm{\Delta }}=2\mathrm{\Delta }`$ .
B. Equal-spin pairing: $`\mathrm{\Delta }_00`$
To obtain the explicit solution we now combine separately the first and third components of Eq.(10) on one side, and the second and the fourth on the other. Adding and subtracting the corresponding terms we obtain:
$$\{\begin{array}{c}(\mathrm{\Delta }_1\lambda )(\psi _1+\psi _3)+(E_𝐤\mu )(\psi _1\psi _3)=0\\ (E_𝐤\mu )(\psi _1+\psi _3)(\mathrm{\Delta }_1+\lambda )(\psi _1\psi _3)=0,\end{array}$$
(27)
and
$$\{\begin{array}{c}(\mathrm{\Delta }_1\lambda )(\psi _2+\psi _4)+(E_𝐤\mu )(\psi _2\psi _4)=0\\ (E_𝐤\mu )(\psi _2+\psi _4)(\mathrm{\Delta }_1+\lambda )(\psi _2\psi _4)=0.\end{array}$$
(28)
Thus, the two pairs of components (27) and (28) separate from each other and it is sufficient to solve e.g. the first system (27) to be able to reproduce the other. Explicitly, the two solutions can be combined into the form, in which the eigenvalues take the form
$$\lambda \lambda _{𝐤1\mathrm{}4}=\pm \sqrt{(E_𝐤\mu )^2+\mathrm{\Delta }_\sigma ^2}\pm \lambda _𝐤^{(\sigma )},$$
(29)
where for each spin orientation $`\sigma =\pm 1`$ of the quasiparticles we have two solutions with the gap $`\pm \sqrt{(E_𝐤\mu )^2+\mathrm{\Delta }_\sigma ^2}`$. The quasiparticle operators $`(\alpha _{𝐤\sigma },\beta _{𝐤\sigma }^{})`$ diagonalizing Hamiltonian (4) in this case are:
$$\alpha _{𝐤\sigma }=u_𝐤^{(\sigma )}\frac{1}{\sqrt{2}}\left(f_{𝐤1\sigma }+f_{𝐤2\sigma }^{}\right)v_𝐤^{(\sigma )}\frac{1}{\sqrt{2}}\left(f_{𝐤1\sigma }f_{𝐤2\sigma }^{}\right),$$
(30)
and
$$\beta _{𝐤\sigma }^{}=v_𝐤^{(\sigma )}\frac{1}{\sqrt{2}}\left(f_{𝐤1\sigma }+f_{𝐤2\sigma }^{}\right)+u_𝐤^{(\sigma )}\frac{1}{\sqrt{2}}\left(f_{𝐤1\sigma }f_{𝐤2\sigma }^{}\right),$$
(31)
with the coherence factors
$$u_𝐤^{(\sigma )}=\frac{1}{\sqrt{2}}\left(1+\frac{\mathrm{\Delta }_\sigma }{\lambda _𝐤^{(\sigma )}}\right)^{1/2},v_𝐤^{(\sigma )}=\frac{1}{\sqrt{2}}\left(1\frac{\mathrm{\Delta }_\sigma }{\lambda _𝐤^{(\sigma )}}\right)^{1/2}.$$
(32)
In general, we have two gaps $`\mathrm{\Delta }_\sigma =(\mathrm{\Delta }_1,\mathrm{\Delta }_1)`$. In the situation $`\mathrm{\Delta }_\sigma =\mathrm{\Delta }_\sigma =\mathrm{\Delta }`$ we have a doubly (spin) degenerate solutions. It can be easily verified that the operators (30) and (31) obey the fermion anticommutation relations. The diagonalized Hamiltonian has the form
$$=\underset{𝐤}{}\lambda _𝐤^{(\sigma )}\left(\alpha _{𝐤\sigma }^{}\alpha _{𝐤\sigma }+\beta _{𝐤\sigma }^{}\beta _{𝐤\sigma }1\right)+E_0.$$
(33)
To determine the gap equation we have to find the transformation which is reverse of (30) and (31). It is of the form
$$\left(\begin{array}{c}f_{𝐤1}\\ f_{𝐤1}\\ f_{𝐤2}^{}\\ f_{𝐤2}^{}\end{array}\right)=\frac{1}{\sqrt{2}}\left(\begin{array}{cccc}u_𝐤^{(+)}+v_𝐤^{(+)},& u_𝐤^{(+)}v_𝐤^{(+)},& 0,& 0\\ 0,& 0,& u_𝐤^{()}+v_𝐤^{()},& u_𝐤^{()}v_𝐤^{()}\\ u_𝐤^{(+)}v_𝐤^{(+)},& u_𝐤^{(+)}v_𝐤^{(+)},& 0,& 0\\ 0,& 0,& u_𝐤^{()}v_𝐤^{()},& u_𝐤^{()}v_𝐤^{()}\end{array}\right)\left(\begin{array}{c}\alpha _𝐤\\ \beta _𝐤^{}\\ \alpha _𝐤\\ \beta _𝐤^{}\end{array}\right).$$
(34)
The difference with the isotropic pairing (21) is that here the coherence factors appear in combinations. Those appear also in the self-consistent equation for the gap
$$<f_{𝐤2\sigma }^{}f_{𝐤1\sigma }^{}>=(u_𝐤^{(\sigma )2}v_𝐤^{(\sigma )2})[1<\alpha _{𝐤\sigma }^{}\alpha _{𝐤\sigma }><\beta _{𝐤\sigma }^{}\beta _{𝐤\sigma }>].$$
(35)
In result, the self-consistent equation will have the following three solutions: $`1^o`$ $`\mathrm{\Delta }_\sigma 0`$ , $`2^0`$ $`\mathrm{\Delta }_{(\sigma )}=0`$, but $`\mathrm{\Delta }_\sigma 0`$ is the solution of equation
$$1=\frac{J}{N}\underset{𝐤}{}\frac{1}{\sqrt{(E_𝐤\mu )^2+\mathrm{\Delta }_\sigma ^2}}\mathrm{tanh}\left(\frac{\beta }{2}\sqrt{(E_𝐤\mu )^2+\mathrm{\Delta }_\sigma ^2}\right)$$
(36)
$`3^o`$ $`\mathrm{\Delta }_\sigma 0`$, $`\mathrm{\Delta }_\sigma 0`$, and each of them is determined from equation (36).
One should note that the coupling constant above $`(J)`$ is the same as for the isotropic phase (cf. Eq.(24)). For the sake of completness, we reproduce the ground-state-energy expression which is
$$\frac{E_G}{N}=\frac{2}{N}\underset{𝐤\sigma }{}\sqrt{(E_𝐤\mu )^2+\mathrm{\Delta }_\sigma ^2}<\alpha _{𝐤\sigma }^{}\alpha _{𝐤\sigma }>+\frac{\mathrm{\Delta }_1^2+\mathrm{\Delta }_1^2}{2J}$$
$$\frac{1}{N}\underset{𝐤\sigma }{}\sqrt{(E_\sigma \mu )^2+\mathrm{\Delta }_\sigma ^2}+\frac{1}{N}\underset{𝐤}{}E_𝐤.$$
(37)
This phase represents a starting point when discussing the coexistence of ferromagmnetism and the spin-triplet superconductivity.
C. Spin-polarized phase: $`\mathrm{\Delta }_0=\mathrm{\Delta }_{}=0`$
In this limit the system is totally spin polarized, i.e. is a spin-saturated superconductor. In that limit we recover again the spectrum both with and without gap, i.e. $`\lambda _{𝐤1,2}=\pm \sqrt{(E_𝐤\mu )^2+\mathrm{\Delta }_{}^2}`$, $`\lambda _{𝐤3,4}=\pm (E_𝐤\mu )`$. Thus paired and unpaired states coexist also in this phase, as can be easily seen from Eqs.(30) - (31), which yield the form written there for $`\sigma =`$ and $`\alpha _𝐤=f_{𝐤1}`$ and $`\beta _𝐤^{}=f_{𝐤2}^{}`$.
Summarizing cases A-C, the lowest energy will have the homogeneous state with $`\mathrm{\Delta }_{}=\mathrm{\Delta }_{}=\mathrm{\Delta }_0`$ so that the effective gap is equal to $`2\mathrm{\Delta }`$. The most interesting feature of the results is that the gapless modes coexist in cases A and C and represent half of the spectrum. Also, the condensed phases described by the cases A-C above correspond roughly to the solutions for superfluid $`{}_{}{}^{3}He`$, which are labelled B, A, and A1. However, under the present circumstances here we have momentum independent gaps, since the pairing is of the local (intrasite, but interorbital) nature. For the sake of comparison we present in Appendix B the case of spin-singlet pairing induced by the same type of local interband pairing induced by antiferromagnetic exchange.
IV. Remark on the triplet pairing for relativistic fermions
The two-band situation with a local ferromagnetic exchange can be easily generalized to the explicitly relativistic form modelling thus the triplet configuration of spin, isospin or color (the singlet case was considered by Nambu and Jona-Lasinio and in ). The paired quasiparticles obey the following modified Dirac wave equation
$$i\mathrm{}_t\mathrm{\Psi }=(c\stackrel{~}{\alpha }\widehat{𝐩}+\stackrel{~}{\beta }mc^2)\mathrm{\Psi }+i(𝐝\stackrel{~}{\alpha })\mathrm{\Sigma }_2\mathrm{\Psi },$$
(38)
where the last term supplements the Dirac equation with the contact pairing. By taking the analogy with the original approach by Nambu one can write down the stationary version of this equation as the following system in the two-component (Weyl) representation
$$\lambda \left(\begin{array}{c}\psi _1\\ \psi _2\end{array}\right)=(c\stackrel{~}{\sigma }\widehat{𝐩}\mu )\left(\begin{array}{c}\psi _1\\ \psi _2\end{array}\right)+mc^2\left(\begin{array}{c}\psi _3\\ \psi _4\end{array}\right)+\left(\begin{array}{c}\mathrm{\Delta }_1\psi _3+\mathrm{\Delta }_0\psi _4\\ \mathrm{\Delta }_0\psi _3+\mathrm{\Delta }_1\psi _4\end{array}\right),$$
(39)
$$\lambda \left(\begin{array}{c}\psi _3\\ \psi _4\end{array}\right)=(c\stackrel{~}{\sigma }\widehat{𝐩}\mu )\left(\begin{array}{c}\psi _3\\ \psi _4\end{array}\right)+mc^2\left(\begin{array}{c}\psi _1\\ \psi _2\end{array}\right)+\left(\begin{array}{c}\mathrm{\Delta }_1\psi _1+\mathrm{\Delta }_0\psi _2\\ \mathrm{\Delta }_0\psi _1+\mathrm{\Delta }_1\psi _2\end{array}\right).$$
(40)
This system of equations can be directly compared with Eq.(10) for nonrelativistic electrons. In the present situation the mass term mixes the upper and the lower two components of the bispinor, as does the pairing part. Apart from a modification in the kinetic-energy part, the system of equations (39)-(40) can be solved in the manner as discussed in Section III. The detailed discussion must include also the gauge fields, which lead to one-exchange pairing (analogous to the phonon-mediated pairing), a topic intensively discussed in recent literature . In general, the singlet pairing , \- is mutually exclusive with the triplet pairing proposed here and therefore, the latter requires first a detailed discussion of exchange interactions represented by relativistic spin $`\{\mathrm{\Sigma }_i\}`$.
V. Discussion
In this paper we have formulated the quasiparticle language for local triplet pairing between fermions (interband pairing in the nonrelativistic case) induced by the local (Hund’s rule or Dirac) exchange. In particular, we have determined explicitly the quasiparticle states and the De Gennes wave equation for them, which can be useful when considering spatially inhomogeneous situations . The principal feature of the results is the existence of the gapless modes, existence of which can also be proved on a phenomenological level . The circumstance that the pairing is induced by the ferromagnetic exchange means that this interaction can lead not only to an itinerant ferromagnetic state, but also to either spin-triplet superconductor or to a coexistence of both these states (for a brief discussion of this issue see Ref.). The present paper represents only the first step in this direction. Furthermore, our mechanism of pairing expressing the fundamental symmetry (the particle interchange) may have an important astrophysical application: the pairing in the neutron-proton matter in pulsar, but this intriguing possibility requires a separate study.
Two problems should be tackled next. First, the analysis of the Meissner effect, since in the present situation the orbital diamagnetism will compete with the ferromagnetic spin polarization (particularly, if the triplet superconducting and ferromagnetic phases can coexist). An intriguing question here is: can we reach the limiting superconducting phase (corresponding to A1 phase in the case of superfluid $`{}_{}{}^{3}He`$), the critical temperature $`T_c`$ of which can be enhanced by the applied magnetic field?
Second, one should derive microscopically the Ginzburg-Landau equation for the condensed pairs. Note that the De Gennes equation (9) or (10) is useful in describing the quasiparticle tunneling, whereas the Ginzburg-Landau equation is useful when considering the Josephon (pair) tunneling. Here an intriguing question to what extent the gapless quasiparticles influence the tunneling between the spin-singlet and the spin-triplet superconductors or between the triplet superconductor and normal metal. We should be able to see the progress in answering those questions in a near future.
Finally, returning to the question of the nature of the paired state in $`Sr_2RuO_4`$ one can make the following two remarks. First, the existence of gapless modes in B- and A1-like phases leads to the persistence of the linear term in the specific heat in the superconduncting phase at its $`50\%`$ value, if the pairing is the pure spin-triplet state of electrons pairs derived from $`d_{zx}`$ bands. The recent measurements in very pure samples contradict such earlier claims that a half of the linear specific heat survives in the superconducting phases. Does that mean that the full phase diagram involves more than one phase depending on the doping degree, as in the heavy-fermion system $`U_{1x}Th_xBe_{13}`$ ? In connection with this one can say that because of the reasons mentioned in Section I it is conceivable that a singlet pairing in $`d_{xy}`$ band induced by antiferromagnetic fluctuations can compete in the triplet state in the other two bands . The nature of the resultant state should be determined then. We should be able to see a progress in those matters in near future.
Second, an important question concerns the nature of the pairing potential. In more standard approach , one introduces the effective triplet pairing via the paramagnon exchange. In that situation the coupling constant is determined by the susceptibility $`\chi (𝐪=𝐤𝐤^{})`$ and hence, is wave vector ($`𝐪`$) dependent. In the approach developed here the exchange interaction itself provides real space pairing, as in the case of high temperature superconductors . In the case of Hund’s rule coupling the pairing potential is then $`𝐤`$-independent. We believe that the latter approach is relevant when the particles are strongly correlated. $`Sr_2RuO_4`$ is a systems close to (but below) the Mott-Hubbard localization threshold, i.e. the halfway between the weak-correlation and the strong-correlation asymptotic regimes. Therefore, the real space pairing is certainly worth of analysing, as it allows for an analytic approach.
One methodological remark at the end. In the analytical analysis of the spin-triplet pairing one usually uses , the $`𝐝`$ vector in expressing the pairing part. Here we decided to use the original BCS gap parameters, a completely equivalent procedure, but probably a bit more direct, at least in the spatially homogenous situation. In connection with this a difference of the present description with that for superfluid $`helium3`$ should be stressed. Namely, in the $`helium3`$ case, the $`L=1`$ orbital moment $`𝐥`$ and the spin vector $`𝐝`$ determine the (many-component) nature of the gap. The $`𝐝.𝐥`$ and $`\mathrm{𝐝𝐱𝐥}`$ combinations determine the order-parameter dynamics. Here, no $`𝐥`$ vector appears and therefore, the order parameter can have up to 3 independent components.
Acknowledgments
I am grateful to my student Andrzej Klejnberg for the cooperation, which lead to this effective model. I am also grateful to Włodek Wójcik for many discussions and a technical help. Additionally, I was partially motivated by the question posed by Mark Jarrell from the University of Cincinnati, who asked if a simpler model representation of the work presented in Ref.12 was possible. The work was supported by KBN, Grant No. 2P03B 092 18.
Appendix A: General solution: $`\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_1`$, $`E_{𝐤1}E_{𝐤2}`$
The most general case of finding the eigenvalues for quasiparticles in the superconducting phase is to diagonalize matrix $`4\times 4`$, i.e. to solve the equation (we assume that $`E_{𝐤1,2}\mu E_𝐤`$ and that $`\mathrm{\Delta }_1=\mathrm{\Delta }_1^{}`$)
$$det\left(\begin{array}{cccc}E_{𝐤1}\lambda ,& 0,& \mathrm{\Delta }_1,& \mathrm{\Delta }_0\\ 0,& E_{𝐤1}\lambda ,& \mathrm{\Delta }_0,& \mathrm{\Delta }_1\\ \mathrm{\Delta }_1,& \mathrm{\Delta }_0,& E_{𝐤2}\lambda ,& 0\\ \mathrm{\Delta }_0,& \mathrm{\Delta }_1,& 0,& E_{𝐤2}\lambda \end{array}\right)=0.$$
(41)
By a straightforward evaluation one obtains the eigenvalues
$$\lambda _{𝐤\mathrm{1..4}}=\frac{1}{2}(E_{𝐤1}E_{𝐤2})\pm \frac{1}{2}\left[(E_{𝐤1}+E_{𝐤2})^2+2\stackrel{~}{\mathrm{\Delta }}^2\pm 2\sqrt{\stackrel{~}{\mathrm{\Delta }}^44\stackrel{~}{\delta }^4}\right]^{1/2},$$
(42)
where
$$\stackrel{~}{\mathrm{\Delta }}=(\mathrm{\Delta }_1^2+\mathrm{\Delta }_1^2+2\mathrm{\Delta }_0^2)^{1/2},$$
(43)
is the total gap, and
$$\stackrel{~}{\delta }=\sqrt{|\mathrm{\Delta }_0^2\mathrm{\Delta }_1\mathrm{\Delta }_1|},$$
(44)
is its anisotropy in fermion-pair spin space. In this general case all four modes are each with a different gap and the results reduce nicely to the eigenvalues discussed as cases A-C in main text. In general, the superconducting coupling at the level of energies amounts to hybridizing the different fermion fields $`(l=1,2)`$ and their spin $`(\sigma =,)`$ states. The general form of the eigenstates can also be found in a straightforward manner, but will not be discussed here.
Appendix B: Local spin-singlet pairing in two-band case
For the sake of comparison with the spin-triplet case we outline here the solution for the corresponding spin-singlet situation. In this case the Hamiltonian with the spin-singlet exchange coupling has the form in the real space
$$=\underset{𝐤\sigma l=1,2}{}E_{𝐤l}n_{𝐤l\sigma }+J\underset{ill^{}}{}(𝐒_{il}𝐒_{il^{}}+\frac{1}{4}n_{il}n_{il^{}}).$$
(45)
Note that now the exchange operator in the present situation differs from (2) introduced in the triplet case. Introducing the corresponding local pairing operators in the singlet state
$$B_i^{}=\frac{1}{\sqrt{2}}\left(a_{i1}^{}a_{i2}^{}a_{i1}^{}a_{i2}^{}\right)$$
(46)
we can write down the second term in (45) as $`2J_iB_i^{}B_i`$.
After taking the space Fourier transform, and including only $`(𝐤,𝐤)`$ pairs we obtain
$$=\underset{𝐤\sigma }{}E_{𝐤l}n_{𝐤l\sigma }\frac{J}{N}\underset{\mathrm{𝐤𝐤}^{}}{}\left(f_{𝐤1}^{}f_{𝐤2}^{}f_{𝐤1}^{}f_{𝐤2}^{}\right)\left(f_{𝐤2}f_{𝐤1}f_{𝐤2}f_{𝐤1}\right).$$
(47)
Making subsequently, as in Section II the BCS approximation, we obtain (the chemical potential is included in $`E_{𝐤l}E_{𝐤l}\mu `$)
$$_{BCS}=\underset{𝐤l\sigma }{}\{E_{𝐤l}n_{𝐤l\sigma }+\mathrm{\Delta }_𝐤^{}(f_{𝐤1}^{}f_{𝐤2}^{}f_{𝐤1}^{}f_{𝐤2}^{})+H.C.\}\frac{\mathrm{\Delta }^2}{2J}N,$$
(48)
where
$$\mathrm{\Delta }\frac{2J}{N}\underset{𝐤}{}<f_{𝐤1}^{}f_{𝐤2}^{}>.$$
(49)
Introducing, as before, the four dimensional vectors $`𝐟^{}(f_{𝐤1}^{},f_{𝐤1}^{},f_{𝐤2},f_{𝐤2})`$ and their conjugate as one column vectors, we can write down the Hamiltonian in the form of the following $`4\times 4`$ matrix
$$_{BCS}=E_0+\left(\begin{array}{cccc}f_{𝐤1}^{},& f_{𝐤1}^{},& f_{𝐤2},& f_{𝐤2}\end{array}\right)\left(\begin{array}{cccc}E_{𝐤1},& 0,& 0,& \mathrm{\Delta }\\ 0,& E_{𝐤1},& \mathrm{\Delta },& 0\\ 0,& \mathrm{\Delta }^{},& E_{𝐤2},& 0\\ \mathrm{\Delta }^{},& 0,& 0,& E_{𝐤2}\end{array}\right)\left(\begin{array}{c}f_{𝐤1}\\ f_{𝐤1}\\ f_{𝐤2}^{}\\ f_{𝐤2}^{}\end{array}\right)$$
(50)
with $`E_0=2_𝐤E_{𝐤2}+\mathrm{\Delta }^2/(2J)`$. Diagonalization of this $`4\times 4`$ matrix leads to the eigenvalues
$$\lambda _{1,2}=\frac{1}{2}\left(E_{𝐤1}E_{𝐤2}\right)\pm \sqrt{\left(\frac{E_{𝐤1}+E_{𝐤2}}{2}\right)^2+|\mathrm{\Delta }|^2}.$$
(51)
Both eigen modes are doubly degenerate and with an isotropic gap $`\mathrm{\Delta }`$. We take the form of usual dispersion relation for degenerate bands $`E_{𝐤1}=E_{𝐤2}`$.
The corresponding combinations of the wave function components are (for $`\mathrm{\Delta }=\mathrm{\Delta }^{}`$)
$$\{\begin{array}{c}(E_{k1}\lambda )(\psi _1+\psi _2)\mathrm{\Delta }(\psi _3\psi _4)=0\\ \mathrm{\Delta }(\psi _1+\psi _2)+(E_{𝐤2}+\lambda )(\psi _3\psi _4)=0,\end{array}$$
(52)
and
$$\{\begin{array}{c}(E_{k1}\lambda )(\psi _1\psi _2)+\mathrm{\Delta }(\psi _3+\psi _4)=0\\ \mathrm{\Delta }(\psi _1\psi _2)(E_{𝐤2}+\lambda )(\psi _3+\psi _4)=0.\end{array}$$
(53)
For the sake of simplicity we consider here only the case $`E_{𝐤1}=E_{𝐤2}=E_𝐤`$, as it provides the main character of the eigenstates. Moreover, it is sufficient to consider only the system (52) due to the double degeneracy of the eigenvalues. By standard method (including the wave function normalization) we obtain the quasiparticle operators
$$\alpha _𝐤=u_𝐤\frac{1}{\sqrt{2}}\left(f_{𝐤1}+f_{𝐤1}\right)+v_𝐤\frac{1}{\sqrt{2}}\left(f_{𝐤2}^{}f_{𝐤2}\right),$$
(54)
and
$$\beta _𝐤^{}=v_𝐤\frac{1}{\sqrt{2}}\left(f_{𝐤1}+f_{𝐤1}\right)+u_𝐤\frac{1}{\sqrt{2}}\left(f_{𝐤2}^{}f_{𝐤2}\right),$$
(55)
where
$$u_𝐤=\frac{1}{\sqrt{2}}\left(1+\frac{E_𝐤}{\sqrt{E_𝐤^2+\mathrm{\Delta }^2}}\right)^{1/2},v_𝐤=\frac{1}{\sqrt{2}}\left(1\frac{E_𝐤}{\sqrt{E_𝐤^2+\mathrm{\Delta }^2}}\right)^{1/2}.$$
(56)
Again, we have a combination of the two types of states. The corresponding wave equation which replaces the Bogolyubov-De Gennes equation for one-band singlet superconductor reads in the effective-mass approximation
$$i\mathrm{}_t\left(\begin{array}{c}\psi _1\\ \psi _2\\ \psi _3\\ \psi _4\end{array}\right)=\left(\frac{\mathrm{}^2}{2m^{}}^2\mu \right)\left(\begin{array}{c}\psi _1\\ \psi _2\\ \psi _3\\ \psi _4\end{array}\right)+\mathrm{\Delta }\left(\begin{array}{c}\psi _4\\ \psi _3\\ \psi _2\\ \psi _1\end{array}\right).$$
(57)
We see that the pairing couples explicitly the particle and hole components ( $`\psi _1`$ with $`\psi _4`$, $`\psi _2`$ with $`\psi _3`$, etc.). This equation forms a basis for the discussion of inhomogeneous paired states.
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# Energetics and geometry of excitations in random systems
## Abstract
Methods for studying droplets in models with quenched disorder are critically examined. Low energy excitations in two dimensional models are investigated by finding minimal energy interior excitations and by computing the effect of bulk perturbations. The numerical data support the assumptions of compact droplets and a single exponent for droplet energy scaling. Analytic calculations show how strong corrections to power laws can result when samples and droplets are averaged over. Such corrections can explain apparent discrepancies in several previous numerical results for spin glasses.
Magnets and superconductors are examples of physical systems where quenched disorder often plays a dominant role. Such systems can exhibit hysteresis effects and long relaxation times that are the manifestation of the large energy barriers created by the quenched disorder. One scenario that makes predictions for the equilibrium and nonequilibrium behavior of disorder dominated phases is the droplet or scaling picture . Predictions in this scenario follow from scaling assumptions for the energetic and geometric properties of excitations. For simple topological reasons, excitations can be defined as regions where the configuration is uniformly related by a symmetry to a global ground state (e.g., spin flipped domains.) In the droplet picture, the low lying excitations of size $`l`$ are connected and compact: they have volume $`l^{d_f}`$, with dimension $`d_f`$ equal to the system dimension $`d`$, and the surface to volume ratio decreases as $`l`$ increases. Droplet boundaries are fractal, with surface dimension $`d_s<d`$. The central ansatz is that the probability distribution $`\rho (\mathrm{\Delta },l)`$ for the energy $`\mathrm{\Delta }`$ of a droplet of size $`l`$ in a given volume $`l^d`$ has a characteristic scale $`l^\theta `$. This distribution is argued to have finite weight at $`\mathrm{\Delta }=0`$. The two exponents $`\theta `$ and $`d_s`$ can be used, for example, to predict many of the properties of a spin glass . This scenario is consistent with numeric results for excitations created by modifying boundary conditions . However, other work has suggested that there may be more than one important energy scaling exponent and more complicated geometries for excitations. The proposed distinct exponents separately describe (i) boundary induced domain wall excitations and (ii) excitations induced by internal constraints or external fields. It has also been suggested that there is distinct scaling for large droplets created by modifying the quenched disorder . It is important to understand these claims, as they suggest that the standard droplet picture is, at best, incomplete.
To provide perspective, it is useful to investigate in detail systems which lend themselves to precise study, where some analytic results are known and large systems can be simulated efficiently. Results are presented here for a 2D elastic medium and a 2D Ising spin glass. Single interior droplets, which include a specified central point, are computed for the elastic medium. In contrast with work on interior droplets in 2D spin glasses , a fast, exact algorithm is used, allowing for precise checks of scaling. The responses of the elastic medium and the spin glass to bulk perturbations are also calculated. The numerical results for droplet energies and geometrical characteristics show that logarithmic or small power law corrections are quite strong. These corrections can be understood in detail by arguments within the droplet picture. Droplets that are not induced by boundary conditions are only bounded above by the system size $`L`$ and below by a discretization scale, so that all scales between must be considered when computing averages. Corrections to scaling for droplets of fixed scale $`l`$, such as $`l^1`$ or $`L^1`$ corrections (e.g., from lattice discreteness) or unknown irrelevant operators, might also be considered. However, the scale averaging corrections are apparently dominant for some quantities. Such corrections lead to an effective energy exponent distinct from $`\theta `$, as boundary condition induced domain walls do not have such corrections. To remove scale averaging corrections, one can group the droplets by scale $`l`$ and study the geometry and energy as a function of $`l`$ (or $`l/L`$ if one is interested in large droplets ), as $`L\mathrm{}`$. With this analysis, the numerical results provide strong evidence that the droplets are “compact”, with fractal domain walls, and that there is a single energy exponent $`\theta `$.
One model that I study here is for a two dimensional elastic medium, with scalar displacement field $`u(x)`$, interacting with quenched periodic disorder. The continuum energy functional is $`[u(x)]=d^2x[u(x)]^2+V(u(x),x)`$, where $`V`$ has short range correlations in its second argument and is periodic in its first argument, $`V(u(x)+1,x)=V(u(x),x)`$. This model has been used for vortex lattices in superconductors, incommensurate charge density waves, and crystal growth on a disordered substrate . The continuum model can be discretized on a scale $`a`$, where the disorder and elastic energies balance . As an effective degree of freeedom $`i`$ is pinned to a preferred configuration (up to periodic shifts), the displacements $`u_i`$ are of the form $`n_i+\beta _i`$, for integer $`n_i`$ and fixed $`\left\{\beta _i\right\}`$. Elastic interactions tend to minimize nearest neighbor differences in $`u_i`$, with excitations of the medium being regions displaced relative to the ground state. Since the $`u_i`$ are discretized, domain walls separate regions relatively shifted by unit amounts. Numerical work for zero temperature ($`T=0`$) has determined properties of the ground state and the scaling of boundary induced domain wall energies . Finite temperature simulations , both Monte Carlo and combinatorial, have shown that the $`T=0`$ phase is stable at finite $`T`$. This model is thus a useful prototype for models with finite $`T`$ transitions, such as the 3D spin glass .
Another model treated here is the 2D Ising spin glass, with Hamiltonian $`=_{ij}J_{ij}s_is_j`$, with spins $`s_i=\pm 1`$ on a triangular lattice and Gaussian distributed $`J_{ij}`$. The ground states $`\left\{s_i^0\right\}`$ for samples in this model were found by a combinatorial method for a standard graph representation .
For the elastic medium, minimal energy domain walls about the center of a sample were studied on a square lattice using a polynomial time algorithm that calculated the energy $`\mathrm{\Delta }_o`$ and the droplet boundary. One method to characterize the compactness of droplets is to compare $`R_O`$, the radius of the smallest circle that encloses the droplet, with $`R_I`$, the radius of the largest circle contained by the boundary vertices. Droplets can be studied in spin glasses by finding the ground state and then recomputing the ground state $`\left\{s_i^ϵ\right\}`$ with modified couplings $`J_{ij}J_{ij}ϵL^ds_i^0s_j^0`$ . This bulk perturbation can introduce excitations on all scales. Sample excitations are depicted in Fig. 1.
The droplet energies are of great interest, as these are believed to determine the static correlation functions at finite temperature and the relaxation to equilibrium. Consider the problem of finding the minimal energy droplet around the origin in a system of size $`L`$ . Assume that there is a factor $`b`$ which gives a separation of length scales: droplets differing in size by $`b`$ are independent. One independent droplet excitation could be excited at each scale, so that sums over all scales must be performed (for similar sums, see Refs. .) At each length scale $`s`$, $`s=1,2,\mathrm{},\mathrm{log}_b(L)`$, the distribution for the energy $`\mathrm{\Delta }(b^s)`$ has characteristic scale $`b^{s\theta }`$ and has finite weight $`b^{s\theta }`$ at $`\mathrm{\Delta }=0`$ . The total density of states $`\rho (\mathrm{\Delta }_o)`$ for $`\mathrm{\Delta }_o<L^\theta `$ is then a sum over $`s`$, giving $`\rho (\mathrm{\Delta }_o)L^\theta \left|1(b^{}L)^\theta \right|`$, for $`\theta 0`$, where $`b^{}>0`$ is set by $`b`$ and the lattice and boundary conditions, which affect the lower and upper ends of the sum. The expected minimum value for $`\mathrm{\Delta }_o`$ scales as $`L^\theta \left|1(b^{}L)^\theta \right|^1`$, the subdominant term reflecting that the minimal energy droplet is chosen from all length scales from $`1`$ to $`L`$. The effective energy exponent is then (for $`\theta <0`$)
$$\theta ^{\mathrm{eff}}=\frac{d\mathrm{ln}(\overline{\mathrm{\Delta }_o})}{d\mathrm{ln}(L)}=\frac{\theta }{1(b^{}L)^\theta }.$$
(1)
Applying Eqn. (1) to the 2D spin glass, taking $`b^{}=2`$ and $`L=16`$, gives an effective exponent $`\theta ^{\mathrm{eff}}=0.45`$, apparently quite different from the domain wall value $`\theta =0.28`$ and consistent with the alternate energy exponent proposed in earlier numerical work . The effective exponent converges to $`\theta `$ quite slowly with $`L`$ (and is relatively insensitive to $`b^{}`$), as $`\theta `$ is near zero.
One case where $`\theta =0`$ for domain walls created by boundary conditions is the 2D elastic medium. Large domain walls can be created by external strains. By statistical tilt invariance of the disorder , the change in the sample averaged energy can be found by computing the elastic energy only, as the change in the sample averaged pinning energy is zero. Displacing one end of a sample by $`\delta u=1`$ to induce one domain wall gives an elastic energy density $`L^2`$ over the volume $`L^2`$, so that the total domain wall energy scales as a constant ($`\theta =0`$.) This result is consistent with previous numerical simulations of boundary induced domain walls . However, the mean interior droplet energy $`\overline{\mathrm{\Delta }_o}(L)`$ can be fit over a decade with $`0.15<\theta <0.23`$, to within a few percent for smaller $`L`$. Arguments similar to those for $`\theta 0`$ can be applied to explain this. There are $`\mathrm{ln}(L/a)/\mathrm{ln}(b)`$ independent scales to choose from, each with identical droplet energy distributions ($`\theta =0`$). The inverse of the minimal droplet energy $`[\overline{\mathrm{\Delta }_o}(L)]^1`$ is therefore linear in $`\mathrm{ln}(L)`$. This result can also be derived using elasticity theory. The displacement at the origin of a region of size $`a`$ costs an elastic energy that scales as $`[\mathrm{ln}(L/a)]^1`$. By tilt symmetry, the pinning can be averaged over, so that $`\overline{\mathrm{\Delta }_o}(L)[\mathrm{ln}(L/a)]^1`$ for interior droplets constrained to contain the origin. The numerical results are quite consistent with these expectations, as shown by the two parameter fit displayed in Fig. 2 (in addition, the computed probability of generating a droplet of size $`l`$ is consistent with a distribution uniform in $`\mathrm{ln}(l)`$ .)
Similar corrections are important for the magnetization $`m(h)`$ of a spin glass in response to an external field $`h`$ . For the case $`\theta 0`$ and $`h<O(L^{\theta d/2})`$, $`m`$ is found by summing over scales the product of the probability $`hl^{d/2\theta }`$ of generating a droplet, its expected contribution $`l^{d/2}L^d`$ to the magnetization, and the number of droplets $`(L/l)^d`$ at scale $`l`$. This gives $`mh\left|(b_h^{}L)^\theta 1\right|`$, with $`b_h^{}>0`$ a constant, to be compared with the uncorrected singular piece $`mhL^\theta `$ (for numerics, see Ref. ). The size of the corrections are quite similar to those for $`\theta `$.
The measurement of geometrical quantities, such as boundary length and droplet area, can also be strongly influenced by scale averaging corrections if one averages over all length scales from $`1`$ to $`L`$. When the sample averaged area $`\overline{A(L)}`$ of interior droplets in the 2D elastic medium is computed as a function of $`L`$ and the local exponent $`d_{\overline{f}}^{\mathrm{eff}}=d[\mathrm{ln}(\overline{A(L)})]/d[\mathrm{ln}(L)]`$ is computed, the local dimension is less than two, which might suggest fractal droplets. This local exponent slowly changes with $`L`$, though (Fig. 3.) The local exponent for $`\overline{A(L)}`$ is $`d_{\overline{f}}^{\mathrm{eff}}=2[\mathrm{ln}(L/a^{})]^1+O(L^2)`$, where $`a^{}`$ depends on the boundary and lattice cutoffs. A useful procedure to reduce the corrections to $`d_f^{\mathrm{eff}}`$ is to separate out the scales and plot the droplet area $`A(R_O,L)`$ as a (binned) function of $`R_0`$. Changing the order of the averages gives local exponents that are much better fit by a constant. For $`L>32`$ and $`R_0>8`$, $`\overline{A(R_o)}R_o^{d_f}`$, with bulk droplet dimension $`d_f=2.01(2)`$ (Fig. .) A similar plot for the perimeter confirms that $`d_s=1.25(1)`$, with the surface to volume ratio vanishing as $`l^{dd_s}`$ for large droplets. Droplet compactness can also be confirmed by plotting the ratio $`k=R_O/R_I`$ binned according to $`R_O`$; it is found that the distribution of $`k`$ converges, with mean $`k(R_O,L)=2.85(5)`$, when $`L/2>R_O>32`$.
Corrections to geometric and energetic quantities are also important when computing link overlaps, such as those found in comparing the unperturbed $`J_{ij}`$ ground state, $`\left\{s_i^0\right\}`$, with the $`ϵ`$-perturbed state. The link overlap $`q_l`$ is the fraction of link values $`s_is_j`$ on bonds $`ij`$ which are unchanged. By summing the contributions over all scales (note the small droplets in Fig. 1(a)), it can be shown that the local exponent for the fraction $`1q_l`$ of changed bonds behaves as
$$\mu _l^{\mathrm{eff}}=\frac{d[\mathrm{ln}\overline{(1q_l)}]}{d\mathrm{ln}(L)}=\mu _l\frac{c^{}}{(cL)^{d\mu _l}1},$$
(2)
where $`\mu _l=\theta +2(dd_s)`$ and $`c,c^{}`$ are constants characterizing the upper and lower cutoffs. The computed local exponent for the 2D spin glass is shown in Fig. 4 ($`\mu _l^{\mathrm{eff}}`$ is relatively insensitive to $`ϵ`$, at least for $`0.35<ϵ<1.3`$.) Only for $`L>100`$ does $`\mu _l^{\mathrm{eff}}`$ approach the large $`L`$ limit of $`\mu _l1.18`$ (using the values $`\theta =0.28`$ and $`d_s=1.27`$.) Similar results are found for the 2D elastic medium . The exponent $`\mu `$ for spin overlaps $`q=L^ds_i^os_i^ϵ`$, with $`1qL^\mu `$, has smaller corrections of this form and may well be dominated by corrections to scaling from unknown operators or inverse lengths. Numerics show that $`\mu `$ converges much more quickly than $`\mu _l`$ in the 2D spin glass. The measurement of $`\mu _l`$ (relative to $`\mu `$) has been used by Palassini and Young to conclude that a second energy exponent $`\theta ^{}`$ affects the response to bulk perturbations in the 3D spin glass. In three dimensions, the scale averaging corrections decrease more quickly with $`L`$, as $`d\mu _l1.3`$ compared with the 2D correction exponent $`d\mu _l0.82`$, but the system sizes that can be simulated are much smaller. It may be that corrections due to irrelevant variables or possible $`1/L`$ effects are dominant, but scale averaging corrections clearly contribute to errors in $`\mu _l`$. A correction of $`\delta \mu _l0.2`$ for $`L8`$, from similar $`c`$ and $`c^{}`$, would invalidate the conclusions of Ref. .
In summary, analysis of numerical data provides a precise confirmation of the droplet picture in the bulk of a sample, both for the scaling of the energies and for the geometrical structure of droplets. In comparing the numerical results with the droplet picture, care must be taken to understand where strong corrections might arise. The corrections arising from averaging over multiple scales can be predicted in some detail and apply to spin glasses and other models where $`\theta `$ is near zero. The corrections to averaged quantities such as the droplet energy $`\overline{\mathrm{\Delta }_o}(L)`$ are satisfactorily explained, for $`L^2>10^2`$ systems, by averaging simple power laws over scales between $`1`$ and $`L`$. The link overlap, which inherently averages over scales, is strongly affected by finite size corrections for small $`\theta `$ and $`dd_s`$, with effective exponent corrections greater than $`0.1`$ for $`L<30`$ in the 2D spin glass. It has been suggested that, for topological reasons, distinct $`\theta `$ exponents exist only in $`d3`$ . The numerics in this Letter are for $`d=2`$, but they and the general analysis suggest that large finite size effects strongly affect results in three dimensions. To reduce these types of corrections, data can be binned over droplet size $`l`$ (or over $`l/L`$) at fixed system size $`L`$, checking for convergence by then increasing $`L`$.
I would like to thank David Huse for bringing to my attention some recent work on droplets, Olivier Martin for a stimulating discussion, and Daniel Fisher for discussions. The LEDA library was most useful in the study of the geometry of the droplets. This work was supported by the National Science Foundation (DMR-9702242).
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# Infinite Dimensional Quantum Information Geometry
## Introduction
In finite dimensional quantum information geometry, the set upon which the geometric structures are defined is simply the set of all (invertible) density matrices on a finite dimensional Hilbert space grasselli:PetzSudar96 . Already in the definition of the underlying set in infinite dimensions, we need to be slightly more careful and take a more restrictive set than just that of all (invertible) density operators. The reason for this is that, if we modify a given state in the set with a small perturbation, we want the perturbed state to have the same properties as the original one.
Let $`𝒞_p,0<p<1`$, denote the set of compact operators $`A:`$ such that $`|A|^p𝒞_1`$, where $`𝒞_1`$ is the set of trace-class operators on $``$. Define $`𝒞_{<1}:=_{0<p<1}𝒞_p`$. We take the underlying set of the quantum information manifold to be $`=𝒞_{<1}\mathrm{\Sigma }`$ where $`\mathrm{\Sigma }𝒞_1`$ denotes the set of density operators. This guarantees that, if $`\rho _0`$, there exists a $`\beta _0<1`$ such that $`\rho _0^{\beta _0}`$ is a density as well as $`\rho _0`$ itself. This set has an affine structure induced from the linear structure of each $`𝒞_p`$ in the following way: let $`\rho _1𝒞_{p_1}\mathrm{\Sigma }`$ and $`\rho _2𝒞_{p_2}\mathrm{\Sigma }`$; take $`p=\mathrm{max}\{p_1,p_2\}`$, then $`\rho _1,\rho _2𝒞_p\mathrm{\Sigma }`$, since $`pq`$ implies $`𝒞_p𝒞_q`$ grasselli:Pietsch72 ; define “$`\lambda \rho _1+(1\lambda )\rho _2,0\lambda 1`$” as the usual sum of operators in $`𝒞_p`$. This is called the $`(1)`$-affine structure. However, this is not the affine structure we use to define a flat connection on the manifold. Instead, we equip $``$ with an exponential affine structure and then define the natural connection associated with it.
To each $`\rho _0𝒞_{\beta _0}\mathrm{\Sigma }`$, $`\beta _0<1`$, let $`H_0=\mathrm{log}\rho _0+cII`$ be a self-adjoint operator with domain $`𝒟(H_0)`$ such that $`\rho _0=Z_0^1e^{H_0}=e^{(H_0+\mathrm{\Psi }_0)}`$.
The idea is to perturb the Hamiltonian $`H_0`$, obtaining a new Hamiltonian $`H_X`$ and then construct a neighbourhood of the point $`\rho _0`$ consisting of the perturbed states $`\rho _X`$. The perturbations considered are of three different types, and that is the content of the next section.
## Perturbations
The most general class of perturbations we use are the form bounded perturbations. Given a positive self-adjoint operator $`H`$ with associated form $`q_H`$ and form domain $`Q(H)`$, we say that a symmetric quadratic form $`X`$ (or the symmetric sesquiform obtained from it by polarization) is $`q_H`$-bounded if
1. $`Q(H)Q(X)`$ and
2. there exist positive constants $`a`$ and $`b`$ such that $`\left|X(\psi ,\psi )\right|aq_H(\psi ,\psi )+b(\psi ,\psi )`$, for all $`\psi Q(H)`$.
Although form bounded perturbations are of much interest in the study of Schrödinger operators for a variety of quantum systems, they provide very little regularity for the quantum information manifold. In grasselli:Streater98 , Streater was able to show that the manifold constructed using then has Lipschitz structure. However, more regularity is needed if we want to define a metric on it by, say, the second derivative of the free energy. This led to the idea of looking at the more restrictive case of operator bounded perturbations. Given operators $`H`$ and $`X`$ defined on dense domains $`𝒟(H)`$ and $`𝒟(X)`$ in a Hilbert space $``$, we say that $`X`$ is $`H`$-bounded if
1. $`𝒟(H)𝒟(X)`$ and
2. there exist positive constants $`a`$ and $`b`$ such that $`X\psi aH\psi +b\psi `$, for all $`\psi 𝒟(H)`$.
For both form bounded and operator bounded perturbations, the infimum of such $`a`$ is called the relative bound of $`X`$ (with respect to $`H`$ or with respect to $`q_H`$, accordingly). If $`a<1`$, the perturbation is said to be small.
Operator bounded perturbations are also used in the study of Schrödinger operators for quite a broad range of quantum systems grasselli:ReedSimon75 , with the additional property of providing enough regularity for the manifold $``$ to have an analytic free energy, for instance.
The following lemma tells us how to characterise operator bounded and form bounded perturbations in terms of certain norms. Before we state it, we need to say what we mean by a form multiplied from both sides by operators. Suppose that $`X`$ is a quadratic form with domain $`Q(X)`$ and $`A,B`$ are operators on $``$ such that $`A^{}`$ and $`B`$ are densely defined. Suppose further that $`A^{}:𝒟(A^{})Q(X)`$ and $`B:𝒟(B)Q(X)`$. Then the expression $`AXB`$ means the form defined by
$$\varphi ,\psi X(A^{}\varphi ,B\psi ),\varphi 𝒟(A^{}),\psi 𝒟(B).$$
Consider now the case where $`H_0I`$ is a self-adjoint operator with domain $`𝒟(H_0)`$, quadratic form $`q_0`$ and form domain $`Q_0=𝒟(H_0^{1/2})`$, and let $`R_0=H_0^1`$ be its resolvent at the origin.
###### Lemma 1
A symmetric operator $`X:𝒟(H_0)`$ is $`H_0`$-bounded if and only if $`XR_0<\mathrm{}`$. Analogously, a symmetric quadratic form $`X`$ defined on $`Q_0`$ is $`q_0`$-bounded if and only if $`R_0^{1/2}XR_0^{1/2}`$ is a bounded symmetric form defined everywhere. Moreover, if $`R_0^{1/2}XR_0^{1/2}<\mathrm{}`$ then the relative bound $`a`$ of $`X`$ with respect to $`q_0`$ satisfies $`aR_0^{1/2}XR_0^{1/2}`$.
The set $`𝒯_\omega (0)`$ of all $`H_0`$-bounded symmetric operators X is a Banach space with norm $`X_\omega (0):=XR_0`$, since the map $`AAH_0`$ from $`()`$ onto $`𝒯_\omega (0)`$ is an isometry. The set $`𝒯_0(0)`$ of all $`q_0`$-bounded symmetric forms $`X`$ is also a Banach space with norm $`X_0(0):=R_0^{1/2}XR_0^{1/2}`$, since the map $`AH_0^{1/2}AH_0^{1/2}`$ from the set of all bounded self-adjoint operators on $``$ onto $`𝒯_0(0)`$ is again an isometry.
Motivated by Banach space interpolation theory, let us consider, for $`\epsilon (0,1/2)`$, the set $`𝒯_\epsilon (0)`$ of all symmetric forms X with $`𝒟(H_0^{\frac{1}{2}\epsilon })Q(X)`$ and such that $`X_\epsilon (0):=R_0^{\frac{1}{2}+\epsilon }XR_0^{\frac{1}{2}\epsilon }`$ is finite. Then the map $`AH_0^{\frac{1}{2}\epsilon }AH_0^{\frac{1}{2}+\epsilon }`$ is an isometry from the set of all bounded self-adjoint operators on $``$ onto $`𝒯_\epsilon (0)`$. Hence $`𝒯_\epsilon (0)`$ is a Banach space with the $`\epsilon `$-norm $`_\epsilon (0)`$. Such an $`X`$ will be called a $`\epsilon `$-bounded perturbation and it is shown to interpolate between the extreme cases of form bounded perturbations, for which $`\epsilon =0`$, and operator bounded perturbations for which $`\epsilon =1/2`$.
###### Lemma 2
For fixed symmetric $`X`$, $`X_\epsilon `$ is a monotonically increasing function of $`\epsilon [0,1/2]`$.
In the next section, we carry out the programme of using these perturbations to obtain the hoods in the manifold. In what follows, we write $`𝒯_{()}(0)`$ to indicate that we can use any of the three Banach spaces $`𝒯_\omega (0)`$, $`𝒯_0(0)`$ or $`𝒯_\epsilon (0)`$.
## Construction of the Manifold
The two technical tools used in the construction of our manifold are the following.
###### Theorem 3 (KLMN)
Let $`H_0`$ be a positive self-adjoint operator with quadratic form $`q_0`$ and form domain $`Q_0`$; let $`X`$ be a $`q_0`$-small symmetric quadratic form. Then there exists a unique self-adjoint operator $`H_X`$ with form domain $`Q_0`$ such that
$$H_X^{1/2}\varphi ,H_X^{1/2}\psi =q_0(\varphi ,\psi )+X(\varphi ,\psi ),\varphi ,\psi Q_0.$$
Moreover, $`H_X`$ is bounded below by $`b`$.
###### Lemma 4 (Streater 98)
Let $`X`$ be a $`q_0`$-small with bound $`a<1\beta _0`$. Denote by $`H_X`$ the unique operator given by the KLMN theorem. Then $`\mathrm{exp}(\beta H_X)`$ is of trace class for all $`\beta >\beta _X=\beta _0/(1a)`$.
The construction of the neighbourhood of $`\rho _0`$ goes as follows. In $`𝒯_{()}(0)`$, take $`X`$ such that $`X_((0)<1\beta _0`$. Since $`X_0(0)X_{()}(0)<1\beta _0`$, $`X`$ is also $`q_0`$-bounded with bound $`a_0`$ less than $`1\beta _0`$. The KLMN theorem then tells us that there exists a unique semi-bounded self-adjoint operator $`H_X`$ with form $`q_X=q_0+X`$ and form domain $`Q_X=Q_0`$. We write $`H_X=H_0+X`$ for this operator and consider the state $`\rho _X=Z_X^1e^{(H_0+X)}=Z_X^1e^{(H_0+X+\mathrm{\Psi }_X)}`$
Then, from lemma 4, $`\rho _X𝒞_{\beta _X}\mathrm{\Sigma }`$, where $`\beta _X=\frac{\beta _0}{1a_0}<1`$. If we add to $`H_X`$ a multiple of the identity, we can still have the same state $`\rho _X`$, by simply adjusting the partition function $`Z_X`$; so we can always assume that, for the perturbed state, we have $`H_XI`$. We take as a hood $`_0`$ of $`\rho _0`$ the set of all such states, that is, $`_0=\{\rho _X:X_{()}(0)<1\beta _0\}`$.
To give a topology to $`_0`$, we first introduce in $`𝒯_{()}(0)`$ the equivalence relation $`XY`$ iff $`XY=\alpha I`$ for some $`\alpha 𝐑`$, precisely because $`\rho _X=\rho _{X+\alpha I}`$, as remarked above. We then identify $`\rho _X`$ in $`_0`$ with the line $`\{Y𝒯_{()}(0):Y=X+\alpha I,\alpha 𝐑\}`$ in $`𝒯_{()}(0)/`$. This is a bijection from $`_0`$ onto the subset of $`𝒯_{()}(0)/`$ defined by $`\{\{X+\alpha I\}_{\alpha 𝐑}:X_{()}(0)<1\beta _0\}`$. The topology in $`_0`$ is then given by transfer of structure. Now that $`_0`$ is a (Hausdorff) topological space, we want to parametrise it by an open set in a Banach space. As in the classical case grasselli:PistoneSempi95 , we choose the Banach subspace of centred variables in $`𝒯_{()}(0)`$; in our terms, perturbations with zero mean (the ‘scores’). The only problem is that, when $`X`$ is not an operator, it is not immediately clear what its mean in the state $`\rho _0`$ should be, let alone the question of whether or not it is finite. To deal with this, define the regularised mean of $`X𝒯_{()}(0)`$ in the state $`\rho _0`$ as $`\rho _0X:=\text{Tr}(\rho _0^\lambda X\rho _0^{1\lambda })`$, for $`0<\lambda <1`$.
Since $`\rho _0𝒞_{\beta _0}\mathrm{\Sigma }`$ and $`X`$ is $`q_0`$-bounded, lemma 5 of grasselli:Streater98 ensures that $`\rho _0X`$ is finite and independent of $`\lambda `$. It was a shown there that $`\rho _0X`$ is a continuous map from $`𝒯_0(0)`$ to $`𝐑`$, because its bound contained a factor $`X_0(0)`$. Exactly the same proof shows that $`\rho _0X`$ is a continuous map from $`𝒯_{()}(0)`$ to $`𝐑`$. Thus the set $`\widehat{𝒯}_{()}(0):=\{X𝒯_{()}(0):\rho _0X=0\}`$ is a closed subspace of $`𝒯_{()}(0)`$ and so is a Banach space with the norm $`_\epsilon `$ restricted to it. We notice that for the case of operator bounded perturbations, the regularised mean of $`X`$ coincides with the usual mean $`\text{Tr}(\rho _0X)`$.
To each $`\rho _X_0`$, consider the point in the line $`\{X+\alpha I\}_{\alpha 𝐑}`$ with $`\alpha =\rho _0X`$. Write $`\widehat{X}=X\rho _0X`$ for this point. The map $`\rho _X\widehat{X}`$ is a homeomorphism between $`_0`$ and the open subset of $`\widehat{𝒯}_{()}(0)`$ defined by $`\{\widehat{X}:\widehat{X}=X\rho _0X,X_{()}<1\beta _0\}`$. The map $`\rho _X\widehat{X}`$ is then a global chart for the Banach manifold $`_0`$ modeled by $`\widehat{𝒯}_{()}(0)`$. The tangent space at $`\rho _0`$ is given by $`\widehat{𝒯}_{()}(0)`$, with the curve $`\left\{\rho (\lambda )=Z_{\lambda X}^1e^{(H_0+\lambda X)},\lambda [\delta ,\delta ]\right\}`$ having tangent vector $`\widehat{X}=X\rho _0X`$.
We extend our manifold by adding new patches compatible with $`_0`$. The idea is to construct a chart around each perturbed state $`\rho _X`$ as we did around $`\rho _0`$. Let $`\rho _X_0`$ with Hamiltonian $`H_XI`$ and consider the Banach space $`𝒯_{()}(X)`$ of all symmetric forms $`Y`$ such that the norm $`Y_{()}(X)`$ is finite, where the expression for this norm is the same as $`Y_{()}(0)`$ but with all the Hamiltonians replaced by $`H_X`$. Then we repeat exactly the same process, namely, take sufficiently small $`Y`$ (with $`Y_{()}(X)<1\beta _X`$), obtain from the KLMN theorem the Hamiltonian $`H_{X=Y}`$ with form $`q_{X+Y}=q_X+Y=q_0+X+Y`$ and form domain $`Q_{X+Y}=Q_X=Q_0`$ and take the hood of $`X`$ to be the set $`_X`$ of all states of the form $`\rho _{X+Y}=Z_{X+Y}^1e^{H_{X+Y}}=Z_{X+Y}^1e^{(H_0+X+Y)}`$. The topology and coordinates for $`_X`$ are then introduced in a completely similar fashion.
We then turn to the union of $`_0`$ and $`_X`$. We need to show that our two previous charts are compatible in the overlapping region $`_0_X`$. For the case of form bounded and operator bounded perturbations, the equivalence of the norms is achieved by a straightforward series of operator identities grasselli:Streater98 ; grasselli:Streater99 . For the case of $`\epsilon `$-bounded perturbations, the argument is more subtle and involves careful consideration of the domains of the operators. Nevertheless, the result still follows grasselli:GrasselliStreater00 .
###### Theorem 5
$`_\epsilon (X)`$ and $`_\epsilon (0)`$ are equivalent norms.
We can repeat the construction again, starting from any point in $`_0_X`$. We then obtain the following definition.
###### Definition 6
The information manifold $`(H_0)`$ defined by $`H_0`$ consists of all states obtainable in a finite numbers of steps, by extending $`_0`$ as explained above.
## Affine Geometry in $`(H_0)`$
The set $`A=\{\widehat{X}\widehat{𝒯}_{()}(0):\widehat{X}=X\rho _0X,X_{()}(0)<1\beta _0\}`$ is a convex subset of the Banach space $`\widehat{𝒯}_{()}(0)`$ and so has an affine structure coming from its linear structure. We provide $`_0`$ with an affine structure induced from $`A`$ using the patch $`\widehat{X}\rho _X`$ and call this the canonical or $`(+1)`$-affine structure. The $`(+1)`$-convex mixture of $`\rho _X`$ and $`\rho _Y`$ in $`_0`$ is then $`\rho _{\lambda X+(1\lambda )Y}`$, $`(0\lambda 1)`$ , which differs from the previously defined $`(1)`$-convex mixture $`\lambda \rho _X+(1\lambda )\rho _Y`$.
Given two points $`\rho _X`$ and $`\rho _Y`$ in $`_0`$ and their tangent spaces $`\widehat{𝒯}_\epsilon (X)`$ and $`\widehat{𝒯}_\epsilon (Y)`$, we define the $`(+1)`$-parallel transport $`U_L`$ of $`(Z\rho _XZ)\widehat{𝒯}_\epsilon (X)`$ along any continuous path $`L`$ connecting $`\rho _X`$ and $`\rho _Y`$ in the manifold to be the point $`(Z\rho _YZ)\widehat{𝒯}_\epsilon (Y)`$. Clearly $`U_L`$ is independent of $`L`$ by construction, thus the $`(+1)`$-affine connection is flat. We see that the $`(+1)`$-parallel transport just moves the representative point in the line $`\{Z+\alpha I\}_{\alpha 𝐑}`$ from one hyperplane to another.
## Discussion
The manifold $`(H_0)`$ constructed here does not cover the whole set $``$ at once. It could not possibly do so, since our small perturbations do not change the domain of the original Hamiltonian $`H_0`$, and certainly $``$ contains states defined by Hamiltonians with plenty of different domains. Also, although we can reach far removed points with a finite number of small perturbations, we can not move in arbitrary directions. For instance, we cannot reach $`X=H_0`$ as the result of our perturbations, since the identity is not an operator of trace class. To cover $``$ entirely, we have to start at several different points. The whole manifold thus obtained consists of several disconnected parts, pointing towards positive directions with respect to the given Hamiltonians.
Finally, it is clear that $`\lambda \rho _X+(1\lambda )\rho _Y`$, for, say, $`\rho _X,\rho _Y_0`$, defines a new state in the underlying set $``$. Nonetheless, we were not yet able to prove that it belongs to the neighbourhood of the original states. This is the main obstacle to define a mixture connection in our manifold and thence develop a quantum version of Amari’s duality theory grasselli:Amari85 . One possibility is to change the definition of the neighbourhoods altogether and use, for instance, the condition that states in the same neighbourhood all have finite relative entropy with respect to each other, besides having finite von Neumann entropy.
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# Measuring sparticle masses in non-universal string inspired models at the LHC
## 1 Introduction
The purpose of this work is to extend the discussion of LHC supersymmetry (SUSY) searches to include string models. We begin by discussing whether string models can be used to motivate previous work on LHC SUSY searches, and then suggest a well-motivated non-universal string model for a new pilot study. We go on to examine how the SUSY particles can be detected and how the model can be distinguished from a similar well studied supergravity (SUGRA) model. We reconstruct sparticle masses by looking for kinematic edges in $`\stackrel{~}{q}_L\stackrel{~}{\chi }_2^0q\stackrel{~}{l}_R^\pm l^{}q\stackrel{~}{\chi }_1^0l^\pm l^{}q`$ and $`\stackrel{~}{q}_L\stackrel{~}{\chi }_2^0q\stackrel{~}{\chi }_1^0Xq\stackrel{~}{\chi }_1^0l^\pm l^{}q`$ decay chains, and in doing so generalize the method of by unifying the treatment of light and heavy sleptons. Additionally, with a novel method based on , we further constrain the $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{l}_R`$ masses by studying the kinematics of events containing pair produced sleptons: $`(gg/q\overline{q})\stackrel{~}{l}_R^+\stackrel{~}{l}_R^{}l^+\stackrel{~}{\chi }_1^0l^{}\stackrel{~}{\chi }_1^0`$. In particular, this allows the mass difference between the $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{l}_R`$ to be determined with sufficient accuracy to permit discrimination between the string model and the most similar standard SUGRA model. We suggest that our analysis is likely to be applicable, not just to string motivated non-universal models, but to other non-universal models as well.
## 2 Theory
### 2.1 Introduction
Throughout this work we assume that the effective theory describing TeV scale physics is the R-parity conserving minimal supersymmetric standard model (MSSM). Within this framework, the collider phenomenology is strongly affected by the SUSY breaking terms
$$=\frac{1}{2}\underset{a=1}{\overset{3}{}}M_a\lambda _a\overline{\lambda }_a\underset{i}{}m_i^2|\varphi _i|^2(A_{ijk}W_{ijk}+B\mu H_1H_2+\text{H.c.}),$$
(1)
where $`i,j,k=Q_L`$, $`u_R^c`$, $`d_R^c`$, $`L_L`$, $`H_1`$, $`H_2`$ and $`\varphi _a,\lambda _a`$ are the scalar and gaugino fields of the MSSM (see e.g. Ref. ). $`W_{ijk}`$ are the trilinear pieces of the MSSM superpotential, which written in terms of superfields is
$$W=𝐘_ELH_1\overline{E}+𝐘_DQH_1^b\overline{D}+𝐘_UQH_2\overline{U}+\mu H_1H_2,$$
(2)
where we have suppressed all gauge and family indices. $`𝐘_D`$, $`𝐘_E`$, $`𝐘_U`$ denote the down quark, charged lepton and up quark Yukawa matrices respectively.
In Eq. (1), a general parameterisation of possible SUSY breaking effects has been employed<sup>1</sup><sup>1</sup>1It has been assumed that non-standard terms such as those discussed in Ref. are disallowed because they cause a naturalness problem in the presence of gauge singlets.. Usually, the parameters are constrained by the condition of universality,
$$m_i=m_0,A_{ijk}=A_0,M_a=M_{1/2}$$
(3)
deriving from simple SUGRA models. Eq. (3) is subject to radiative corrections and should be imposed at the string scale $`M_S`$. In the usual formulation of perturbative string theory, this corresponds to $`M_S5\times 10^{17}`$ GeV, but Eq. (3) is usually applied at the grand unified scale $`M_{GUT}2\times 10^{16}`$ GeV as an approximation. Once $`B`$ and $`\mu `$ are constrained by radiative electroweak symmetry breaking , the SUSY breaking sector is then characterised by one sign: sgn$`\mu `$, and four scalar parameters: $`m_0,A_0,M_{1/2},`$ and $`\mathrm{tan}\beta `$, the ratio $`v_2/v_1`$ of the two MSSM Higgs vacuum expectation values (VEVs). Once these are specified and current data are used to predict supersymmetric couplings such as the top Yukawa coupling and gauge couplings, the sparticle spectrum and decay chains are specified. Five points (denoted S1-S5) in the SUGRA parameter space $`m_0,A_0,M_{1/2},`$ sgn$`\mu `$, $`\mathrm{tan}\beta `$ have been suggested for study of SUSY production at the LHC and are catalogued in Table 1. These models have been well studied in the context of the LHC , and we shall use them as a reference to compare and contrast with new models, which do not necessarily obey Eq. (3).
### 2.2 SUGRA point compatibility with string models
In Refs. , the authors study the phenomenological viability of string and M-theory scenarios coming from the desirable absence of dangerous charge and colour breaking (CCB) minima or unbounded from below (UFB) directions in the effective potential. One of the models considered in is weakly coupled string theory with orbifold compactifications. In this case, the soft masses evaluated at the string scale are dependent upon the modular weights $`n_i`$ of the $`\varphi _i`$ fields and do not necessarily conform with Eq. (3. Other string scenarios exist which erase the UFB/CCB global minima which we do not explicitly investigate. For all modular weights equal to –1 however, the tree-level pattern of soft-masses conforms with Eq. (3), with the additional constraint
$$M_{1/2}=A_0=\sqrt{3}m_0.$$
(4)
In Table 1, we display under the “Weak” column whether each standard SUGRA point is approximately compatible with this sub-class of universal string models. None of S1-S5 fit exactly with this scenario, but S5 is the closest and would have a similar sparticle spectrum if $`m_0`$ were 173 GeV instead of 100 GeV. In regions of consistent radiative electroweak symmetry breaking (REWSB), the string model version of S5 has dangerous UFB minima . We reject this class of model for further study, partly because S5 gives a similar spectrum, and partly because it is already well studied, but mainly because of the UFB problem in the potential mentioned above.
The “M-theory” column of Table 1 refers to compatibility with the strong-coupling limit of $`E_8\times E_8`$ Heterotic string theory . Given simple assumptions (that SUSY is spontaneously broken by the auxiliary components of the bulk moduli superfields), the soft terms of M-theory valid at $`M_{GUT}`$ are
$`M_{1/2}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}M_{3/2}}{1+ϵ}}\left(\mathrm{sin}\theta +{\displaystyle \frac{ϵ}{\sqrt{3}}}\mathrm{cos}\theta \right)`$ (5)
$`m_0^2`$ $`=`$ $`M_{3/2}^2\left[1{\displaystyle \frac{3}{(3+ϵ)^2}}\left(ϵ(6+ϵ)\mathrm{sin}^2\theta +(3+2ϵ)\mathrm{cos}^2\theta 2\sqrt{3}ϵ\mathrm{sin}\theta \mathrm{cos}\theta \right)\right]`$ (6)
$`A_0`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}M_{3/2}}{3+ϵ}}\left[(32ϵ)\mathrm{sin}\theta +\sqrt{3}ϵ\mathrm{cos}\theta \right].`$ (7)
So, the SUGRA parameters $`M_{1/2}`$, $`m_0`$, $`A_0`$ become replaced by the goldstino angle $`\mathrm{sin}\theta `$, the gravitino mass $`M_{3/2}`$ and a ratio of moduli VEVs $`0<ϵ1`$. $`ϵ=0`$ corresponds to the weakly-coupled perturbative string. For S1-S4, we notice from Table 1 that $`A_0=0`$, allowing us to solve Eq. (7):
$$\mathrm{tan}\theta =\frac{\sqrt{3}ϵ}{2ϵ3}.$$
(8)
Specifying $`ϵ`$ then determines the ratio $`|m_0/M_{1/2}|^2`$, which is displayed in Fig. 1a. $`|m_0/M_{1/2}|^2`$ is always greater than one except as $`ϵ\mathrm{}`$, ruling out M-theory derivations of S1 and S2. S3 and S4 are compatible with $`ϵ=0.21`$ and $`0.53`$ respectively, as indicated in the figure.
Eq. (8) is not relevant for S5 because $`A_00`$, but the condition $`M_{1/2}=A_0`$ may be solved yielding
$$\mathrm{tan}\theta =\frac{ϵ(3+2ϵ)}{\sqrt{3}(3+ϵϵ^2)}.$$
(9)
$`|m_0/M_{1/2}|^2`$ is then plotted using this relation against $`ϵ`$ in Figure 1b. The figure indicates that M-theory does not reproduce S5 for any realistic value of $`ϵ`$.
To summarise, Table 1 shows that S3-S4 are compatible with Eq.s (5)-(7), and therefore M-theory. However, Refs. show that each of the points corresponding to S3-S4 is in conflict either with UFB, CCB or REWSB constraints. In fact, all of the M-theory parameter space examined in Ref. was shown to be in conflict with one of these constraints.
We have therefore shown that while the LHC SUGRA points include models which may be derived from weakly or strongly coupled strings, they possess potentially catastrophic global CCB and/or UFB minima. The existence of a global CCB or UFB minimum does not necessarily rule out a model. Some models have meta-stable minima with lifetimes longer than the current age of the universe . The question of which minimum the VEV of scalar fields rest in is one of cosmology, and beyond the scope of this paper. We therefore take the view that if the bounds from CCB/UFB global minima are not valid, examples of weakly/strongly coupled string models are approximated by or included within S1-S5 so there is no need to construct another similar model to examine the string-derived SUSY phenomenology at the LHC. If one should take the global CCB/UFB minimum bounds strictly however, Ref. indicates that no variation of parameters in the class of string models considered above will result in a model without CCB/UFB problems. We thus conclude that it is useful to examine a new class of string model which does not have difficulty evading CCB/UFB constraints. We note that it is also possible to evade the CCB/UFB constraints by lowering the string scale in type I string models .
### 2.3 Optimized string model
We now turn to the analysis of a model specifically designed to provide a large region of parameter space without CCB/UFB problems . It is a weak coupling model, where the modular weights have been chosen so that
$$n_{Q_L}=n_{d_R^c}=n_{u_R^c}=2,n_{L_L}=n_{e_R^c}=n_{H_{1,2}}=1.$$
(10)
The model is non-universal, but still in a family independent way, and so avoids serious problems associated with flavour changing neutral currents. As a consequence of the assignments in Eq. (10), the string scale soft masses are
$`m_{H_{1,2}}^2=m_{L_L}^2=m_{e_R^c}^2=M_{3/2}^2\mathrm{sin}^2\theta ,`$
$`A_{L_LH_1e_R^c}=M_{1/2}=\sqrt{3}M_{3/2}\mathrm{sin}\theta ,`$
$`m_{Q_L}=m_{d_R^c}=m_{u_R^c}=0,A_t=M_{3/2}(\sqrt{2}{\displaystyle \frac{\sqrt{3}}{2}}),`$ (11)
so that the squarks are light and the sleptons heavy at the unification scale. This has the effect of ameliorating the CCB/UFB problem. It should be noted that the $`A`$-terms of squarks other than the stop are not calculable for small $`\mathrm{tan}\beta `$. We shall approximate them here to be equal to $`A_t`$, but in fact they have a negligible effect upon the phenomenology/spectrum unless $`\mathrm{tan}\beta `$ is large, in which case $`A_b=A_t`$.
To be definite, we choose model parameters $`\mathrm{tan}\beta =10`$, $`M_{3/2}=250`$ GeV and $`\theta =\pi /4`$. We call this optimized model O1 and numerically its parameters are
$`m_{H_{1,2}}=m_{L_L}=m_{e_R^c}=177\text{ GeV},`$
$`A_{L_LH_1e_R^c}=M_{1/2}=306\text{ GeV},`$
$`m_{Q_L}=m_{d_R^c}=m_{u_R^c}=0,A_{Q_LH_2u_R^c}=A_{Q_LH_1d_R^c}=137\text{ GeV.}`$ (12)
We make the approximation that these relations hold at $`M_{GUT}2\times 10^{16}`$ GeV, but it should be borne in mind that logarithmic corrections from renormalisation between $`M_{\text{string}}5\times 10^{17}`$ and $`M_{GUT}`$ are expected. The spectrum and decay chains of the sparticles are calculated using ISASUGRA and ISASUSY using Eq. (12) as input.
Using three different Monte Carlo programs and the CTEQ3L parton distribution functions , we calculate the total cross sections of SUSY particles at the LHC for models S1-S5 and O1. Table 2 shows the comparison of HERWIG6.0 <sup>2</sup><sup>2</sup>2This version of HERWIG, not officially released, was a developmental version of HERWIG6.1 . total cross-section with those calculated by ISAJET7.40 and SPYTHIA . In each case, we have used the CTEQ3L parton distribution functions and the spectrum is calculated by ISASUGRA with $`m_t=175`$ GeV. As can be seen from the table, the three Monte-Carlo programs agree to about 10$`\%`$.
Fig. 2 displays the result of the calculation using HERWIG6.0. It is noticeable from the figure that S5 and O1 have broadly similar SUSY production cross-sections except for the sleptons which are noticeably lower. This can be understood by comparing the spectra of the two models, displayed in Table 3. The spectra are approximately similar, except for the sleptons which are heavier in O1 relative to the squarks. This is a consequence of the different choice of modular weights for the sleptons and the squarks. We therefore propose to analyse LHC SUSY production in O1 using S5 as a benchmark or comparison. We will show that the two models can be distinguished experimentally.
## 3 Experimental observability
### 3.1 Method
The primary experimental aim is to take previously developed model-independent methods for measuring SUSY particle masses (which were developed at standard SUGRA points) and by testing them in the context of a new optimised model, identify where these methods need to be generalised to perform well in both models. Secondarily it is to be shown that, after modifying these methods’ treatments of the slepton sector, their performance is sufficient to distinguish between the optimised and standard SUGRA scenarios.
To accomplish these aims, some model-dependent assumptions have to be made. In this analysis, R-parity is taken to be conserved and certain sparticle decay chains are assumed to exist. In R-parity conserving (RPC) models, SUSY particles are only produced in pairs, and the lightest SUSY particle (LSP) is stable. SUSY events contain two of these LSPs which, being only weakly interacting, escape detection and lead to SUSY events with large amounts of missing transverse energy – the standard signature for RPC models. The fact that these massive particles go missing from all RPC SUSY events means that in these models it is not usually possible to measure particle masses by reconstructing entire decay chains. Thus, as in , endpoints in kinematic variables constructed from SUSY decay chains must instead be examined. Specifically, the “sequential” decay mode $`\stackrel{~}{q}_L\stackrel{~}{\chi }_2^0q\stackrel{~}{l}_R^{}l_{\mathrm{near}}^\pm q\stackrel{~}{\chi }_1^0l_{\mathrm{far}}^{}l_{\mathrm{near}}^\pm q`$ (see Figure 3) and the “branched” decay modes $`\stackrel{~}{q}_L\stackrel{~}{\chi }_2^0q\stackrel{~}{\chi }_1^0Xq\stackrel{~}{\chi }_1^0l^\pm l^{}q`$ (see Figure 4) form the starting point for this investigation.
The kinematic edges used in to identify particle masses at S5 contain:
* $`l^+l^{}`$ edge: This picks out, from the “sequential” decays, the position of the very sharp edge in the dilepton invariant mass spectrum caused by $`\stackrel{~}{\chi }_2^0l\stackrel{~}{l}`$ followed by $`\stackrel{~}{l}l\stackrel{~}{\chi }_1^0`$.
* $`l^+l^{}q`$ edge: In “sequential” decays, the $`llq`$ invariant mass spectrum contains a linearly vanishing upper edge due to successive two-body decays. Our theoretical model of the $`l^+l^{}q`$ edge is more model independent than that used in earlier studies as it does do not assume any relation between the sparticle masses, other than the hierarchy $`0<m_{\stackrel{~}{\chi }_1^0}<m_{\stackrel{~}{l}_R}<m_{\stackrel{~}{\chi }_2^0}<m_{\stackrel{~}{q}}`$ necessary for the “seqential” decay chain to exist at all.
* $`l^+l^{}q`$ threshold: This is the non-zero minimum in the “sequential” $`llq`$ invariant mass spectrum, for the subset of events in which the angle between the two leptons (in the centre of mass frame of the slepton) is greater than $`\pi /2`$. This translates into the direct cut on $`m_{ll}`$ described in Section 3.4.1.
* $`hq`$ edge: This is one of the two instances of the $`Xq`$ edge. In general, the $`Xq`$ edge is the upper edge of the distribution of the invariant mass of three visible particles in the “branched” decays of Figure 4. The position of this edge is again determined by two-body kinematics. Depending on the Higgs mass and the mass difference between the $`\stackrel{~}{\chi }_2^0`$ and the $`\stackrel{~}{\chi }_1^0`$, one of the two “branched” decay chains will be strongly suppressed with respect to the other, so typically only one edge will be visible. At S5 the Higgs mode dominates, while the $`Z`$ mode will dominate at O1.
As detailed later in this analysis, new edges are added to the list above. Two of them are more general versions of other edges treated in , while another is entirely new. They are the $`l^\pm q`$ high-, $`l^\pm q`$ low- and $`M_{T2}`$ edges respectively. To summarise, the overall method we apply in this paper consists of:
* finding a model-independent set of cuts which can be used to select events from which the endpoints of all the kinematic edges may be measured,
* obtaining an estimate of the accuracy with which the edge positions might be determined by an LHC experiment,
* performing chi-squared fits of the expected positions of these edges (as functions of the sparticle masses) to a set of “simulated edge measurements” one might expect from an ensemble of such experiments, and
* interpreting the results as the statistics contribution to model independent sparticle mass measurements of a typical LHC experiment.
Were all squarks to have the same mass, then (with perfect detector resolution and with infinite statistics) the endpoints of the edges listed above would be found at the positions given in Table 4. Note that these positions depend on four unknown parameters, namely the masses of the squark, the slepton and the two neutralinos participating in each decay. The mass of the lightest Higgs boson also appears, but we assume this will already be known or will be obtained by other methods (e.g. ) to within 2%.
In a more realistic “non-degenerate squark masses” scenario, every distribution with a squark related edge will actually be a superposition of many underlying distributions (one for each squark mass), each having its edge at a slightly different location. This results in some smearing of the edges, even before detector effects (resolutions, acceptances, jet energy calibrations etc.) are taken into account.
It is not possible to measure each of the squark masses separately. Consequently the quantity $`m_{\stackrel{~}{q}}`$ in Table 4 will, after being obtained from the “smeared” edges, represent a squark mass scale rather than a specific squark mass. In all squark related edges, except the $`l^+l^{}q`$ threshold, the contribution to the outermost part of each edge is provided by the heaviest squarks. In the case of the $`l^+l^{}q`$ threshold, the “true” endpoint is set by the lightest squark. However, if (as at O1 and S5) the other eleven squarks are much heavier (see Table 3), it is easier to measure the contribution coming from them. A strong correlation between $`m_{\stackrel{~}{q}}`$ and the mass of the heaviest squark would therefore be expected.
More work would be required to fully understand the “theoretical” systematic errors including the use of a full simulation of the expected edge positions. Better edge models than those used later in this analysis will definitely be needed to permit measurements of edge positions to be better correlated with functions of the particle masses. Other sources of systematic errors which require further analysis are the detector effects mentioned above, combinatorial backgrounds near the edges and possible cut biases. Such systematic errors can only meaningfully be studied when real data are available.
Although they are beyond the scope of this paper, it is assumed that the above investigations could be performed so as to leave the eventual edge resolutions determined only by statistics and detector resolution. In this analysis, then, all edges are fitted with simple shapes (see Section 7) with the intention of extracting only an estimate of the uncertainty on the edge position, and not to obtain the edge position itself.
The first four edges listed in Table 4 thus constitute a minimal constraint on the four unknown sparticle mass parameters, which may then be further constrained by the other measurements which follow.
### 3.2 Other measurements
#### 3.2.1 Near, far, high and low $`l^\pm q`$ edges
In “sequential” decays there are three observable outgoing particles. There are only four different ways of grouping these particles together, and so only four different invariant masses may be formed from them: $`m_{ll}`$, $`m_{llq}`$, $`m_{l_{\mathrm{near}}q}`$ and $`m_{l_{\mathrm{far}}q}`$. The first two of these, $`m_{ll}`$ and $`m_{llq}`$, may be formed from the observed momenta without knowledge of which lepton was $`l_{\mathrm{near}}`$ and which was $`l_{\mathrm{far}}`$; only the total lepton four-momentum is required. Edges in these invariant mass distributions have already been described. If it were possible to tag the near- and far- leptons separately, the third and fourth invariant mass combinations could also be formed unambiguously, allowing the positions of two more edges, $`m_{l_{\mathrm{near}}q}^{\mathrm{max}}`$ and $`m_{l_{\mathrm{far}}q}^{\mathrm{max}}`$, to play a part in the final fit. However, such tagging is impossible. If further information is to be gathered from these decays in a model independent way, it is then necessary to look for edges in variables (functions of $`m_{l_{\mathrm{near}}q}`$ and $`m_{l_{\mathrm{far}}q}`$) which are observable. On an event-by-event basis, then, we define
$`m_{lq(\mathrm{high})}=\mathrm{max}(m_{l^+q},m_{l^{}q})\mathrm{max}(m_{l_{\mathrm{near}}q},m_{l_{\mathrm{far}}q})`$ (13)
and
$`m_{lq(\mathrm{low})}=\mathrm{min}(m_{l^+q},m_{l^{}q})\mathrm{min}(m_{l_{\mathrm{near}}q},m_{l_{\mathrm{far}}q}).`$ (14)
The simplest theoretical predictions for the positions of the corresponding edges, $`m_{lq(\mathrm{high})}^{\mathrm{max}}`$ and $`m_{lq(\mathrm{low})}^{\mathrm{max}}`$, are listed in Table 4, along with the positions of the near- and far-edges. Note that although the position of the high-edge is just the higher of $`m_{l_{\mathrm{near}}q}^{\mathrm{max}}`$ and $`m_{l_{\mathrm{far}}q}^{\mathrm{max}}`$, the position of the low-edge has a more interesting form. This asymmetry arises because it is not always kinematically possible for the invariant mass of the lepton/quark pair coming from the lower near/far distribution to approach $`\mathrm{min}[m_{l_{\mathrm{near}}q}^{\mathrm{max}},m_{l_{\mathrm{far}}q}^{\mathrm{max}}]`$ arbitrarily closely, while simultaneously requiring that this invariant mass is less than that of the other lepton/quark pair.
#### 3.2.2 $`Zq`$ edge
Since the neutralino mass difference at O1 ($`109\mathrm{GeV}`$) is too small to permit $`\stackrel{~}{\chi }_2^0h\stackrel{~}{\chi }_1^0`$, “branched” decays are mediated by the $`Z`$. (Particle masses may be seen in Table 3.) The cuts developed in for picking out this special case of the $`Xq`$ edge at S2 are found to perform equally well at O1, so they are adopted unchanged. Although it might be possible to benefit from adapting these cuts to O1 slightly, this temptation is resisted in order to retain model independence.
#### 3.2.3 $`M_{T2}`$
In order to constrain the slepton and neutralino masses better, we construct another variable, $`M_{T2}(\chi )`$, whose distribution relates just these two masses. We base this on the variable, proposed in , which looks at events containing two identical two body decays: $`xyXYa_1a_2Yb_1c_1b_2c_2`$, where the particles of types $`a`$ and $`c`$ are of unknown mass, where particles of type $`c`$ are undetectable, where the longitudinal momentum of the incoming particles is also unknown, and where it is assumed that $`Y`$ does not contain any unobservable particles such as neutrinos. In such cases, the variable provides a kinematic constraint on the masses of $`a`$ and $`c`$. We seek to apply this primarily to LHC dislepton events of the form $`q\overline{q}\stackrel{~}{l}_R^+\stackrel{~}{l}_R^{}l^+\stackrel{~}{\chi }_1^0l^{}\stackrel{~}{\chi }_1^0`$ and so define our $`M_{T2}(\chi )`$ by:
$`M_{T2}^2(\chi )\underset{\text{ / }𝐩_1+\text{ / }𝐩_2=\text{ / }𝐩_T}{\mathrm{min}}\left[\mathrm{max}\{m_T^2(𝐩_T^{l_1},\text{ / }𝐩_1,\chi ),m_T^2(𝐩_T^{l_2},\text{ / }𝐩_2,\chi )\}\right]`$ (15)
where
$`m_T^2(𝐩_T^l,𝐪_T,\chi )m_l^2+\chi ^2+2(E_T^lE_T^\chi 𝐩_T^l𝐪_T),`$ (16)
$`E_T^l=\sqrt{𝐩_{T}^{l}{}_{}{}^{2}+m_l^2}\text{and}E_T^\chi =\sqrt{𝐪_T^2+\chi ^2}.`$ (17)
This definition includes the lepton masses for completeness, although these are neglected in all computations.
The value $`M_{T2}(\chi )`$ takes for a given candidate dislepton event is a function of: the transverse missing-momentum vector, $`\text{ / }𝐩_T`$; the transverse momentum vectors of the two leptons, $`𝐩_T^{l_1}`$ and $`𝐩_T^{l_2}`$; and one other parameter – an estimate of the neutralino mass, $`\chi `$ (not to be confused with the actual mass of the neutralino, $`m_{\stackrel{~}{\chi }_1^0}`$). Unlike the other parameters, the value of $`\chi `$ is not measured in each event – events may be reinterpreted for different values of $`\chi `$. $`M_{T2}(\chi )`$ has the property that, for signal events in a perfect detector,
$`\underset{\mathrm{events}}{\mathrm{max}}\left[M_{T2}(m_{\stackrel{~}{\chi }_1^0})\right]=m_{\stackrel{~}{l}}.`$ (18)
Thus when $`\chi `$ is indeed $`m_{\stackrel{~}{\chi }_1^0}`$, then the distribution of $`M_{T2}`$ has an end point at the slepton mass. Since the other observables allow $`m_{\stackrel{~}{\chi }_1^0}`$ to be measured in a model independent way, $`M_{T2}`$ can then be included in the analysis to provide an additional constraint on the slepton/neutralino mass difference.
In practice the $`M_{T2}`$ edge can be distorted by the finite resolution of the detector or missing energy from soft underlying events. Additionally, standard model (SM) backgrounds provide constraints on minimum detectable slepton/neutralino mass differences (see Section 10). However the most important factor to control is the ability to correctly identify which particle species is contributing to an observed $`M_{T2}`$ edge – particularly since it is possible to have multiple edges in the non standard model contributions to $`M_{T2}`$ distributions. At S5, for example, where both the right- and left-sleptons are lighter than at O1 (see Table 3), both $`\stackrel{~}{l}_R^+\stackrel{~}{l}_R^{}`$ and $`\stackrel{~}{l}_L^+\stackrel{~}{l}_L^{}`$ events pass the cuts, and two edges are generated. The edge coming from the (lighter) right-slepton falls well under the SM background and so is not measurable, while the (heavier) left-slepton still has a cross section for pair-production high enough to let it form a good edge of its own at $`M_{T2}(m_{\stackrel{~}{\chi }_1^0})=m_{\stackrel{~}{l}_L}`$. It is important that this edge is not mistaken for the $`\stackrel{~}{l}_R`$, so methods are needed to dismiss it. A detailed prescription of how to go about such a dismissal is beyond the scope of this paper, but it would clearly be accomplished by looking for inconsistency between a given edge-particle hypothesis and all the other sparticle masses, the branching ratios and (in particular) the strongly mass dependent pair-production cross sections, which could all be measured by other means.
### 3.3 Event generation and detector simulation
All events, except those from $`q\overline{q}W^+W^{}`$ background processes, are simulated by HERWIG6.0. The $`W`$-pair events are generated by ISAJET7.42 . The detector chosen for simulation is the ATLAS detector , one of the two general purpose experiments scheduled for the LHC. The LHC is expected to start running at a luminosity of about $`10^{33}\mathrm{cm}^2\mathrm{s}^1`$ and this is expected to be increased over a period of about three years to the design luminosity of $`10^{34}\mathrm{cm}^2\mathrm{s}^1`$. These two periods are referred to, respectively, as the periods of low and high luminosity running.
The performance of the detector is simulated by ATLFAST2.16 which is primarily a fast calorimeter simulation which parametrises detector resolution and energy smearing and identifies jets and isolated leptons, in both the low and high luminosity environments. Throughout this analysis, the parameters controlling ATLFAST’s jet and lepton isolation criteria are left with the default values appropriate to the apparatus: i.e. jets must satisfy $`p_T^j15\mathrm{GeV}`$ and must lie in the pseudo-rapidity range $`5\eta ^j5`$, while electrons must have $`p_T^l5\mathrm{GeV}`$, muons $`p_T^l6\mathrm{GeV}`$ and both must lie in $`2.5\eta ^l2.5`$. For lepton isolation a maximum energy of $`10\mathrm{GeV}`$ may be deposited in a cone about the lepton of radius 0.2 in $`(\eta ,\varphi )`$-space ($`\varphi `$ being the azimuthal angle) while its separation from other jets must be at least 0.4 in the same units.
At high luminosity approximately 20 minimum bias events (“pile-up” events) are expected to occur in each bunch crossing. Pile-up events are not simulated by ATLFAST2.16, although it does alter its reconstruction resolutions to reflect the two different luminosity environments. It must be checked that any cuts applied at high luminosity will not be affected substantially by pile-up events. Within this article, events corresponding to $`100\mathrm{fb}^1`$ are generated and are reconstructed in the high luminosity environment. This approximately corresponds to one year of high luminosity running.
### 3.4 Event selection cuts
Section 3.4.1 summarises all the cuts used to obtain the edges. The $`l^+l^{}`$ edge, $`l^+l^{}q`$ edge and $`l^+l^{}q`$ threshold cuts do not differ significantly from those used in . The $`l^\pm q`$ edge cuts, however, do. The most significant change is the relaxation of the splitting requirement. The original splitting requirement tries to guarantee that the jet which comes from the quark produced in association with the observed dilepton pair is correctly identified. It achieves this by insisting that:
* both the dilepton pair and one of the two hardest jets $`j_1`$ and $`j_2`$ (ranked by $`p_T`$) are consistent with being the decay products of a squark (i.e. $`m_{llj_i}<m_{\mathrm{cutoff}}`$), and
* the invariant mass $`m_{llj_j}`$ of the two leptons and the other of the two highest $`p_T`$ jets is inconsistent with these being the decay products of a squark (i.e. $`m_{llj_j}>m_{\mathrm{cutoff}}`$).
Although the above demand for inconsistency increases the purity of the signal, it has a significant detrimental effect on the efficiency. In our analysis the consistency requirement is retained but the inconsistency requirement dropped. Instead we require that $`m_{ll}`$ is inside the expected region, given the $`l^+l^{}`$ edge measurement. We also perform background subtraction, modelling the opposite-sign same-lepton-family (OSSF) backgrounds by the distributions produced by those opposite-sign different-lepton-family (OSDF) events which pass the same cuts.
The production cross section for dislepton events suitable for $`M_{T2}`$ analysis is, in all models, by far the smallest (see Figure 2). So to show that $`M_{T2}`$ events passing cuts are not in danger of being swamped by small variations in backgrounds or the cuts themselves, two set of cuts referred to as “hard” and “soft” are developed. The “hard” cuts attempt to maximise the signal to background ratio in the vicinity of the dislepton edge in the $`M_{T2}`$ spectrum, while the much looser “soft” cuts try to maximise statistics at the edge at the expense of allowing in additional SUSY backgrounds.
#### 3.4.1 Cuts listed by observable
This section summarises the cuts used in the analysis. For notational purposes, reconstructed leptons and jets are sorted by $`p_T`$. For example $`j_2`$, with corresponding transverse momentum $`𝐩_T^{j_2}`$, is the reconstructed jet with the second highest $`p_T`$. All cuts are requirements unless stated otherwise.
$`n_{\mathrm{leptons}}=2`$, both leptons OSSF and $`p_T^{l_1}p_T^{l_2}10\mathrm{GeV}`$.
$`n_{\mathrm{jets}}2`$ and $`p_T^{j_1}p_T^{j_2}150\mathrm{GeV}`$. $`\text{ / }p_T300\mathrm{GeV}`$.
The kinematics of OSSF leptons coming from background processes which produce uncorrelated leptons (for example tau decays) are modelled well by OSDF lepton combinations. Consequently, edge resolution is improved by “flavour subtracting” OSDF event distributions from OSSF event distributions.
$`n_{\mathrm{leptons}}=2`$, both leptons OSSF and $`p_T^{l_1}p_T^{l_2}10\mathrm{GeV}`$.
$`n_{\mathrm{jets}}4`$, $`p_T^{j_1}100\mathrm{GeV}`$, and $`p_T^{j_1}p_T^{j_2}p_T^{j_3}p_T^{j_4}50\mathrm{GeV}`$.
$`\text{ / }p_T\mathrm{max}(100\mathrm{GeV},0.2M_{\mathrm{effective}})`$ and $`M_{\mathrm{effective}}400\mathrm{GeV}`$, where $`M_{\mathrm{effective}}`$ is defined by the scalar sum:
$$M_{\mathrm{effective}}=\text{ / }E_T+p_T^{j_1}+p_T^{j_2}+p_T^{j_3}+p_T^{j_4}.$$
(19)
Since the desired edge is a maximum, only the smaller of the two $`m_{llq}`$ combinations which can be formed using $`j_1`$ or $`j_2`$ is used.
Flavour subtraction is employed here as at the $`l^+l^{}`$ edge.
All the cuts for the $`l^+l^{}`$ edge as above.
In addition $`m_{ll}^{\mathrm{max}}/\sqrt{2}m_{ll}m_{ll}^{\mathrm{max}}`$, where the value of $`m_{ll}^{\mathrm{max}}`$ would be obtained from the $`l^+l^{}`$ edge in practice, but was determined theoretically in this analysis.
Since the desired edge is a minimum, only the larger of the two $`m_{llq}`$ combinations which can be formed using $`j_1`$ or $`j_2`$ is used.
$`n_{\mathrm{leptons}}=0`$ and $`\text{ / }p_T>300\mathrm{GeV}`$.
Exactly two $`b`$ jets with $`p_T^{j_{b_1}}p_T^{j_{b_2}}50\mathrm{GeV}`$. No other $`b`$ jets, regardless of $`p_T`$.
At least two non-$`b`$ jets with $`p_T^{j_{q_1}}p_T^{j_{q_2}}100\mathrm{GeV}`$, with at least one inside $`2.0\eta ^j2.0`$.
$`m_{bb}`$ within $`17\mathrm{GeV}`$ of Higgs peak in $`m_{bb}`$ spectrum.
Since the desired edge is a maximum, the non-$`b`$ jet (i.e. $`j_{q_1}`$ or $`j_{q_2}`$) chosen to form the $`m_{hq}`$ invariant mass is that which minimises $`m_{hq}`$.
$`n_{\mathrm{leptons}}=2`$, both leptons OSSF and $`p_T^{l_1}p_T^{l_2}10\mathrm{GeV}`$.
$`m_{ll}`$ within 2.5 GeV of the centre of the Z-mass peak in the $`m_{ll}`$ spectrum.
At least two non-$`b`$ jets exist with $`p_T^j100\mathrm{GeV}`$ and $`\text{ / }p_T300\mathrm{GeV}`$.
Since the desired edge is a “maximum”, the jet chosen to form the $`m_{Zq}`$ invariant mass is the one (from those with $`p_T^j100\mathrm{GeV}`$) which minimises $`m_{Zq}`$.
Again, flavour subtraction is used.
All the cuts for the $`l^+l^{}q`$ edge above are required.
Additionally, events consistent with the $`l^+l^{}`$ edge measurement are selected by asking for $`m_{ll}m_{ll}^{\mathrm{max}}+1\mathrm{GeV}`$.
To choose the jet from which to form $`m_{l^\pm q}`$ we select $`j_i=j_1`$ or $`j_2`$, whichever gives the smaller value of $`m_{llj}`$, and require $`m_{llj_i}<m_{\mathrm{cutoff}}`$, where $`m_{\mathrm{cutoff}}`$ is chosen to be above the $`l^+l^{}q`$ edge, but is otherwise arbitrary. For easy comparison with , $`m_{\mathrm{cutoff}}=600\mathrm{GeV}`$ was used in this analysis.
Finally, the two invariant mass combinations, $`m_{l^\pm q}`$, are assigned to $`m_{lq(\mathrm{high})}`$ and $`m_{lq(\mathrm{low})}`$ as defined earlier by Equations (13) and (14).
Events are required to have exactly one OSSF pair of isolated leptons satisfying $`p_T^{l_1}>50\mathrm{GeV}`$ and $`p_T^{l_2}>30\mathrm{GeV}`$.
$`\delta _T<20\mathrm{GeV}`$ is required, where $`\delta _T=|𝐩_T^{l_1}+𝐩_T^{l_2}+\text{ / }𝐩_T|`$. Large $`\delta _T`$ in signal events is indicative of an unidentifiable transverse boost to the centre of mass frame, perhaps due to initial state radiation. Large $`\delta _T`$ in other (not necessarily signal) events simply suggests an inconsistency with the desired event topology.
Events containing one or more jets with $`p_T^j>40\mathrm{to}50\mathrm{GeV}`$ are vetoed. This cut also helps to reduce standard model backgrounds, notably $`t\overline{t}`$. The lower this cut is placed, the better for the background rejection. However the cut cannot be placed too low (especially at high luminosity) due to the significant number of jets coming from the underlying event and other minimum bias events in the same bunch crossing. For order 25 minimum bias events per bunch crossing, estimates that about 10% (1%) of bunch crossings will include a jet from the underlying event with a $`p_T`$ greater than 40 GeV (50 GeV). Our results are not sensitive to variation of the jet veto cut between 40 and 50 GeV, where at least 90% signal efficiency is expected.
Events with $`|m_{l_1l_2}m_Z|<5\mathrm{GeV}`$ are vetoed to exclude lepton pairs from $`Z`$ bosons.
$`m_{l_1l_2}>80\mathrm{GeV}`$ and $`\text{ / }p_T>80\mathrm{GeV}`$ are also required.
These are as above, but
* the $`\text{ / }p_T`$ requirement is lowered from 80 to 50 GeV,
* the upper limit for $`\delta _T`$ is extended from 20 to 90 GeV, and
* the dilepton invariant mass cut is removed altogether.
### 3.5 Edge resolutions
By applying at S5 the cuts listed in Section 3.4.1 we reproduce the $`m_{ll}`$ and $`m_{llq}`$ distributions whose edges are analyzed in (see Figure 5: a, b and d). We also confirm that the same cuts may also be used to generate edges of similar quality at O1 (see Figure 6: a, b and d). In addition, plots c1 and c2 in Figures 5 and 6 display the $`m_{lq(\mathrm{high})}`$ and $`m_{lq(\mathrm{low})}`$ distributions generated from the modified $`l^\pm q`$ edge cuts also listed above. Numerical results are summarised in Table 5.
Similar pictures are seen at S5 and O1. The most obvious difference between the two sets of data is the peak at the $`Z`$ mass present in Figure 6a but absent in Figure 5a. This peak comes from direct $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0Z\stackrel{~}{\chi }_1^0l^+l^{}`$ decays, which have a branching fraction of 39% at O1 compared with 0.65% at S5. The most common direct neutralino decay mode at S5 is through the $`h_0`$ ( 65% ) of which only a tiny proportion ( 0.02% ) decay to the two light leptons – the partial decay width $`\mathrm{\Gamma }(h_0ff)`$ going approximately as $`m_f^2`$. Scenarios in which the $`l^+l^{}`$ edge happens to coincide with the $`Z`$ peak require special treatment and are not considered here.
In order to make useful statements about the degree to which model parameters may be extracted from the endpoints and edges of these distributions, it is necessary to obtain an estimate of the accuracy with which these observables may be measured. It is expected that the errors on all of the observables considered here will eventually be statistics dominated, so simple fits have been made to the data to obtain estimates of the statistical errors on the edge or end point locations. The shapes fitted to the data (see Figure 7) and the algorithms for determining the boundaries of the fitted regions have been kept as simple and generic as possible, with the intention of making the fit results both conservative and simple to interpret. Given real data, it would be worth putting substantial effort into understanding how detector and physics effects affect the shape of each distribution in turn. This could significantly improve particle mass resolutions. The resolutions obtained from the fits to the distributions in Figure 5 and Figure 6 are listed in Table 5.
The $`M_{T2}`$ distributions obtained after the cuts listed in Section $`M_{T2}`$ edge (soft cuts) are shown in Figure 8. The signal dislepton region having the desired edge is the unhatched region on each plot. Conveniently, most of the other SUSY events passing the cuts (mainly events containing gauginos) are also distributed with an edge located at a similar position to the disleptons’. This is due to the fact that such events often include pairs of slepton decays, and these may sometimes occur in combination with low jet activity and without additional leptons being produced inside detector acceptance. Consequently, for much of the SUSY background which passes cuts, the cuts are accepting $`M_{T2}`$ values which do not degrade edge performance significantly. Note that $`\mathrm{\Delta }M(\chi )=M_{T2}(\chi )\chi 0`$ is plotted evaluated at $`\chi =m_{\stackrel{~}{\chi }_1^0}`$, so this should have an edge at $`\mathrm{\Delta }M_{\mathrm{max}}=m_{\stackrel{~}{l}}m_{\stackrel{~}{\chi }_1^0}`$. Of course, $`m_{\stackrel{~}{\chi }_1^0}`$ is not actually known a priori, so generating this graph in a real experiment requires an estimate of $`m_{\stackrel{~}{\chi }_1^0}`$ to be obtained by other means. For our purposes it will be sufficient if the position of the edge of the distribution can be measured to about 10%, and this determines the accuracy required of the neutralino mass estimate $`\chi `$. The width $`\mathrm{\Delta }M(\chi )=M_{T2}(\chi )\chi `$ remains approximately stable at the 10% level for $`m_{\stackrel{~}{\chi }_1^0}/2<\chi <2m_{\stackrel{~}{\chi }_1^0}`$. To illustrate the relative insensitivity of $`\mathrm{\Delta }M_{\mathrm{max}}(\chi )`$ to $`\chi `$ near $`m_{\stackrel{~}{\chi }_1^0}`$, plots generated from the same data as before but with $`\chi m_{\stackrel{~}{\chi }_1^0}\pm 50\mathrm{GeV}`$ are shown in Figure 9. Satisfying the 10% requirement above by obtaining a suitable value of $`\chi `$ within such a $`\pm 40\%`$ range is not difficult and may be accomplished by first performing a “cut down” version of the analysis described later in the text, but omitting the $`M_{T2}`$ data, and then choosing the estimate, $`\chi `$, to be the resulting reconstructed neutralino mass.
Determining the likely resolution for the $`M_{T2}`$ edge requires a slightly different approach to that used for the other edges. Whereas the cuts used to obtain the invariant mass distributions have high SM rejections (primarily due to the presence of at least one high mass $`\text{ / }p_T`$ or $`M_{\mathrm{effective}}`$ cut), the $`M_{T2}`$ cuts cannot be so hard, primarily because the desired dislepton events have have very little jet activity. With the dislepton production cross sections typically two orders of magnitude smaller than the squark/gluino production cross sections (Table 2) a low efficiency is not affordable. There are also irreducible SM backgrounds (primarily $`W^+W^{}l^+l^{}\nu \overline{\nu }`$ but also $`t\overline{t}b\overline{b}W^+W^{}jjl^+l^{}\nu \overline{\nu }`$ in cases where jets are below the $`p_T`$ cut or outside detector acceptance) which have signatures identical to dislepton events. These backgrounds are clearly visible in Figure 8 and would cause problems for naive straight line fitting techniques.
The standard model backgrounds in Figure 8, although large, are easy to control. The SM edge is very clean, since the events which pass the cuts are well reconstructed and there are no noticeable tails. The SM edge will fall at $`m_W`$ for $`\chi =0`$, corresponding to the mass of the missing neutrino in SM events. As $`\chi `$ increases, the SM edge recedes to lower masses as shown in Figure 10, falling to $`60\mathrm{GeV}`$ for $`\chi =100\mathrm{GeV}`$. A significant excess of events above this threshold would be a clear signal for non-SM processes. Such an excess will appear when the $`\stackrel{~}{l}`$-$`\stackrel{~}{\chi }_1^0`$ mass difference exceeds $`80\mathrm{GeV}`$ for low $`m_{\stackrel{~}{\chi }_1^0}`$ and $`50\mathrm{GeV}`$ for high $`m_{\stackrel{~}{\chi }_1^0}`$.
We estimate that, using $`M_{T2}`$ with either set of cuts, it is possible to measure $`\mathrm{\Delta }M_{\mathrm{max}}=m_{\stackrel{~}{l}}m_{\stackrel{~}{\chi }_1^0}`$ to 10% or better. The “hard” $`M_{T2}`$ cuts successfully remove almost all SUSY background above the SM threshold at the expense of only retaining half of the events passing the “soft” cuts. For both sets of cuts the SM threshold, at about 60 GeV, would be approximately three sigma away from $`\mathrm{\Delta }M_{\mathrm{max}}`$ at this accuracy.
### 3.6 Reconstructing sparticle masses
Sparticle masses are reconstructed by performing a chi-squared fit with between four and six free parameters $`p`$. The main parameters of the fit are $`m_{\stackrel{~}{\chi }_1^0}`$, $`m_{\stackrel{~}{l}}`$, $`m_{\stackrel{~}{\chi }_2^0}`$ and $`m_{\stackrel{~}{q}}`$. Where appropriate, $`m_h`$ and $`m_Z`$ also appear as fit parameters, although if they do, they are strongly constrained (particularly in the case of the $`Z`$) by present LEP measurements or expected LHC errors (0.0077% and 2.0% respectively).
The chi-squared, as a function of the free parameters $`p`$, is then formulated as:
$$\chi ^2(p)=\underset{1}{\overset{n}{}}\frac{(O_i^{\mathrm{sm}}O_i(p))^2}{\sigma _{i}^{}{}_{}{}^{2}},$$
(20)
where $`O_i^{\mathrm{sm}}=O_i(p_{\mathrm{model}})+\sigma _iX_i`$ are $`n`$ “smeared observables”, $`X_iN(0,1)`$ are the $`n`$ random variables from the Normal distribution (with mean 0 and standard deviation 1) which accomplish the smearing, $`\sigma _i`$ is the anticipated statistical error for the $`i^{\mathrm{th}}`$ observable (approximated by the fit error listed in Table 5), and $`O_i(p)`$ is the value one would expect for the $`i^{\mathrm{th}}`$ observable given the parameters $`p`$ (for examples see Table 4). $`p_{\mathrm{model}}`$ indicates the actual masses of the sparticles in the particular model being considered. The results of the fit, $`p_{\mathrm{fit}}`$, are then those which minimise the chi-squared. $`p_{\mathrm{fit}}`$ may be interpreted as the sparticle masses which might be reconstructed by an LHC experiment after obtaining $`100\mathrm{f}\mathrm{b}^1`$ at high luminosity, with the $`X_i`$ parametrising the experimental errors. In order to determine the accuracy with which this reconstruction can be performed, the above fitting process is repeated many times for different values of the $`X_i`$, producing the distributions of reconstructed particle masses which follow.
Figures 11 and 12 show the probability distributions expected for the reconstructed $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{l}`$, $`\stackrel{~}{\chi }_2^0`$ and $`\stackrel{~}{q}`$ masses, while Figures 13 and 14 show the corresponding fractional errors, $`\epsilon _p`$, for the same quantities. All are approximately Gaussian and have reassuringly small tails. Statistics summarising these plots (widths, means and RMS values) are listed in Table 6.
It may be seen that in both O1 and S5 the widths of the mass distributions for all four particles are very similar. This is because, in localized regions of parameter space, the edges tend to constrain mass differences far better than absolute masses. Evidence of this may be seen in Figure 15 which shows the scatter of reconstructions in the $`m_{\stackrel{~}{\chi }_1^0}`$-$`m_{\stackrel{~}{l}_R}`$ plane. It is interesting to note that without the $`M_{T2}`$ constraint at O1, the fit’s chi-squared commonly has two distinct and competing comparable minima – one at high and one at low values of $`m_{\stackrel{~}{l}_R}`$. The simultaneous existence of these minima is a direct consequence of putting the near and far $`l^\pm q`$ edges on an equal footing in this analysis, allowing more than one interpretation for each of the high and low edges. The need to resolve this kind of ambiguity in model-independent investigations of this type illustrates the importance of establishing reliable model-independent ways of measuring the absolute scale of $`m_{\stackrel{~}{l}_R}m_{\stackrel{~}{\chi }_1^0}`$ (or a related quantity) even if only to an accuracy of 20-30%.
The clear gap between the reconstruction regions for O1 and S5 in Figure 15 supports the original claim that, systematic errors permitting, it will be possible to distinguish between the two scenarios.
## 4 Conclusions
Some of the five standard LHC SUGRA points are compatible with universal perturbative string and M-theory, but dangerous CCB/UFB breaking minima are present in each example. We therefore studied a perturbative string model which is optimized to ameliorate the CCB/UFB problems present in the other models. The optimized model is non-universal because the squarks and sleptons are split in mass at the string scale. We identify the SUGRA point with the most similar spectrum and hard SUSY production cross sections (S5) to compare the optimized model with. The main difference is that the sleptons are heavier and therefore have lower production cross-sections. We have demonstrated the existence of a method by which an LHC experiment will be able to measure the masses of the (lighter) sleptons and the two neutralinos at O1 in a largely model independent way. In a specific comparison of S5 and O1 we have shown that this method will be able to distinguish a SUGRA model from an optimised string model with very similar properties.
More importantly, we expect that the techniques developed here are general enough to be used to discriminate between other pairings of optimised and non-optimised models with similar characteristics. The optimized model analysis applies to a more general class of models than the string model itself. We could apply it to models with non-universal SUSY breaking terms at $`M_{GUT}`$ in which the squarks and sleptons are explicitly split in mass. These constitute a superset of the particular string model considered here.
## Acknowledgements
This work was partially supported by the U.K. Particle Physics and Astronomy Research Council. We thank D.J. Summers for helpful discussions. CGL also wishes to thank A.J. Barr, L.M. Drage, J.P.J. Hetherington and C. Jones for their help on numerous occasions.
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# hep-ph/0007199 LA-UR-00-2142 Pseudovector mesons, hybrids and glueballs
## Abstract
We consider glueball– (hybrid) meson mixing for the low–lying four pseudovector states. The $`h_1^{^{}}(1380)`$ decays dominantly to $`K^{}K`$ with some presence in $`\rho \pi `$ and $`\omega \eta `$. The newly observed $`h_1(1600)`$ has a $`D`$– to $`S`$–wave width ratio to $`\omega \eta `$ which makes its interpretation as a conventional meson unlikely. We predict the decay pattern of the isopartner conventional or hybrid meson $`b_1(1650)`$. A notably narrow $`s\overline{s}`$ partner $`h_1^{^{}}(1810)`$ is predicted.
The pseudovector ($`J^{PC}=1^+`$) $`s\overline{s}`$ ground state has the interesting property that its OZI allowed decay to open strangeness, i.e. $`K^{}K`$, which is a priori expected to be dominant, is severely suppressed by phase space. This not only makes the state anomalously narrow , but opens up the possibility that other decays could be significant. These can arise from $`u\overline{u}`$, $`d\overline{d}`$ components in the state, which can come from mixing with a glueball.
We solve Schwinger–type mass equations with linear masses, pioneered in refs. and motivated in refs. . In this approach the underlying nature of the meson, whether conventional or hybrid, is not specified. The primitive (bare) states are ideally mixed. Primitive isoscalar and isovector $`u\overline{u}`$, $`d\overline{d}`$ states are degenerate. In this work we further only allow $`SU(3)`$ symmetric glueball–meson coupling, with no meson–meson coupling. We restrict to ground state and first excited state mesons. It is known that such restriction is quite accurate if the glueball mass is far from those of the states , as is the case here.
The numerical input is as follows. The ratio between pseudovector and scalar glueball masses is evaluated in lattice QCD as $`1.70\pm 0.05`$ or $`1.73\pm 0.09`$ . Taking the world average scalar glueball mass as 1.6 GeV , this implies a (input) pseudovector glueball mass of 2.7 GeV. Our conclusions do not critically depend on this value. The primitive $`u\overline{u}+d\overline{d}`$ ground state is input as the $`b_1`$ mass . The physical masses of $`h_1(1170)`$, $`h_1^{^{}}(1380)`$ and the newly discovered $`h_1(1600)`$ at $`1594\pm 15_{60}^{+10}`$ MeV are used as input. We further assume that the difference between the primitive $`s\overline{s}`$ and $`u\overline{u}+d\overline{d}`$ masses is the same for the ground states and excited states. Lastly, the primitive excited $`u\overline{u}+d\overline{d}`$ mass, the most uncertain input, is taken as $`1650\pm 50`$ MeV. This is hence the assumed mass region for the yet undiscovered excited $`b_1`$ resonance.
The output of our analysis is as follows. The experimenally unobserved physical excited $`s\overline{s}`$ state ($`|h_1^{^{}}_2`$) is predicted at $`1810\pm 40`$ MeV. The difference between the primitive $`s\overline{s}`$ and $`u\overline{u}+d\overline{d}`$ masses, for both the ground and excited states, is $`180\pm 10`$ MeV, yielding a primitive $`s\overline{s}`$ ground state ($`|s\overline{s}_1`$) at $`1410\pm 10`$ MeV. This is consistent with $`1445\pm 41`$ MeV derived from quark model relations<sup>3</sup><sup>3</sup>3Combining $`K(^1P_1)+K(^3P_1)=K(1270)+K(1400)`$, $`b_1+(s\overline{s})_1=2K(^1P_1)`$ and $`K(^1P_1)b_1=K(^3P_1)a_1`$. Here all items are the corresponding masses. The $`{}_{}{}^{1}P_{1}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ kaon masses before mixing are $`K(^1P_1)`$ and $`K(^3P_1)`$ respectively, and the primitive $`s\overline{s}`$ ground state mass is $`(s\overline{s})_1`$.. The coupling, in the notation of refs. is $`g_1=0.19\pm 0.01`$ GeV for the ground states and $`g_2=0.19\pm 0.12`$ GeV for the excited states. The accurate former value is larger than values found for scalar and tensor mesons .
The valence content of the physical mesons is
$$|h_1_1=(0.22_{0.01}^{+0.02})|g+(0.06_{0.05}^{+0.03})|s\overline{s}_2+(0.12_{0.09}^{+0.05})|u\overline{u}_2+(0.17_{0.00}^{+0.01})|s\overline{s}_1+(\mathbf{0.95}_{0.01}^{+0.01})|u\overline{u}_1,$$
$$|h_1^{}_1=(0.13_{0.03}^{+0.02})|g+(0.06_{0.05}^{+0.03})|s\overline{s}_2+(0.13_{0.11}^{+0.08})|u\overline{u}_2+(\mathbf{0.96}_{0.03}^{+0.02})|s\overline{s}_1+(0.22_{0.03}^{+0.02})|u\overline{u}_1,$$
$$|h_1_2=(0.19_{0.08}^{+0.15})|g+(0.16_{0.16}^{+0.09})|s\overline{s}_2+(\mathbf{0.94}_{0.08}^{+0.07})|u\overline{u}_2+(0.20_{0.11}^{+0.16})|s\overline{s}_1+(0.14_{0.06}^{+0.11})|u\overline{u}_1,$$
$$|h_1^{}_2=(0.12_{0.01}^{+0.07})|g+(\mathbf{0.97}_{0.03}^{+0.04})|s\overline{s}_2+(0.21_{0.13}^{+0.22})|u\overline{u}_2+(0.06_{0.00}^{+0.03})|s\overline{s}_1+(0.05_{0.00}^{+0.03})|u\overline{u}_1.$$
where the states on the left and right are respectively the physical and primitive states. The first three physical states are the experimental states $`h_1(1170)`$, $`h_1^{^{}}(1380)`$ and $`h_1(1600)`$ .
Decays are now studied by using a finite width for the initial meson, and unless otherwise indicated, a narrow width approximation for the final mesons. Finite widths are implemented by smearing over relativistic Breit–Wigner shapes with Quigg – von Hippel energy dependent widths. Whenever a decay is OZI allowed from an ideally mixed initial state, we assume, for simplicity, that the initial state is $`100\%`$ ideally mixed. OZI forbidden decays are implemented by using the (small) valence contents above to calculate connected decays .
The decays of conventional mesons are studied in the $`{}_{}{}^{3}P_{0}^{}`$ model using the methods, conventions and parameters of refs. . Making the usual identification that the primitive ground state mesons are P–wave quark model states, we obtain the decay widths in Table 1. We note that although the experimentally observed $`K^{}K`$ mode is dominant, and similar to the total width of the state<sup>4</sup><sup>4</sup>4We find that mock meson phase space gives a $`K^{}K`$ partial width of $`191\pm 18`$ MeV, inconsistent with the experimental total width of the state . For near threshold decays of this type mock meson phase space always gives a substantially larger width than relativistic phase space. Mock meson phase space results are hence not quoted for near threshold decays in the tables. We note that the $`K^{}K`$ partial width calculation in ref. misses a flavour factor of two, and is hence unreliable. , the $`\rho \pi `$ mode is detectable. It is not as large relative to $`K^{}K`$ as one might expect from the limited phase space of $`K^{}K`$. Identification in $`\rho \pi `$ is complicated by the huge $`360\pm 40`$ MeV width of the $`h_1(1170)`$ mainly in $`\rho \pi `$ . This makes $`\rho \pi `$ an unattractive search channel for $`h_1^{^{}}(1380)`$, since no viable production processes are known which strongly produce the dominant $`s\overline{s}`$ component in $`h_1^{^{}}(1380)`$ as opposed to the dominant $`u\overline{u}+d\overline{d}`$ component in $`h_1(1170)`$. Although $`\omega \eta `$ is small, $`h_1^{^{}}(1380)`$ has recently been observed in this mode . Additional decay modes that have not been calculated but are expected to be small are decays to $`h_1(\pi \pi )_S`$ and direct three–body decays like $`\pi ^0\pi ^+\pi ^{}`$.
We proceed to analyse $`h_1(1600)`$. One has to allow for the possibility that the excited $`u\overline{u}`$, $`d\overline{d}`$ and $`s\overline{s}`$ states are hybrid mesons. The calculations for this possibility are performed in the Isgur–Paton flux–tube model with the standard parameters of ref. . The results are displayed in Table 2. The $`h_1(1600)`$ is predicted to decay from most to least prevalent to $`\rho \pi /\rho (1450)\pi ,K^{}K`$ and $`\omega \eta `$ in all interpretations of the state. A minor feature that distinguishes interpretations is the relative size of the $`\rho \pi `$ and the $`\rho (1450)\pi `$ modes. The main distinguishing feature is the ratio of $`D`$–wave to $`S`$–wave widths, which is consistently larger for the meson than the hybrid interpretation. For the meson interpretation the S–wave width is suppressed due to a node in the amplitude, making it sensitive to the wave function parameter $`\beta `$ employed. Table 2 shows that the $`D/S`$–wave width ratio in $`\omega \eta `$ is inconsistent with the experimental result $`0.3_{0.[13]}^{+0.1}`$ if $`h_1(1600)`$ is a conventional meson. In order to confirm this result, we perform three further checks. Firstly, we evaluate the ratio by taking the wave function parameter $`\beta `$ to be different for different mesons participating in the decay. Varying $`\beta `$ in the reasonable range $`0.350.45`$ GeV confirms the result. Secondly, using the full valence content of $`h_1(1600)`$ above, and allowing decay via the ground state P–wave meson component, confirms the result. Thirdly, experimental data has few D–wave events above 1.8 GeV . Restricting the Breit–Wigner smearing to invariant masses less than 1.8 GeV gives the nearest ratio to experiment in all these simulations, $`0.9_{0.5}^{+1.2}`$. This ratio is still outside the range allowed by experiment, although it is not far outside the range. Experimentally, it has not been established that the D–wave exists , so that the very small ratio predicted for a hybrid meson in Table 2 could be consistent with experiment. Thus current experiment makes the conventional meson interpretation of the $`h_1(1600)`$ unlikely, but allows the hybrid interpretation. This assumes that the state observed in experiment cannot be resolved into two separate states. Since the $`{}_{}{}^{3}P_{0}^{}`$ model has only been tested for a few $`D/S`$–wave width ratios , one needs further information. The total width $`384\pm 60_{100}^{+70}`$ MeV of $`h_1(1600)`$ is slightly more consistent with the hybrid interpretation. Future searches for $`h_1(1600)`$ should focus on obtaining the $`D/S`$–wave width ratio in the sizable $`\rho \pi `$ channel. The $`b_1\pi `$ mode distinguishes the two models of hybrid decay in Table 2.
We note that since the $`\rho `$ Regge trajectory dominates the $`\rho (1450)`$ and $`b_1`$ trajectories, and $`h_1(1600)`$ has a healthy $`\rho \pi `$ coupling for all interpretations, one expects the $`h_1(1600)`$ to be produced via natural parity exchange in the $`\pi ^{}p`$ collisions it has been observed in. This is confirmed in the experimental analysis , providing an independent check on our calculations. The non–observation of $`h_1(1600)`$ in unnatural parity exchange may put bounds on its $`b_1\pi `$ coupling, discriminating between different hybrid decay models.
In Table 3 the widths for the isopartner state $`b_1(1650)`$ are calculated. The channels that distinguish between conventional and hybrid meson interpretations of the state, $`\omega (1420)\pi `$ and $`\rho \rho `$, are difficult to access experimentally. However, $`D/S`$–wave width ratios remain an excellent distinguishing feature. Possible search channels are $`\omega \pi `$ and $`\rho \eta `$.
The widths for the undiscovered excited $`s\overline{s}`$ state $`h_1^{^{}}(1810)`$ are indicated in Table 4. It is interesting to note that the flux–tube model selection rule, which states that decays to $`S+S`$ states ($`K^{}K,\varphi \eta `$) are suppressed relative to $`P+S`$ states ($`K_1(1270)K`$) , is apparently violated. This is due to phase space. Whether the $`h_1^{^{}}(1810)`$ is a conventional or hybrid meson, it is surprisingly narrow. Excellent search channels are $`K^{}K`$ and $`\varphi \eta `$. The latter is especially interesting since it cannot come from a $`u\overline{u}+d\overline{d}`$ state via OZI allowed decay. Small OZI forbidden modes like $`\rho \pi `$ could also effect detection. A natural place to search for $`s\overline{s}`$ states is at Jefferson Lab, where the photon has a sizable coupling to $`s\overline{s}`$. Production is likely to be via diffractive exchange, as meson exchange involves OZI forbidden or evading processes.
We thank C. Amsler, P. Eugenio, E. Klempt and D. Weygand for useful discussions on their experimental data. This research is supported by the Department of Energy under contract W-7405-ENG-36.
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